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# 1 Introduction
## 1 Introduction
In the light of the no-hair and no-go theorems for the classical glueball solutions with or without gravity, the discovery of both smooth and black-hole solutions of the self-gravitating non-Abelian gauge theories was a big surprise (for a review see ). Among such solutions a physically interesting case is that of a spontaneous broken EYMH theory examined in , where both regular and black-hole solutions, i.e. gravitating sphalerons and sphaleron black-holes have been found. The stability analysis of this system has shown that the solutions are unstable .
Because of the physical importance of these objects, it is worthwhile to study generalizations of the couplings of the flat-space Lagangian to gravity. One of the simplest and best motivated extensions is the inclusion of an explicit coupling between the scalar field $`\mathrm{\Phi }`$ and the curvature of the spacetime $``$ of the form $`\xi \mathrm{\Phi }^{}\mathrm{\Phi }`$, where $`\xi `$ is a dimensionless coupling constant. There are many reasons to believe that a nonminimal coupling term appears. A nonminimal coupling is generated by quantum corrections even if it is absent in the classical action and is required in order to renormalize the theory .
In many physical situations, inclusion of a $`\xi 0`$ term leads to new interesting physical effects even at the classical level. Examples are the Bronnikov-Melnikov-Bacharova-Bekenstein conformal scalar hair , the inflationary scenario with a nonminimally coupled ”inflaton” field , and boson star solutions . For a review of the present situation see . Two cases occur most frequently in the literature: ”minimal coupling” ($`\xi `$=0) and ”conformal coupling” ($`\xi `$=1/6). The conformal invariance dictates $`\xi `$=1/6 for a massless scalar field , while Nambu-Goldstone bosons have a minimal coupling $`\xi `$=0 . However, there is no preferential value of $`\xi `$ for a Higgs field in a unified gauge theory of electroweak interactions.
In this paper we study numerically regular and black-hole solutions of the coupled EYMH field equations with a nonminimal coupling to gravity, extending the results of ref. to this case. Ref. presented strong numerical arguments for the existence of both regular and black-hole solutions in a minimally coupled EYMH theory. For each fixed value of the Higgs vacuum expectation value $`v`$, solutions have been found, that can be indexed by the number of nodes $`k`$ of the Yang-Mills potential function. For each $`k`$ there are two branches of solution, depending on the behavior of $`v0`$. For example, the so-called quasi $`k=0`$ branch of solutions approaches the Schwarzschild solution as $`v0`$; whereas the regular $`k=1`$ branch approaches the first colored black-hole solution of the Einstein-Yang-Mills system in the same limit. The two branches of solutions converge for some values of the theory parameters .
Although most of the phenomena discussed in for the $`\xi `$=0 case repeat themselves in the general case, there are some important differences. For a nonminimal coupling, the time component of $`T_{\mu \nu }`$, which in Einstein’s gravity would correspond to the local energy density, may be non-positive. Indeed, as we shall see later on, there are regions in space, where this quantity is negative. The reason is that, as a result of the nonminimal coupling with gravity, there are contributions to $`T_{\mu \nu }`$ from the gravitational field itself . However, the local energy densities do not directly determine the sign of the asymptotic ADM mass, which is found to be positive.
Also, the parameter range of the solutions found in ref. remains no longer valid and a new range has to be found for every choice of $`\xi `$. The existence of a nonminimal coupling between the Higgs field and the gravitational field implies a decrease of the maximal allowed vacuum expectation value of the Higgs field.
The paper is structured as follows: in section II we present the general framework and an analysis of the field equations, while in section III we adress the problem of the numerical construction of solutions. In section IV the possibility of the existence of Lorentzian wormholes is considered with a negative result. We conclude with section V where the results are compiled.
## 2 GENERAL FRAMEWORK AND BASIC EQUATIONS
Our study of the EYMH system is based upon the action
$$S=d^4x\sqrt{g}[\frac{}{16\pi G}\frac{1}{4\pi }((D_\mu \mathrm{\Phi })^{}(D^\mu \mathrm{\Phi })+V(\mathrm{\Phi })+\xi \mathrm{\Phi }^{}\mathrm{\Phi })\frac{1}{4\pi }\frac{1}{4}F^2].$$
(1)
Here $`G`$ is the gravitational constant, $`D_\mu `$ is the usual gauge-covariant derivative expressed in the anti-hermitian basis of SU(2) ($`\tau _a=i\sigma _a/2`$)
$$D_\mu =_\mu +g\tau A_\mu ,$$
(2)
$`g`$ is the gauge coupling constant. Following , we assume that $`\mathrm{\Phi }`$ possesses only one degree of freedom
$$\mathrm{\Phi }=\frac{1}{\sqrt{2}}\left(\genfrac{}{}{0pt}{}{0}{\varphi (x)}\right),$$
(3)
with $`\varphi `$ real and time independent, and with the Higgs potential
$$V(\varphi )=\frac{\lambda }{4}(\varphi ^2v^2)^2,$$
(4)
where $`v`$ denotes the vacuum expectation value of $`\mathrm{\Phi }`$. The action (1) becomes
$$S=d^4x\sqrt{g}[\frac{}{16\pi G}\frac{1}{4\pi }\frac{1}{2}((_\mu \varphi )(^\mu \varphi )+(\frac{g\varphi }{2})^2A_\mu A^\mu +V(\varphi )+\xi \varphi ^2)\frac{1}{4\pi }\frac{1}{4}F^2].$$
(5)
There are considerable modifications in the Einstein equations due to the new energy-momentum tensor:
$`8\pi T_{\mu \nu }=`$ $`8\pi T_{\mu \nu }^{(minimal)}+2\xi (G_{\mu \nu }\varphi ^2+g_{\mu \nu }_\gamma ^\gamma \varphi ^2\varphi _{,\mu ;\nu }^2)`$ (6)
$`8\pi T_{\mu \nu }^{(minimal)}=`$ $`2F_{\mu \gamma }F_\nu ^\gamma {\displaystyle \frac{1}{2}}g_{\mu \nu }F^2+2({\displaystyle \frac{g\varphi }{2}})^2A_\mu A_\nu g_{\mu \nu }({\displaystyle \frac{g\varphi }{2}})^2A_\gamma A^\gamma `$ (7)
$`+2(_\mu \varphi )(_\nu \varphi )g_{\mu \nu }((_\gamma \varphi )(^\gamma \varphi )+2V(\varphi )).`$
where $`G_{\mu \nu }`$ is the Einstein tensor. As we assume spherical symmetry it is convenient to use the usual metric form:
$$ds^2=R^2(r)dr^2+r^2(d\theta ^2+sin^2\theta d\phi ^2)\frac{dt^2}{T^2(r)}$$
(8)
where $`R(r)=(12m(r)/r)^{1/2}`$ and $`m(r)`$ may be interpreted as the total mass-energy within the radius $`r`$. To describe the black-hole solutions we define $`\delta =ln(R/T)`$; thus:
$$ds^2=(1\frac{2m(r)}{r})^1dr^2+r^2(d\theta ^2+sin^2\theta d\phi ^2)(1\frac{2m(r)}{r})e^{2\delta (r)}dt^2.$$
(9)
The event horizon is at $`r=r_h`$ where $`1/R^2(r_h)=0`$. In case there are several such zeroes, the horizon corresponds to the outer one. Regularity at the origin is satisfied when T(0)$`<\mathrm{}`$ and $`R^{}(0)=T^{}(0)=0`$, while regularity at the event horizon r=$`r_h`$ requires $`\delta `$($`r_h)<\mathrm{}`$. In this paper we deal with nonextremal black-holes only, i.e. near the event horizon
$$12m(r)/rrr_h.$$
(10)
A suitable rescaling of the time coordinate $`t`$ implies:
$`R(0)=1,T(0)=1,m(r_h)=r_h/2,\delta (r_h)=0.`$ (11)
For the Yang-Mills field, it is convenient to use the ansatz discussed in ; thus a suitable parametrization of the Yang-Mills connection is:
$$A=\frac{1}{g}(1+\omega )[\widehat{\tau }_\phi d\theta +\widehat{\tau }_\theta \mathrm{sin}\theta d\phi ].$$
(12)
The $`\widehat{\tau }_i`$ are appropriately normalised spherical generators of the SU(2) group in the notation of ref. , $`e.g.`$ $`\widehat{\tau }_r=\widehat{r}\tau `$, $`[\tau _a,\tau _b]=ϵ_{abc}\tau _c`$, while $`\varphi (r)`$ is the Higgs field.
Expressing the curvature scalar $``$ in terms of the metric function $`R(r)`$ and $`T(r)`$, we obtain the following expression of the reduced action of our static spherically symmetric system:
$`S={\displaystyle }drdt[{\displaystyle \frac{1}{2G}}{\displaystyle \frac{1}{T}}(R{\displaystyle \frac{1}{R}}+2r{\displaystyle \frac{R^{}}{R}}){\displaystyle \frac{1}{2}}({\displaystyle \frac{(\varphi ^{})^2r^2}{RT}}+{\displaystyle \frac{\varphi ^2(1+\omega )^2}{2}}{\displaystyle \frac{R}{T}})`$
$`V(\varphi )r^2{\displaystyle \frac{R}{T}}{\displaystyle \frac{1}{g^2}}({\displaystyle \frac{(\omega ^{})^2}{RT}}+{\displaystyle \frac{(1\omega ^2)^2}{2r^2}}{\displaystyle \frac{R}{T}})\xi {\displaystyle \frac{\varphi ^2}{RT}}(R^21+2r{\displaystyle \frac{R^{}}{R}})+2\xi \varphi \varphi ^{}{\displaystyle \frac{r^2T^{}}{RT^2}}]`$ (13)
for a regular spacetime, while a suitable form of the reduced action for a black-hole spacetime is
$`S={\displaystyle }drdte^\delta [{\displaystyle \frac{m^{}(12\xi G\varphi ^2)}{G}}{\displaystyle \frac{1}{2}}((\varphi ^{})^2r^2(1{\displaystyle \frac{2m}{r}})+{\displaystyle \frac{\varphi ^2(1+\omega )^2}{2}})`$
$`V(\varphi )r^2{\displaystyle \frac{1}{g^2}}((\omega ^{})^2(1{\displaystyle \frac{2m}{r}})+{\displaystyle \frac{(1\omega ^2)^2}{2r^2}})+2\xi \varphi \varphi ^{}(r^2\delta ^{}(1{\displaystyle \frac{2m}{r}})+m^{}rm)]`$ (14)
where the prime denotes derivative with respect to r.
A usual rescaling
$`rrg/\sqrt{G},\varphi \varphi /\sqrt{G}`$ (15)
reveals that the solutions depend essentially on two dimensionless parameters $`\alpha `$ and $`\beta `$, expressible through the mass ratios
$`\alpha ={\displaystyle \frac{M_W}{vM_{Pl}}};\beta ={\displaystyle \frac{M_H}{M_W}}`$ (16)
with $`M_W=gv`$, $`M_H=\sqrt{\lambda }v`$ and $`M_{Pl}=\frac{1}{\sqrt{G}}`$ . $`V(\varphi )=\frac{\beta ^2}{4}(\varphi ^2\alpha ^2)^2`$ is the standard double well field potential. The field equations imply the relations
$$\omega ^{\prime \prime }=\omega ^{}(\frac{R}{R^{}}+\frac{T}{T^{}})+\frac{\omega (\omega ^21)}{r^2}R^2+\frac{\varphi ^2(\omega +1)}{4}R^2$$
(17)
for the gauge field, and
$$\varphi ^{\prime \prime }=\varphi ^{}(\frac{R}{R^{}}+\frac{T}{T^{}}\frac{2}{r})+R^2[\frac{\varphi (1+\omega )^2}{2r^2}+\frac{dV}{d\varphi }+\xi \varphi ]$$
(18)
for the Higgs field, where
$`={\displaystyle \frac{1}{\frac{1}{2}+(6\xi 1)\xi \varphi ^2}}((16\xi )(\varphi ^{})^2(1{\displaystyle \frac{2m}{r}})+4V(\varphi )+{\displaystyle \frac{\varphi ^2(1+\omega )^2}{2r^2}}`$
$`6\xi \varphi {\displaystyle \frac{dV(\varphi )}{d\varphi }}3\xi \varphi ^2{\displaystyle \frac{(1+\omega )^2}{r^2}})`$ (19)
is the spacetime curvature. The $`(rr)`$ and $`(tt)`$ Einstein equation are
$`(12\xi \varphi ^2)m^{}=(1{\displaystyle \frac{2m}{r}})({\displaystyle \frac{(\varphi ^{})^2r^2}{2}}+(\omega ^{})^22\xi \varphi \varphi ^{}r^2{\displaystyle \frac{T^{}}{T}}2\xi r^2(\varphi ^{})^2)+Vr^2+`$
$`{\displaystyle \frac{\varphi ^2(1+\omega )^2}{4}}+{\displaystyle \frac{(1\omega ^2)^2}{2r^2}}2\xi r^2(\varphi {\displaystyle \frac{dV}{d\varphi }}+{\displaystyle \frac{\varphi ^2(1+\omega )^2}{2r^2}}+\xi \varphi ^2),`$ (20)
$`(\xi \varphi \varphi ^{}r({\displaystyle \frac{1}{2}}\xi \varphi ^2))2r{\displaystyle \frac{T^{}}{T}}=({\displaystyle \frac{1}{2}}\xi \varphi ^2){\displaystyle \frac{2m}{r}}{\displaystyle \frac{1}{1\frac{2m}{r}}}+4\xi \varphi \varphi ^{}r+(\omega ^{})^2+{\displaystyle \frac{r^2(\varphi )^2}{2}}`$
$`{\displaystyle \frac{1}{(1\frac{2m}{r})}}(Vr^2+{\displaystyle \frac{\varphi ^2(1+\omega )^2}{4}}+{\displaystyle \frac{(1\omega ^2)^2}{2r^2}}).`$ (21)
For the black-hole solutions we replace the auxiliary $`T^{}`$ equation with an equation for $`\delta ^{}`$
$`r\delta ^{}(1+2\xi \varphi ^2)=a_1+{\displaystyle \frac{a_2}{1\frac{2m}{r}}}{\displaystyle \frac{\frac{1}{2}+\xi \varphi \varphi ^{}r\xi \varphi ^2}{\frac{1}{2}+\xi \varphi \varphi ^{}r+\xi \varphi ^2}}(a_3+{\displaystyle \frac{a_4}{1\frac{2m}{r}}})`$ (22)
where
$`a_1=`$ $`(\omega ^{})^2+{\displaystyle \frac{(\varphi ^{})^2r^2}{2}}2\xi r^2(\varphi ^{})^2,`$ (23)
$`a_2=`$ $`({\displaystyle \frac{1}{2}}\xi \varphi ^2){\displaystyle \frac{2m}{r}}+Vr^2+{\displaystyle \frac{\varphi ^2(1+\omega )^2}{4}}`$ (24)
$`+{\displaystyle \frac{(1\omega ^2)^2}{2r^2}}2\xi r^2(\varphi {\displaystyle \frac{dV}{d\varphi }}+{\displaystyle \frac{\varphi ^2(1+\omega )^2}{2r^2}}+\xi \varphi ^2),`$
$`a_3=`$ $`4\xi \varphi \varphi ^{}r+(\omega ^{})^2+{\displaystyle \frac{r^2(\varphi )^2}{2}},`$ (25)
$`a_4=`$ $`({\displaystyle \frac{1}{2}}\xi \varphi ^2){\displaystyle \frac{2m}{r}}(Vr^2+{\displaystyle \frac{\varphi ^2(1+\omega )^2}{4}}+{\displaystyle \frac{(1\omega ^2)^2}{2r^2}}).`$ (26)
Following the analysis in , we can already predict the boundary conditions and some general features of the finite energy solutions. Since the Yang-Mills equations are unaffected by the presence of the $`\xi \mathrm{\Phi }^2`$ term in (5), the analysis presented by Greene, Mathur and O’Neill remains valid: $`\omega =1`$ is the only acceptable value and $`\omega 1`$ is required for finite energy solutions. If we assume that $`\varphi `$ is $`O(\alpha )`$ in the region $`r1`$, we obtain:
$`rO(1/\alpha ):1<\omega <{\displaystyle \frac{1}{2}}(1+\sqrt{1(\alpha r)^2})`$
$`rO(1/\alpha ):\omega ^{}<0,\omega ^{\prime \prime }>0`$ (27)
These constraints are valid for both regular and black-hole solutions. The Higgs equation can be written in the form
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d}{dr}}({\displaystyle \frac{r^2(\varphi ^2)^{}}{RT}})={\displaystyle \frac{1}{\frac{1}{2}+(6\xi 1)\xi \varphi ^2}}[{\displaystyle \frac{1}{2}}{\displaystyle \frac{(\varphi ^{})^2r^2}{RT}}+{\displaystyle \frac{Rr^2}{T}}({\displaystyle \frac{\varphi }{2}}{\displaystyle \frac{dV}{d\varphi }}+{\displaystyle \frac{(1+\omega )^2\varphi ^2}{4r^2}})`$
$`+2\xi \varphi ^2r^2{\displaystyle \frac{R}{T}}(2V{\displaystyle \frac{\varphi }{2}}{\displaystyle \frac{dV}{d\varphi }})].`$ (28)
The obvious requirement for finite energy solution is
$`\varphi {\displaystyle \frac{dV}{d\varphi }}(12\xi \varphi ^2)<0`$ (29)
which implies that $`\varphi `$ is restricted to lie between the minima of the potential, $`\alpha \varphi \alpha `$. Relation (29) provides also an upper bound on the range of $`\xi `$ :
$$\xi <\frac{1}{2\alpha ^2}$$
(30)
It is worth noting, that an investigation of classical stability of a scalar field in a curved spacetime with a general coupling to gravity found, that the Higgs fields in the standard model must have $`\xi 0`$ or $`\xi 1/6`$.
¿From the relation (5) we can see that an effective gravitational constant is given by
$$G_{eff}=\frac{G}{12\xi \varphi ^2}.$$
(31)
Thus the condition (30) implies positivity of the effective gravitational constant. It is not possible to obtain an explicit lower bound for $`\xi `$ for a given value of $`\alpha `$. The field equations imply that $`\pm \alpha `$ are the only allowed values of $`\varphi `$ as $`r\mathrm{}`$. We focus here on solutions with $`\varphi (\mathrm{})=\alpha `$ without loss of generality. The vacuum values $`\omega (\mathrm{})=1`$ and $`\varphi (\mathrm{})=\alpha `$ are shared both by black-holes and regular solutions. The analysis of the field equations as $`r\mathrm{}`$ gives
$`m(r)`$ $`M+{\displaystyle \frac{1}{12\xi \alpha ^2}}(2a\sqrt{2}\xi \alpha ^2\beta cr^2e^{\sqrt{2}\alpha \beta cr}`$ (32)
$`{\displaystyle \frac{\alpha \beta }{2\sqrt{2}c}}((c^2+1)(14\xi )+{\displaystyle \frac{48\xi ^3\alpha ^2}{\frac{1}{2}+(6\xi 1)\xi \alpha ^2}}))a^2r^2e^{2\sqrt{2}\alpha \beta cr},`$
$`lnT(r)`$ $`ln(T_0)+{\displaystyle \frac{M}{r}},`$ (33)
$`\delta (r)`$ $`\delta _02\sqrt{2}{\displaystyle \frac{\xi \alpha ^2\beta c}{12\xi \alpha ^2}}are^{\sqrt{2}\alpha \beta cr}`$ (34)
$`+{\displaystyle \frac{\alpha \beta }{2\sqrt{2}c}}((c^2+1)(14\xi )+{\displaystyle \frac{48\xi ^3\alpha ^2}{\frac{1}{2}+(6\xi 1)\xi \alpha ^2}}){\displaystyle \frac{a^2re^{2\sqrt{2}\alpha \beta cr}}{12\xi \alpha ^2}},`$
$`\omega (r)`$ $`1+be^{\frac{\alpha r}{2}},`$ (35)
$`\varphi (r)`$ $`\alpha +ae^{\sqrt{2}\alpha \beta cr},`$ (36)
where $`c=\sqrt{\frac{12\xi \alpha ^2}{1+2\xi \alpha ^2(6\xi 1)}}`$; $`M,b,a`$ are constants; $`b>0,a<0`$.
Relation (32) implies an asymptotic violation of the weak energy condition (WEC) for negative values of $`\xi `$, since $`m^{}(r)<0`$. There exist also other classical field theories that violate the WEC. Examples are theories containing $`+^2`$ terms in the action , an antisymmetric 3-form axion field coupled to scalar fields , the Brans-Dicke scalar-tensor theory , and Einstein-dilaton theory with curvature-squared terms of Gauss-Bonnet type .
Since $`\mu (\xi )=\alpha \beta \sqrt{\frac{12\xi \alpha ^2}{1+2\xi \alpha ^2(6\xi 1)}}`$ corresponds to the mass of the Higgs field at infinity, the following relation holds
$$\mu (\xi )\mu (0)=\alpha \beta .$$
(37)
Thus any nonminimal coupling decreases the asymptotic value of the Higgs field mass. To get further insight into the meaning of a large value of $`\xi `$ it is worthwhile to consider the rescaling $`rr/\sqrt{(\xi )}`$, $`\varphi \varphi /\sqrt{(\xi )}`$. For $`\xi \mathrm{}`$ we find that $`\alpha 0`$ necessarily. Thus, for a large negative $`\xi `$, we expect a decrease of the maximal allowed value of the parameter $`\alpha `$. Furthermore, there is an effective decoupling of the Yang-Mills and gravitational fields and the effective coupling of the Higgs field to matter becomes of gravitational strength. In the limit of infinite negative $`\xi `$, the following field equations are obtained
$`({\displaystyle \frac{1}{2}}+\varphi ^2)2r{\displaystyle \frac{R^{}}{R}}=({\displaystyle \frac{1}{2}}+\varphi ^2)(1R^2)+2\varphi \varphi ^{}r^2{\displaystyle \frac{T^{}}{T}},`$ (38)
$`(\varphi \varphi ^{}r+{\displaystyle \frac{1}{2}}+\varphi ^2)2r{\displaystyle \frac{T^{}}{T}}=({\displaystyle \frac{1}{2}}+\varphi ^2)(1R^2)+4\varphi \varphi ^{}r,`$ (39)
$`\omega ^{\prime \prime }=\omega ^{}({\displaystyle \frac{R^{}}{R}}+{\displaystyle \frac{T^{}}{T}})+{\displaystyle \frac{\omega (\omega ^21)R^2}{r^2}}+{\displaystyle \frac{\varphi ^2(1+\omega )^2}{4}}R^2,`$ (40)
$`\varphi ^{\prime \prime }=\varphi ^{}({\displaystyle \frac{R^{}}{R}}+{\displaystyle \frac{T^{}}{T}}{\displaystyle \frac{2}{r}}){\displaystyle \frac{(\varphi ^{})^2}{\varphi }}.`$ (41)
By using a usual power series expansion near the origin or the event horizon it can be proven that there are no initial conditions consistent with the requirement of energy finiteness. However a simpler proof is to observe that equation (41) implies the relation $`(\varphi ^2)^{}=const.\frac{RT}{r^2}`$, which is also consistent with the general equation (2). There are no nonsingular solution of this equation consistent with the requirement of metric regularity at the origin or with a regular event horizon. Thus we conclude that nontrivial solutions are absent in the case of an infinite negative $`\xi `$. In practice, it becomes increasingly difficult to solve the field equations for large negative values of $`\xi `$, with a fast convergence to the asymptotic values $`\omega (\mathrm{})`$, $`\varphi (\mathrm{})`$.
A particularly interesting case of the general theory is obtained for $`G\mathrm{}`$, i.e. in the absence of the Einstein term in the action (5), corresponding to the spontaneous symmetry breaking theory of gravity, with the standard Higgs field as the origin of the Plank mass (for the remainder of this paragraph we do not consider the rescaling (15)).
There has recently been an increased interest in induced gravity in the standard model, that might help to solve some problems of particle physics and cosmology. Typical problems are the necessity of the Higgs mass to be order of the theory cut-off , the missing mass problem, Mach’s principle , and the inflationary scenario . The existence of magnetic monopole solutions in an induced gravity YMH theory has also been discussed . Unfortunately, it can be proven that, in the absence of the Einstein term, only the trivial case $`\varphi =\pm \alpha `$ is consistent with the requirement of energy finiteness. When we consider the equation (18) and use the trace of Einstein equations to eliminate the term $`\xi \varphi `$ we obtain the general equation
$$\frac{1}{2}_\mu ^\mu \varphi ^2=\frac{1}{16\xi }(\frac{}{2G}+\varphi \frac{dV}{d\varphi }4V).$$
(42)
For our ansatz we obtain the relation
$$\frac{d}{dr}(\frac{r^2(\varphi ^2)^{}}{RT})=\frac{Rr^2}{T}\frac{1}{16\xi }(\frac{}{2G}+\varphi \frac{dV}{d\varphi }4V).$$
(43)
which should be satisfied for all $`r`$. Since clearly $`4V\varphi \frac{dV}{d\varphi }>0`$ for the considered potential, it follows that in the absence of an Einstein term in the original action only $`\varphi =\pm \alpha `$ is consistent with the assumption of finite energy, both for regular and black-hole solutions. However, for $`\varphi =\pm \alpha `$ we obtain the Bartnik-McKinnon solutions and their black-hole generalizations; there are no spherically symmetric gravitating sphaleron or sphaleron black-hole solutions. One can conjecture that similar to the boson star case, it may be possible to obtain nontrivial solutions by considering a time dependence of the matter field.
## 3 NUMERICAL SOLUTION
Nontrivial solutions are not known in closed form, and so a numerical method of solution is necessary.
### 3.1 REGULAR SOLUTIONS
For regular solutions, finite $`T_{tt}`$ and regularity of the metric at $`r=0`$ give two possible sets of initial conditions
$`2m(r)=O(r^3)`$ (44)
$`lnT(r)=O(r^2)`$ (45)
$`\left({\displaystyle \genfrac{}{}{0pt}{}{\omega (r)}{\varphi (r)}}\right)=\left({\displaystyle \genfrac{}{}{0pt}{}{1+O(r^2)}{\varphi _0+O(r^2)}}\right),`$ (46)
or
$`\left({\displaystyle \genfrac{}{}{0pt}{}{\omega (r)}{\varphi (r)}}\right)=\left({\displaystyle \genfrac{}{}{0pt}{}{1+O(r^2)}{O(r)}}\right).`$ (47)
The general properties of the solutions are the same as for the minimally-coupled EYMH theory. Solutions are again characterized by $`w(r)`$ oscillations in the region $`r>1`$ and classified by the node number k which may be even or odd. The formal power series describing the above boundary conditions at $`r=0`$ is
$`2m(r)=ar^3+O(r^5),`$ (48)
$`lnT(r)={\displaystyle \frac{1}{2}}fr^2+O(r^4),`$ (49)
$`\omega =1+br^2+O(r^4),`$ (50)
$`\varphi =\varphi _0+er^2+O(r^4),`$ (51)
with
$`a={\displaystyle \frac{4b^2+\frac{2}{3}V_0}{12\xi \varphi _0^2}}{\displaystyle \frac{4}{3}}{\displaystyle \frac{\xi \varphi _0}{12\xi \varphi _0^2}}{\displaystyle \frac{(\frac{1}{2}\xi \varphi _0^2)V_0^{}+4V_0\varphi _0\xi }{\frac{1}{2}+(6\xi 1)\xi \varphi _0^2}},`$ (52)
$`f={\displaystyle \frac{1}{2}}{\displaystyle \frac{4b^2+\frac{2}{3}V_0}{12\xi \varphi _0}}+{\displaystyle \frac{2}{3}}{\displaystyle \frac{\xi \varphi _0}{12\xi \varphi _0^2}}{\displaystyle \frac{(\frac{1}{2}\xi \varphi _0^2)V_0^{}+4V_0\varphi _0\xi }{\frac{1}{2}+(6\xi 1)\xi \varphi _0^2}}`$
$`+{\displaystyle \frac{8\xi \varphi _0eV_0+2b^2}{2\xi \varphi _0^21}},`$ (53)
$`e={\displaystyle \frac{1}{6}}V_0^{}+{\displaystyle \frac{\xi \varphi _0}{\frac{1}{2}+(6\xi 1)\xi \varphi _0^2}}(4V_06\xi \varphi _0V_0^{}),`$ (54)
for even-k solutions ($`V_0,V_0^{}`$ are the potential and its derivative with respect to $`\varphi `$ at $`\varphi =\varphi _0`$) and
$`2m(r)=(4b^2+{\displaystyle \frac{2}{3}}V_0+e^24\xi e^2)r^3+O(r^5),`$ (55)
$`lnT(r)=(2b^2{\displaystyle \frac{1}{3}}V_0+2\xi e^2)r^2+O(r^4),`$ (56)
$`\omega =1br^2+O(r^4),`$ (57)
$`\varphi =er+O(r^3),`$ (58)
for odd-k solutions. The shooting parameters are $`(\varphi _0,b)`$ and $`(e,b)`$ respectively. Using a standard ordinary differential equation solver, we evaluate the initial conditions at $`r=10^3`$ for global tolerance $`10^{12}`$, adjusting for fixed shooting parameters and integrating towards $`r\mathrm{}`$. The difficulty of the two-dimensional shooting problem in the presence of two free parameters is increased by the presence of a nonminimal coupling between the Higgs field and gravity which leads to a slow convergence of the mass function $`m(r)`$.
We limit the discussion to results containing only one or two nodes. The results obtained for the $`k=1`$ and $`k=2`$ solutions retain the general characteristics of the minimally coupled case. In order to define the terminology, which is somewhat confusing, we review here the $`\xi =0`$ case ().
In the $`k=1`$ case there are two possible solutions that are called the quasi k=0 solution and the proper $`k=1`$ solution. The different character of the two branches becomes apparent in the limit $`\alpha 0`$. Remembering that $`\alpha =v\sqrt{G}`$, there are two physically inequivalent ways for $`\alpha `$ to approach $`0`$. The first is for the Newton constant to vanish, the second for the Higgs vacuum expectation value to vanish. In the first case we have weakly coupled gravity and the solution approaches the standard model sphaleron. The quasi-$`k=`$0 branch of solutions has this limit for $`\alpha 0`$. For $`\alpha 0`$ the node that is present actually moves to infinity in this branch, therefore the name quasi-$`k=`$0. The proper $`k=1`$ branch approaches the first Bartnik-McKinnon solution when $`\alpha 0`$. It corresponds to taking $`v0`$, but always having gravity present. For small enough $`\alpha `$, depending on $`\beta `$, both solutions are present and different. For finite $`\alpha `$ the behaviour is dependent on the value of $`\beta `$. For $`\beta `$ larger than a critical value $`\beta _{crit}0.12`$ there is a first maximum value of $`\alpha `$, where the proper $`k=1`$ branch disappears. For a larger value of $`\alpha `$ also the quasi-$`k=0`$ solution disappears. In the case $`\beta `$ smaller than $`\beta _{crit}`$ the maximum value of $`\alpha `$ is the same for both branches. Here the solutions merge; the shooting parameters approach each other. The situation is clarified in $`figure`$ 1.
For the case of two nodes the situation is similar. There are again two branches, called the quasi$`k=1`$ solution and the proper $`k=2`$ solution. The quasi-$`k=1`$ solution approaches the first Bartnik-McKinnon solution in the limit $`\alpha 0`$, one of the nodes moving to infinity. The proper $`k=2`$ solution approaches the second Bartnik-McKinnon solution. The behaviour as a function of $`\alpha `$ and $`\beta `$ is qualitatively similar to the one-node case ($`figure`$ 1).
To compare numerically the results with those found in we focused on solutions with $`\beta ^2=1/8`$ (although similar results have been obtained for other choices of $`\beta `$) and with $`k=1`$ and $`k=2`$ only. The results of the numerical integration for $`\alpha `$=0.1, $`\beta ^2`$=1/8 and a range of $`\xi `$ are presented in $`figure`$ 2.
In this figure $`\xi `$ is relatively small (i.e. $`\xi \frac{1}{2\alpha ^2}`$) . As a consequence the results of remain approximately valid. The correction to the shooting parameters is very small. A general feature of the solutions is the small influence of the term $`\xi \mathrm{\Phi }^2`$ on the value of the ADM mass. A negative value of $`\xi `$ seems to decrease the asymptotic value of the $`m(r)`$ function while a positive $`\xi `$ determines a higher ADM mass with respect to the minimally coupled case. The effect is particularly small for the quasi-$`k=0`$ solution. This is understandable as this solution corresponds to the flat space sphaleron, for which gravity is a small effect. For the considered range of $`\xi `$, a nonminimal coupling term has a small effect on the shape of the $`w(r)`$ function. However, for suitable values of $`\xi `$ we have noticed a strong influence of this term on the behavior of the Higgs field in the intermediate region. Also, for k=2 solutions, the initial value of the Higgs field $`\varphi _0`$ is strongly dependent on the value of $`\xi `$, with $`\varphi _0(\xi 0)<\varphi _0(\xi =0)`$ ($`figure`$ 2 c, d) . As we expected from (32), for negative $`\xi `$ we obtain a violation of the WEC beyond a certain limit of the radial coordinate, corresponding to a peak of the mass function $`m(r)`$. The height of the peak is proportional to the absolute value of $`\xi `$.
In $`figure`$3 we study the solutions as a function of $`\alpha `$. Significant changes occur for $`\xi \frac{1}{2\alpha ^2}`$ and for large negative values of $`\xi `$. The parameter range obtained in for the two sheets of solutions does not remain valid; it is necessary to establish a different value of $`\alpha _{max}`$ for every choice of $`\xi `$. A general feature for a nonminimal coupling is the decrease of the maximal allowed value of the parameter $`\alpha `$ . For example for the proper $`k=1`$ branch, Greene, Mathur and O’Neill have found $`0<\alpha <0.599`$; for $`\xi =0.1`$ we have obtained $`0<\alpha <0.48`$, while for $`\xi =2`$ the limiting value of the parameter $`\alpha `$ is 0.51 ($`figure`$ 3a, b). The quasi-k=0 branch with $`\xi =0`$ has $`0<\alpha <0.619`$; for $`\xi =1`$ we have found $`\alpha _{max}=0.616`$ while for $`\xi =1/6`$ we have $`0<\alpha <0.525`$ ($`figure`$ 3c, d). A minimally coupled even $`k`$ configuration has $`0<\alpha <0.120`$ (proper $`k=2`$ branch) or $`0<\alpha <0.122`$ (quasi-$`k=1`$ branch); in these cases, for $`\xi =1`$ or $`\xi =1/6`$ we did not notice a significant deviation of the $`\alpha _{max}`$ value ($`figure`$ 3 e-h). Different limiting values occur for the shooting parameters $`b,\varphi _0`$ and $`e`$ also.
### 3.2 BLACK-HOLE SOLUTIONS
Similar results can be obtained for numerical black-hole solutions. We use the following expansion near the event horizon:
$`m(r)={\displaystyle \frac{r_h}{2}}+m^{}(r_h)(rr_h),`$ (59)
$`\delta (r)=0+\delta ^{}(r_h)(rr_h),`$ (60)
$`\omega (r)=\omega (r_h)+\omega ^{}(r_h)(rr_h),`$ (61)
$`\varphi (r)=\varphi (r_h)+\varphi ^{}(r_h)(rr_h),`$ (62)
with
$`m^{}(r_h)={\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{a_2(r_h)a_4(r_h)}{12\xi \varphi ^2(r_h)}},`$ (63)
$`\omega ^{}(r_h)={\displaystyle \frac{r_h(12\xi \varphi ^2(r_h))}{a_2(r_h)a_4(r_h)}}({\displaystyle \frac{\omega (r_h)(\omega ^2(r_h)1)}{r_h^2}}`$
$`+{\displaystyle \frac{\varphi ^2(r_h)(1+\omega (r_h))}{4}}),`$ (64)
$`\varphi ^{}(r_h)={\displaystyle \frac{12\xi \varphi ^2(r_h)}{2\xi \varphi (r_h)r_h}}{\displaystyle \frac{a_2(r_h)+a_4(r_h)}{a_2(r_h)a_4(r_h)}},`$ (65)
$`\delta ^{}(r_h)={\displaystyle \frac{1}{r_h(12\xi \varphi ^2(r_h))}}({\displaystyle \frac{f^{}(r_h)r_h}{12m^{}(r_h)}}{\displaystyle \frac{1}{\frac{1}{2}+\xi \varphi ^2(r_h)+\xi \varphi (r_h)\varphi ^{}(r_h)r_h}}`$
$`+{\displaystyle \frac{(a_1(r_h)+a_3(r_h))(\frac{1}{2}\xi \varphi ^2(r_h))\xi \varphi (r_h)\varphi ^{}(r_h)r_h(a_1(r_h)a_3(r_h))}{\frac{1}{2}+\xi \varphi ^2(r_h)+\xi \varphi (r_h)\varphi ^{}(r_h)r_h}})`$ (66)
where
$`f(r)=a_2(r)+a_4(r))({\displaystyle \frac{1}{2}}\xi \varphi ^2)\xi \varphi \varphi ^{}r(a_2(r)a_4(r)).`$ (67)
The new shooting parameters are $`\omega (r_h)`$ and $`\varphi (r_h)`$ (we have studied the case $`r_h=1`$ only).
The non-minimal gravitational coupling allows for a not necessarily positive field energy. Therefore one loses one of the earlier tools for proving the no hair theorems, which already failed for the minimally coupled EYMH system. The bypassing of the usual no-hair theorems in the considered system can be proven by using the method of for the $`\xi =0`$ case.
Starting from the solutions (59-62) we integrated the system (2, 22, 17, 18) towards $`r\mathrm{}`$ using an automatic step procedure and accuracy $`10^{12}`$. The integration stops when the flat spacetime asymptotic limit (32, 33, 34, 35) is reached.
The behaviour of the black-hole solutions as a function of $`\xi `$ and $`\alpha `$ is similar to the regular solutions. Two solution branches appear for each $`k`$ corresponding to two different values of the shooting parameters $`\omega (r_h)`$ and $`\varphi (r_h)`$. As $`\alpha 0`$, the proper $`k=1,2`$ branches approaches the corresponding Einstein-Yang-Mills black-hole solutions ().
In the same limit, the quasi-$`k=0`$ branch is distinguished by its Schwarzschild solution limit ($`\omega =1`$, $`\varphi =0`$) and the last node of the quasi-$`k=0`$ and quasi-$`k=1`$ branches is again pushed out to infinity. The corresponding limit of the quasi-$`k=1`$ branch is the n=1 Einstein-Yang-Mills black-hole solution.
Again, every branch the solutions exist only for a finite range of the parameter $`0\alpha \alpha _{max}(\beta ,k)`$ with different values of $`\alpha _{max}`$ for every solution branch.
The results for $`k=`$1, 2, $`\beta ^2`$=1/8 and various values of the parameter $`\xi `$ are presented in $`figure`$ 4. As we expected, for a nonzero $`\xi `$ it is necesary to establish new limiting values of the values of the normalised vacuum expectation values $`\alpha `$. For example, a minimally coupled solution has necessarily $`0<\alpha <0.331`$ (proper $`k=1`$ branch), $`0<\alpha <0.356`$ (quasi-$`k=0`$ branch), $`0<\alpha <0.0475`$ (proper $`k=2`$ branch) and $`0<\alpha <0.0486`$ for the quasi-$`k=2`$ branch.
We have found $`0<\alpha <0.325`$ ($`\xi =1`$) and $`0<\alpha <0.352`$ for $`\xi =1/6`$ (proper $`k=1`$ branch); some results for the quasi-$`k=0`$ branch are: $`0<\alpha <0.356`$ ($`\xi =0.1`$), $`0<\alpha <0.352`$ ($`\xi =1/6`$). For the proper $`k=2`$ branch we have found $`0<\alpha <0.0382`$ ($`\xi =60`$), $`0<\alpha <0.336`$ ($`\xi =60`$), while the quasi-$`k=0`$ branch with $`\xi =10`$ has the limiting value $`\alpha <0.044`$; for a negative coupling constant $`\xi =5`$ we have found $`0<\alpha <0.047`$ (see also $`figure`$ 5). Different ranges for the shooting parameters $`\omega (r_h)`$ and $`\varphi (r_h)`$ are to be imposed.
Similar to the case of regular solutions, we notice the occurence of negative energy densities. Anyway, an unexpected feature is the violation of the WEC even in the vicinity of the event horizon (for positive values of $`\xi `$ and quasi$`k=0`$ branch), which is supposed to destabilize the black-hole and to lead to a traversable wormhole .
Another interesting problem is the effect of the nonminimal coupling on the properties of a black-hole. Not suprisingly, for the quasi$`k=0`$ branch we have noticed a violation of the generic relation
$$T_H=\frac{1}{4\pi r_H}e^{\delta (r_H)}(12m^{}(r_h))\frac{1}{4\pi r_h}=T_H^{vac}.$$
(we use units $`k_B=\mathrm{}`$=1) where $`T_H^{vac}`$ is the Hawking temperature of a Schwarzchild black-hole with the same area.
We have found that generally a positive $`\xi `$ will increase the value of the Hawking temperature (the only disturbing exception is the quasi $`k=1`$ case). However, following the validity of the generic relation $`S=\frac{A}{4}`$ (where $`A`$ is the event horizon area) can easily be proven for the Lagrangian density (5).
## 4 FURTHER DISCUSSION
Further insight into the meaning of the nonminimal coupling in EYMH theory can be obtained by using the conformal rescaling of the action (5):
$`\overline{g}_{\mu \nu }=\mathrm{\Omega }^2g_{\mu \nu },\mathrm{\Omega }^2=12\xi G\varphi ^2.`$ (68)
The use of this conformal transformation together with a redefinition of the scalar field for the case of nonminimal coupling has a long history; ref. , for instance, presents a large set of references on this subject. The usual condition $`\mathrm{\Omega }^2>0`$ has a clear physical meaning since it is satisfied by finite energy solutions only. The pairs of variables (metric $`g_{\mu \nu }`$, scalar $`\varphi `$, $`SU(2)`$ field $`F_{\mu \nu }`$) defined originally constitute what is called a Jordan frame. Consider now the transformation
$`\psi ={\displaystyle 𝑑\varphi F(\varphi )},`$
$`F^2(\varphi )={\displaystyle \frac{12\xi G\varphi ^2(16\xi )}{(12\xi G\varphi ^2)^2}}`$ (69)
such that, in the redefined action
$`S={\displaystyle }d^4x\sqrt{\overline{g}}[{\displaystyle \frac{\overline{}}{16\pi G}}{\displaystyle \frac{1}{4\pi }}{\displaystyle \frac{1}{2}}(_\mu \psi )(^\mu \psi )`$
$`{\displaystyle \frac{1}{4\pi }}{\displaystyle \frac{1}{2}}({\displaystyle \frac{g\varphi }{2}})^2{\displaystyle \frac{A^2}{12\xi G\varphi ^2}}{\displaystyle \frac{1}{4\pi }}\overline{V}(\psi ){\displaystyle \frac{1}{4\pi }}{\displaystyle \frac{1}{4}}F^2]`$ (70)
$`\psi `$ becomes minimally coupled to $`\overline{}`$, with
$$\overline{V}(\psi )=\frac{V(\varphi )}{\mathrm{\Omega }^4}.$$
(71)
The new variables (metric $`\overline{g}_{\mu \nu }`$, scalar $`\psi `$, $`SU(2)`$ field $`F_{\mu \nu }`$) are said to constitute an Einstein frame. The transformation given by eqs. (68, 4) therefore maps a solution of the field equations imposed by (5) to a solution that extremizes (4). The transformation is independent of any assumption of symmetry, and in this sense is covariant; one can easily infer that the transformation is one-to-one in general. Also, the transformation preserves symmetries, which means that if $`g_{\mu \nu }`$ admits a Killing vector $`\eta `$ such that $`\mathrm{\pounds }_\eta \varphi =0`$ , then $`\eta `$ is also a Killing vector of $`\overline{g}_{\mu \nu }`$ and $`\mathrm{\pounds }_\eta \psi =0`$. There is a long debate in the literature on the problem of which of these two frames is physical (for a review see ). For example, in ref. it has been shown in a more general context that all thermodynamical variables defined in the original frame are the same as those in the Einstein frame, if spacetimes in both frames are asymptotically flat, regular and posses event horizons with non-zero temperature. We know that $`\mathrm{\Omega }^2`$ goes to some finite positive value at infinity. Since this value is not unity, the asymptotically Minkowskian metric $`g_{\mu \nu }`$ will be mapped into a generally non-asymptotically Minkowskian line element $`\overline{g}_{\mu \nu }`$. However, one needs only to redefine globally the units of length and time to obtain an Einstein-frame standard Minkowski form at infinity. Considering an expansion of the Higgs field $`\varphi `$ around the minimum $`\varphi =v+\eta `$, for large enough negative values of $`\xi `$ we obtain the following first order Einstein frame expression
$$L_{YM}=\frac{1}{4\pi }\frac{1}{4}F^2+\frac{1}{4\pi }\frac{1}{4\xi G}(\frac{g}{4})^2A^2,$$
(72)
with an effective decoupling of the YM and Higgs fields and a massive Yang-Mills theory, along the line suggested in .
The Weyl rescaling (68) helps us to rule out the existence of traversable wormhole solutions, since one can conclude that when we know all Einstein-frame solutions with a given symmetry we automatically know all Jordan-frame solutions with the same symmetry.
A spacetime wormhole is usually introduced as a topological handle connecting two universes or distant places in the same universe. Over the last decade following the seminal papers of Morris, Thorne and Yurtsever , considerable interest has grown in the domain of traversable wormhole physics (for a review see ). We recall that a Lorentzian wormhole solution is said to be traversable if it does not contain horizons that prevent the crossing of the throat. A remarkable result is that, assuming Einstein gravity, the WEC is violated at throat of a traversible static wormhole .
Since we have found that the violation of the WEC is possible, it is natural to look for spherically symmetric, traversable wormhole solutions of the coupled EYMH equations. Further, it has been conjectured that a violation of the WEC in the vicinity of the event horizon is quite likely to destabilize the horizon and lead to a traversable wormhole . Thus, we suppose the existence of a traversable wormhole solution in the original Jordan frame, therefore the violation of the energy condition at or near the throat of the wormhole. The case of a static spherically symmetric Lorentzian wormhole corresponds to the following choice of the metric functions in the general ansatz (8)
$`R(r)=(1b(r)/r)^{1/2},T(r)=e^{f(r)},`$ (73)
where $`b(r)`$ is called the shape function as it describes the shape of the spatial geometry of the wormhole in an embedding diagram and $`f(r)`$ describes the gravitational redshift in this spacetime (with $`e^{f(r)}>0`$) . In this case the coordinate $`r`$ is constrained to run between $`r_0<r<\mathrm{}`$, where $`r_0`$ is the throat radius ($`b(r_0)=r_0`$). By following the analysis of Sec. 2, we can again predict the general features of the possible solutions and boundary conditions. It follows that the general relations obtained as $`r\mathrm{}`$ are still valid. If we do not allow for $`\varphi `$ or $`\varphi ^{}`$ to take an infinite value at the wormhole throat, a similar analysis of the Higgs field equation (2) implies that the generic conditions $`\mathrm{\Omega }^2>0`$, $`\xi <\frac{1}{2\alpha ^2}`$ hold also .
Using the conformal transformation (68) we convert the theory to the Einstein frame. The existence of wormhole solutions is not affected by the transformation (68) that preserves the traversability for $`\mathrm{\Omega }^2>0`$, i.e. a positive effective gravitational constant and no event horizon in the new frame. It can easily be proven that for the rescaled action (4), the dominant energy condition holds, and thus there are no traversable wormhole solutions, i.e. no traversable wormhole solutions in the Jordan frame also.
One can wonder whether this absence of traversable wormhole solutions is a general feature of arbitrary nonminimal scalar couplings to Einstein gravity. The nonminimal coupling of the scalar field considered in this paper is a particular case of a more general theory, where the term $`\frac{}{16\pi G}(12\xi G\varphi ^2)`$ is replaced by a more general function $`\frac{}{16\pi G}f(\varphi )`$.
By using a conformal transformation $`\overline{g}_{\mu \nu }=f(\varphi )g_{\mu \nu }`$ and a redefinition of the scalar field ()
$`\psi ={\displaystyle 𝑑\varphi (\frac{f(\varphi )+\frac{3}{4G}(\frac{df(\varphi )}{d\varphi })^2}{f(\varphi )^2})^{1/2}}`$ (74)
we can convert the action to the Einstein frame
$`S={\displaystyle }d^4x\sqrt{\overline{g}}[{\displaystyle \frac{\overline{}}{16\pi G}}{\displaystyle \frac{1}{4\pi }}{\displaystyle \frac{1}{2}}(_\mu \psi )(^\mu \psi )`$
$`{\displaystyle \frac{1}{4\pi }}{\displaystyle \frac{1}{2}}({\displaystyle \frac{g\varphi }{2}})^2{\displaystyle \frac{A^2}{f(\varphi )}}{\displaystyle \frac{1}{4\pi }}\overline{V}(\psi ){\displaystyle \frac{1}{4\pi }}{\displaystyle \frac{1}{4}}F^2]`$ (75)
(with $`\overline{V}(\psi )=\frac{V(\varphi )}{f(\varphi )^2}`$).
If $`f(\varphi )>0`$ (i.e. a positive effective Newton constant) the WEC will be satisfied in the Einstein frame. The positivity of the effective Newton constant followed in our case from the demand that the energy of the solution is finite; implying $`\xi <\frac{1}{2\alpha ^2}`$. We have not been able to show that an equivalent condition holds in the general case. The possibility of traversable wormholes in general is therefore left open, though it seems likely they will be absent.
## 5 CONCLUSIONS
In this paper we have studied static, spherically symmetric classical solutions of spontaneously broken SU(2) gauge theory with a nonminimally coupled Higgs field and presented strong numerical arguments for the existence of both regular and black-hole solutions for suitable values of the coupling parameter $`\xi `$. The main properties of these solutions, such as their nodal structure and discrete mass spectrum, are generic and shared by practically all known solitons with gravitating non-abelian gauge fields.
It should be stressed that it is the non-abelian nature of the Yang-Mills field that allows the existence of nontrivial solutions since a nonminimally coupled Higgs charged hair has been ruled out for Abelian Higgs theory (, see also ).
The absence of gravitating sphalerons and sphaleron black-holes in a spontaneously broken theory of gravity has also been proven. As a new feature we have established a violation of the WEC for a certain range of $`\xi `$. The nonexistence of traversable wormhole solutions has been shown using a conformal map to convert the problem to the one with minimal coupling to gravity. For small values of the parameter $`\xi `$, the effect of the nonminimal coupling on the asymptotic features of a finite energy solution is rather benign. Although we did not address the nature of the solutions for $`r<r_h`$, we expect that a term $`\xi \mathrm{\Phi }^2`$ can strongly influence the properties of the inner black-hole solutions.
We have not considered the question of stability of solutions in this paper. Since in the case of minimal coupling the solutions were found to be unstable, we see no reason to expect something special to happen for $`\xi 0`$.
Acknowledgement
This work was performed in the context of the Graduiertenkolleg of the Deutsche Forschungsgemeinschaft (DFG): Nichtlineare Differentialgleichungen: Modellierung,Theorie, Numerik, Visualisierung.
Figure Captions
Figure 1: The parameter beta versus the maximal value of the parameter $`\alpha `$ for one node graviting sphaleron solutions; qualitative picture for the minimally coupled case.
Figure 2: One- and two-node sphaleron solutions of the nonminimally coupled EYMH theory for $`\alpha =0.1`$, $`\beta ^2=1/8`$ and various values of $`\xi `$.
Figure 3: One- and two-node sphaleron solutions of the nonminimally coupled EYMH theory for $`\beta ^2=1/8`$ and various values of $`\xi `$. Here and in $`figure`$ 5, the parameter $`\alpha `$ varies between zero and the maximum allowed value $`\alpha _{max}`$; increasing $`\alpha `$ corresponds to a decrease of the value of the radial coordinate at which the solution exponentially decays to its vacuum value.
Figure 4: One- and two-node black hole solutions of the nonminimally coupled EYMH theory for $`\beta ^2=1/8`$ and various values of $`\xi `$.
Figure 5: One- and two-node black hole solutions of the nonminimally coupled EYMH theory for $`\beta ^2=1/8`$ and various values of $`\xi `$.
Figure 1.
Figure 2a. Proper $`k=1`$ Regular $`\xi =7`$, 0, -10
Figure 2b. Quasi$`k=0`$ Regular $`\xi =0.7`$, 0, -7
Figure 2c. Proper $`k=2`$ Regular $`\xi =2.5`$, 0, -4.5
Figure 2d. Quasi$`k=1`$ Regular $`\xi =5`$, 0, -10
Figure 3a. Proper $`k=1`$ Regular $`\xi =0.1`$; $`\alpha =0.005`$, 0.25, 0.47
Figure 3b. Proper $`k=1`$ Regular $`\xi =2`$; $`\alpha =0.005`$, 0.2, 0.5
Figure 3c. Quasi$`k=0`$ Regular $`\xi =1/6`$; $`\alpha =0.005`$, 0.2, 0.5
Figure 3d. Quasi$`k=0`$ Regular $`\xi =1`$; $`\alpha =0.05`$, 0.2, 0.615
Figure 3e. Proper $`k=2`$ Regular $`\xi =1/6`$; $`\alpha =0.005`$, 0.07, 0.121
Figure 3f. Proper $`k=2`$ Regular $`\xi =1`$; $`\alpha =0.001`$, 0.01, 0.12
Figure 3g. Quasi$`k=1`$ Regular $`\xi =1/6`$; $`\alpha =0.01`$, 0.1, 0.121
Figure 3h. Quasi$`k=1`$ Regular $`\xi =1`$; $`\alpha =0.005`$, 0.05, 0.12
Figure 4a. Proper $`k=1`$ Black Hole; $`\alpha =0.1`$; $`\xi =7`$, 0, -7
Figure 4b. Quasi$`k=0`$ Black Hole; $`\alpha =0.1`$; $`\xi =7`$, 0, -6
Figure 4c. Proper $`k=2`$ Black Hole; $`\alpha =0.01`$; $`\xi =60`$, 0, -200
Figure 4d. Quasi$`k=1`$ Black Hole; $`\alpha =0.01`$; $`\xi =20`$, 0, -20
Figure 5a. Proper $`k=1`$ Black Hole $`\xi =1/6`$; $`\alpha =0.005`$, 0.15, 0.351
Figure 5b. Proper $`k=1`$ Black Hole $`\xi =1`$; $`\alpha =0.005`$, 0.15, 0.31
Figure 5c. Quasi$`k=0`$ Black Hole $`\xi =1/6`$; $`\alpha =0.005`$, 0.15, 0.35
Figure 5d. Quasi$`k=0`$ Black Hole $`\xi =0.1`$; $`\alpha =0.005`$, 0.2, 0.356
Figure 5e. Proper $`k=2`$ Black Hole $`\xi =60`$; $`\alpha =0.005`$, 0.02, 0.038
Figure 5f. Proper $`k=2`$ Black Hole $`\xi =60`$; $`\alpha =0.005`$, 0.02, 0.033
Figure 5g. Quasi-$`k=1`$ Black Hole $`\xi =10`$; $`\alpha =0.005`$, 0.02, 0.042
Figure 5h. Quasi-$`k=1`$ Black Hole $`\xi =5`$; $`\alpha =0.005`$, 0.02, 0.046
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# 1 Problem
## 1 Problem
In the usual formulation of the Kirchhoff diffraction integral, a scalar field with harmonic time dependence at frequency $`\omega `$ is deduced at the interior of a charge-free volume from knowledge of the field (or its normal derivative) on the bounding surface. In particular, the field is propagated forwards in time from the boundary to the desired observation point.
Construct a time-reversed version of the Kirchhoff integral in which the knowledge of the field on the boundary is propagated backwards in time into the interior of the volume.
Consider the example of an optical focus at the origin for a system with the $`z`$ axis as the optic axis. In the far field beyond the focus a Gaussian beam has cone angle $`\theta _0\sqrt{2}\sigma _\theta `$, and the $`x`$ component of the electric field in a spherical coordinate system is given approximately by
$$E_x(r,\theta ,\varphi ,t)=E(r)e^{i(kr\omega t)}e^{\theta ^2/\theta _0^2},$$
(1)
where $`k=\omega /c`$ and $`c`$ is the speed of light. Deduce the field near the focus.
Since the Kirchhoff diffraction formalism requires the volume to be charge free, the time-reversed technique is not applicable to cases where the source of the field is inside the volume. Nonetheless, the reader may find it instructive to attempt to apply the time-reversed diffraction integral to the example of an oscillating dipole at the origin.
## 2 The Kirchhoff Integral via Green’s Theorem
A standard formulation of Kirchhoff’s diffraction integral for a scalar field $`\psi (𝐱)`$ with time dependence $`e^{i\omega t}`$ is
$$\psi (𝐱)\frac{k}{2\pi i}_S\frac{e^{ikr^{}}}{r^{}}\psi (𝐱^{})𝑑\mathrm{Area}^{},$$
(2)
where the spherical waves $`e^{i(kr^{}\omega t)}/r^{}`$ are outgoing, and $`r^{}`$ is the magnitude of vector $`𝐫^{}=𝐱𝐱^{}`$.
For a time-reversed formulation in which we retain the time dependence as $`e^{i\omega t}`$, the spherical waves of interest are the incoming waves $`e^{i(kr^{}+\omega t)}/r^{}`$. In brief, the desired time-reversed diffraction integral is obtained from eq. (2) on replacing $`i`$ by $`i`$:
$$\psi (𝐱)\frac{ik}{2\pi }_S\frac{e^{ikr^{}}}{r^{}}\psi (𝐱^{})𝑑\mathrm{Area}^{}.$$
(3)
For completeness, we review the derivation of eqs. (2)-(3) via Green’s theorem. See also, sec. 10.5 of ref. .
Green tells us that for any two well-behaved scalar fields $`\varphi `$ and $`\psi `$,
$$_V(\varphi ^2\psi \psi ^2\varphi )𝑑\mathrm{Vol}=_S(\varphi ^{}\psi \psi ^{}\varphi )𝑑𝐒^{}.$$
(4)
The surface element $`d𝐒^{}`$ is directly outward from surface $`S`$. We consider fields with harmonic time dependence at frequency $`\omega `$, and assume the factor $`e^{i\omega t}`$. The wave function of interest, $`\psi `$, is assumed to have no sources within volume $`V`$, and so obeys the Helmholtz wave equation,
$$^2\psi +k^2\psi =0.$$
(5)
We choose function $`\varphi (𝐱)`$ to correspond to waves associated with a point source at $`𝐱^{}`$. That is,
$$^2\varphi +k^2\varphi =\delta ^3(𝐱𝐱^{}).$$
(6)
The well-known solutions to this are the incoming and outgoing spherical waves,
$$\varphi _\pm (𝐱,𝐱^{})=\frac{e^{\pm ikr^{}}}{r^{}},$$
(7)
where the \+ sign corresponds to the outgoing wave. We recall that
$$^{}r^{}=\frac{𝐫^{}}{r^{}}=\widehat{𝐧}_o,$$
(8)
where $`\widehat{𝐧}_o`$ points towards the observer at x. Then,
$$^{}\varphi _\pm =ik\widehat{𝐧}_o\left(1\pm \frac{1}{ikr^{}}\right)\varphi .$$
(9)
Inserting eqs. (5)-(9) into eq. (4), we find
$$\psi (𝐱)=\frac{1}{4\pi }_S\frac{e^{\pm ikr^{}}}{r^{}}\widehat{𝐧}^{}\left[^{}\psi \pm ik\widehat{𝐧}_o\left(1\pm \frac{1}{ikr^{}}\right)\psi \right]𝑑\mathrm{Area}^{},$$
(10)
where the overall minus sign holds with the convention that $`\widehat{𝐧}^{}`$ is the inward normal to the surface.
We only consider cases where the source of the wave $`\psi `$ is far from the boundary surface, so that on the boundary $`\psi `$ is well approximated as a spherical wave,
$$\psi (𝐱^{})A\frac{e^{ikr_s}}{r_s},$$
(11)
where $`r_s`$ is the magnitude of the vector $`𝐫_s=𝐱^{}𝐱_s`$ from the effective source point $`𝐱_s`$ to the point $`𝐱^{}`$ on the boundary surface. In this case,
$$^{}\psi =ik\widehat{𝐧}_s\left(1\pm \frac{1}{ikr_s}\right)\psi ,$$
(12)
where $`\widehat{𝐧}_s=𝐫_𝐬/r_s`$
We also suppose that the observation point is far from the boundary surface, so that $`kr^{}1`$ as well as $`kr_s1`$. Hence, we neglect the terms in $`1/ikr^{}`$ and $`1/ikr_s`$ to find
$$\psi (𝐱)=\frac{ik}{4\pi }_S\frac{e^{\pm ikr^{}}}{r^{}}\widehat{𝐧}^{}(\widehat{𝐧}_s\pm \widehat{𝐧}_o)\psi (𝐱^{})𝑑\mathrm{Area}^{}.$$
(13)
The usual formulation, eq. (2), of Kirchhoff’s law is obtained using outgoing waves (+ sign), and the paraxial approximation that $`\widehat{𝐧}^{}\widehat{𝐧}_o\widehat{𝐧}_s`$. The latter tacitly assumes that the effective source is outside volume $`V`$.
Here, we are interested in the case where the effective source is inside the volume $`V`$, so that the paraxial approximation is $`\widehat{𝐧}^{}\widehat{𝐧}_o\widehat{𝐧}_s`$. When we use the incoming wave function to reconstruct $`\psi (𝐱,t)`$ from information on the boundary at time $`t^{}>t`$, we use the $``$ sign in eq. (13) to find eq. (3).
Note that in this derivation, we assumed that $`\psi `$ obeyed eq. (5) throughout volume $`V`$, and so the actual source of $`\psi `$ cannot be within $`V`$. Our time-reversed Kirchhoff integral (3) can only be applied when any source inside $`V`$ is virtual. This includes the interesting case of a focus of an optical system (secs. 4 and 5). However, we cannot expect eq. (3) to apply to the case of a physical source, such as an oscillating dipole, inside volume $`V`$ (sec. 6). The laws of diffraction do not permit electromagnetic waves to converge into a volume smaller than a wavelength cubed, and so eq. (3) cannot be expected to describe the near fields around a source smaller than this.
## 3 A Plane Wave
The time-reversed Kirchhoff integral (3) for the $`x`$ component of the electric field is
$$E_x(\mathrm{obs},\mathrm{now})=\frac{ik}{2\pi }\frac{e^{ikr^{}}}{r^{}}E_x(r,\theta ,\varphi ,\mathrm{future})\mathrm{dArea},$$
(14)
where $`r^{}`$ is the distance from the observation point $`𝐫_{\mathrm{obs}}=(x,y,z)`$ in rectangular coordinates to a point $`𝐫=r(\mathrm{sin}\theta \mathrm{cos}\varphi ,\mathrm{sin}\theta \mathrm{sin}\varphi ,\mathrm{cos}\theta )`$ on a sphere of radius $`r`$ in the far field.
As a first example, consider a plane electromagnetic wave,
$$E_x=E_0e^{i(kz\omega t)}=E_0e^{i(kr\mathrm{cos}\theta \omega t)},$$
(15)
where the second form holds in a spherical coordinate system $`(r,\theta ,\varphi )`$ where $`\theta `$ is measured with respect to the $`z`$ axis. We take the point of observation to be $`(x,y,z)=(0,0,r_0)`$, and evaluate the diffraction integral (14) over a sphere of radius $`rr_0`$. In the exponential factor in the Kirchhoff integral, we approximate $`r^{}`$ as
$$r^{}r\widehat{𝐫}𝐫_{\mathrm{obs}}=rr_0\mathrm{cos}\theta ,$$
(16)
while in the denominator we approximate $`r^{}`$ as $`r`$. Then,
$`E_x(\mathrm{obs})`$ $``$ $`{\displaystyle \frac{ik}{2\pi }}{\displaystyle _1^1}r^2d\mathrm{cos}\theta {\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle \frac{e^{ik(rr_0\mathrm{cos}\theta )}}{r}}E_0e^{ikr\mathrm{cos}\theta }`$
$`=`$ $`{\displaystyle \frac{r}{r+r_0}}E_0[e^{ikr_0}e^{ik(2r+r_0)}]`$
$``$ $`E_0e^{ikr_0},`$
where we ignore the rapidly oscillating term $`e^{ik(2r+r_0)}`$ as unphysical.
This verifies that the time-reversed diffraction formula works for a simple example.
## 4 The Transverse Field near a Laser Focus
We now consider the far field of a laser beam whose optic axis is the $`z`$ axis with focal point at the origin. The polarization is along the $`x`$ axis, and the electric field has Gaussian dependence on polar angle with characteristic angle $`\theta _01`$. Then, we can write
$$E_x(r,\theta ,\varphi )=E(r)e^{ikr}e^{\theta ^2/\theta _0^2},$$
(18)
where $`E(r)`$ is the magnitude of the electric field on the optic axis at distance $`r`$ from the focus. In the exponential factor in the Kirchhoff integral (14), $`r^{}`$ is the distance from the observation point $`𝐫_{|rmobs}=(x,y,z)`$ to a point $`𝐫=r(\mathrm{sin}\theta \mathrm{cos}\varphi ,\mathrm{sin}\theta \mathrm{sin}\varphi ,\mathrm{cos}\theta )`$ on the sphere. We approximate $`r^{}`$ as
$$r^{}r\widehat{𝐫}𝐫_{\mathrm{obs}}=rx\mathrm{sin}\theta \mathrm{cos}\varphi y\mathrm{sin}\theta \mathrm{sin}\varphi z\mathrm{cos}\theta ,$$
(19)
while in the denominator we approximate $`r^{}`$ as $`r`$. Inserting eqs. (18) and (19) into (14), we find
$`E_x(\mathrm{obs})`$ $`=`$ $`{\displaystyle \frac{ikrE(r)}{2\pi }}{\displaystyle _1^1}e^{ikz\mathrm{cos}\theta }e^{\theta ^2/\theta _0^2}d\mathrm{cos}\theta {\displaystyle _0^{2\pi }}e^{ikx\mathrm{sin}\theta \mathrm{cos}\varphi +iky\mathrm{sin}\theta \mathrm{sin}\varphi }𝑑\varphi `$ (20)
$`=`$ $`ikrE(r){\displaystyle _1^1}e^{ikz\mathrm{cos}\theta }e^{\theta ^2/\theta _0^2}J_0(k\rho \mathrm{sin}\theta )d\mathrm{cos}\theta ,`$
where
$$\rho =\sqrt{x^2+y^2},$$
(21)
and $`J_0`$ is the Bessel function of order zero.
Since we assume that the characteristic angle $`\theta _0`$ of the laser beam is small, we can approximate $`\mathrm{cos}\theta `$ as $`1\theta ^2/2`$ and $`k\rho \mathrm{sin}\theta `$ as $`k\rho \theta `$. Then, we have
$`E_x(\mathrm{obs})`$ $``$ $`ikrE(r)e^{ikz}{\displaystyle _0^{\mathrm{}}}e^{(2/\theta _0^2+ikz)\theta ^2/2}J_0\left(\sqrt{2}k\rho \sqrt{\theta ^2/2}\right)d(\theta ^2/2)`$ (22)
$`=`$ $`{\displaystyle \frac{ik\theta _0^2rE(r)e^{ikz}e^{k^2\theta _0^2\rho ^2/4(1+ik\theta _0^2z/2)}}{2(1+ik\theta _0^2z/2)}},`$
where the Laplace transform, which is given explicitly in , can be evaluated using the series expansion for the Bessel function. This expression can be put in a more familiar form by introducing the Rayleigh range (depth of focus),
$$z_0=\frac{2}{k\theta _0^2},$$
(23)
and the so-called waist of the laser beam,
$$w_0=\theta _0z_0=\frac{2}{k\theta _0}.$$
(24)
We define the electric field strength at the focus $`(\rho =0,z=0)`$ to be $`E_0`$, so we learn that the far-field strength is related by
$$E(r)=i\frac{z_0}{r}E_0.$$
(25)
The factor $`i=e^{i\pi /2}`$ is the $`90^{}`$ Guoy phase shift between the focus and the far field. Then, the transverse component of the electric field near the focus is
$`E_x(x,y,z)`$ $``$ $`E_0{\displaystyle \frac{e^{\rho ^2/w_0^2(1+iz/z_0)}e^{ikz}}{(1+iz/z_0)}}`$ (26)
$`=`$ $`E_0{\displaystyle \frac{e^{\rho ^2/w_0^2(1+z^2/z_0^2)}e^{i\mathrm{tan}^1z/z_0}e^{i\rho ^2z/w_0^2z_0(1+z^2/z_0^2)}e^{ikz}}{\sqrt{1+(z/z_0)^2}}}.`$
This is the usual form for the lowest-order mode of a linearly polarized Gaussian laser beam . Figure 1 plots this field.
The Gaussian beam (26) could also be deduced by a similar argument using eq. (2), starting from the far field of the laser before the focus. The form (26) is symmetric in $`z`$ except for a phase factor, and so is a solution to the problem of transporting a wave from $`z=r`$ to $`z=+r`$ such that the functional dependence on $`\rho `$ and $`z`$ is invariant up to a phase factor. One of the earliest derivations of the Gaussian beam was based on the formulation of this problem as an integral equation for the eigenfunction (26).
## 5 The Longitudinal Field
Far from the focus, the electric field E(r) is perpendicular to the radius vector r. For a field linearly polarized in the $`x`$ direction, there must also be a longitudinal component $`E_z`$ related by
$$𝐄\widehat{𝐫}=E_x\mathrm{sin}\theta \mathrm{cos}\varphi +E_z\mathrm{cos}\theta =0.$$
(27)
Thus, far from the focus,
$$E_z(𝐫)=E_x(𝐫)\mathrm{tan}\theta \mathrm{cos}\varphi .$$
(28)
Then, similarly to eqs. (14) and (20), we have
$`E_z(\mathrm{obs})`$ $`=`$ $`{\displaystyle \frac{ik}{2\pi }}{\displaystyle \frac{e^{ikr^{}}}{r^{}}E_z(𝐫)𝑑\mathrm{Area}}`$ (29)
$`=`$ $`{\displaystyle \frac{ikrE(r)}{2\pi }}{\displaystyle _1^1}e^{ikz\mathrm{cos}\theta }e^{\theta ^2/\theta _0^2}\mathrm{tan}\theta d\mathrm{cos}\theta {\displaystyle _0^{2\pi }}e^{ikx\mathrm{sin}\theta \mathrm{cos}\varphi +iky\mathrm{sin}\theta \mathrm{sin}\varphi }\mathrm{cos}\varphi d\varphi `$
$`=`$ $`{\displaystyle \frac{ikxz_0E_0}{\rho }}{\displaystyle _1^1}e^{ikz\mathrm{cos}\theta }e^{\theta ^2/\theta _0^2}\mathrm{tan}\theta J_1(k\rho \mathrm{sin}\theta )d\mathrm{cos}\theta ,`$
using eq. (3.937.2) of .
We again note that the integrand is significant only for small $`\theta `$, so we can approximate eq. (29) as the Laplace transform
$`E_z(x,y,z)`$ $``$ $`ik^2xz_0E_0e^{ikz}\sqrt{2}{\displaystyle _0^{\mathrm{}}}e^{(2/\theta _0^2+ikz)\theta ^2/2}\sqrt{\theta ^2/2}J_1\left(\sqrt{2}k\rho \sqrt{\theta ^2/2}\right)d(\theta ^2/2)`$ (30)
$`=`$ $`{\displaystyle \frac{ik^2\theta _0^4xz_0E_0e^{ikz}e^{\rho ^2/w_0^2(1+iz/z_0)}}{4(1+iz/z_0)^2}}`$
$`=`$ $`i\theta _0{\displaystyle \frac{x}{w_0}}{\displaystyle \frac{E_x(x,y,z)}{(1+iz/z_0)}},`$
with $`E_x`$ given by eq. (26). Figure 2 plots this field.
Together, the electric field components given by eqs. (26) and (30) satisfy the Maxwell equation $`𝐄=0`$ to order $`\theta _0^2`$ .
## 6 Oscillating Dipole at the Origin
We cannot expect the Kirchhoff diffraction integral to apply to the example of an oscillating dipole, if our bounding surface surrounds the dipole. Let us see what happens if we try to use eq. (3) anyway.
The dipole is taken to be at the origin, with moment $`p`$ along the $`x`$ axis. Then, the $`x`$ component of the radiation field is
$$E_x=k^2p\mathrm{sin}\theta _x\frac{e^{ikr}}{r}.$$
(31)
where $`\theta _x`$ is the angle between the $`x`$ axis and a radius vector to the observer. We consider an observer near the origin at $`(x,y,z)=(0,0,r_0)`$, for which $`\mathrm{sin}\theta _x=1`$, and so
$$E_x(\mathrm{obs})=k^2p\frac{e^{ikr_0}}{r_0}.$$
(32)
We now attempt to reconstruct this field near the origin from its value on a sphere of radius $`r`$ using the time-reversed Kirchhoff integral (3). We use a spherical coordinate system $`(r,\theta ,\varphi )`$ that favors the $`z`$ axis. Then, the $`x`$ component of the radiation field on the sphere of radius $`r`$ is
$$E_x(r,\theta ,\varphi )=k^2p\sqrt{1\mathrm{sin}^2\theta \mathrm{cos}^2\varphi }\frac{e^{ikr}}{r}.$$
(33)
This form cannot be integrated analytically, so we use a Taylor expansion of the square root, which will lead to an expansion in powers of $`1/r_0`$. It turns out that the coefficient of the $`1/r_0`$ term, which is our main interest, is very close to that if we simply approximate the square root by unity. For brevity, we write
$$E_x(r,\theta ,\varphi )k^2p\frac{e^{ikr}}{r}.$$
(34)
In the time-reversed Kirchhoff integral (3), we make the usual approximation that $`r^{}=rr_0\mathrm{cos}\theta `$ in the exponential factor, but $`r^{}=r`$ in the denominator. Then, using eq. (34) we have
$`E_x(\mathrm{obs})`$ $``$ $`{\displaystyle \frac{ik^3pe^{ikr}}{2\pi r}}{\displaystyle _1^1}r^2d\mathrm{cos}\theta {\displaystyle _0^{2\pi }}𝑑\varphi e^{ikr_0\mathrm{cos}\theta }{\displaystyle \frac{e^{ikr}}{r}}`$ (35)
$`=`$ $`k^2p{\displaystyle \frac{e^{ikr_0}}{r_0}}k^2p{\displaystyle \frac{e^{ikr_0}}{r_0}}`$
$`=`$ $`2ik^3p{\displaystyle \frac{\mathrm{sin}kr_0}{kr_0}}.`$
The first, outgoing wave in middle line of eq. (35) is the desired form, but the second, incoming wave is of the same magnitude. Together, they lead to the form $`\mathrm{sin}(kr_0)/kr_0`$ which is nearly constant for kr0
<
1
<
𝑘subscript𝑟01kr_{0}\mathrel{\vbox{\kern 0.0pt\hbox{$<$} \kern 0.0pt\hbox{$\sim$} }}1. The presence of outgoing as well as incoming waves is to be expected because dipole radiation is azimuthally symmetric about the $`x`$ axis. In the absence of a charged source at the origin, an outgoing wave at $`\theta =\pi `$ must correspond to an incoming wave at $`\theta =0`$.
The result that the reconstructed field is uniform for distances within a wavelength of the origin is consistent with the laws of diffraction that electromagnetic waves cannot be focused to a region smaller than a wavelength. Far fields of the form (31) could only be propagated back to the form of dipole fields near the origin with the addition of nonradiation fields tied to a charge at the origin. Such a construction is outside the scope of optics and diffraction.
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# Another Faint UV Object Associated with a Globular Cluster X-Ray Source: The Case of M921footnote 11footnote 1Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by AURA, Inc., under NASA contract NAS5-26555
## 1 Introduction
Despite their rarity, it is the exotic creatures that attract the crowds at the zoo; similarly the exotic objects in the stellar zoo attract our attention. Unusual environments often lead to relatively large populations of the exotic. So it is with the cores of the Galactic Globular Cluster (GGCs), which have long been thought to harbor a variety of exotic objects—blue stragglers, low mass X-ray binaries, cataclysmic variables, millisecond pulsars, etc. Most of these objects are thought to result from various kinds of binary systems whose nature and even existence can be strongly affected by dynamics in the dense cluster cores.
When a binary system contains a compact object (like a neutron star or white dwarf) and a close enough secondary, mass transfer can take place. The streaming gas, its impact on the compact object, or the presense of an accretion disk can give such systems observational signatures which make them stand out above ordinary cluster stars. These signatures might include X-ray emission, significant radiation in the ultraviolet (UV), emission lines, or rapid time variations. The first evidence for such objects in globular clusters was the discovery of X-ray sources. One population of X-ray sources with $`L_X>10^{34.5}\mathrm{erg}\mathrm{sec}^1`$ (the so-called Low Mass X-ray Binaries, LMXB) are thought to be binary systems with an accreting neutron star because of their X-ray bursts. LMXBs are very overabundant (a factor 100) in GGCs with respect to the field, presumably because the high stellar density has led to many capture binaries.
Given the existence of neutron star systems in GGCs one might expect to find many more analogous systems involving white dwarfs (WDs). In the field, binary systems in which a WD is accreting material from a late type dwarf, i.e., a main sequence or subgiant star, are observed cataclysmic variables (CVs). CVs are well-studied objects in the field, where they are thought to form by the evolution of primordial binaries. They come in many varieties depending on stellar masses, mass transfer rates, magnetic field strength, etc. In GGCs one can expect even more variety because CVs located in dense clusters could have been created by dynamical processes (Hut & Verbunt, 1983, Bailyn 1995), while the CVs in low-density clusters result from primordial binary systems (Verbunt & Meylan 1988).
Numerical simulations (e.g., DiStefano & Rappaport 1994) suggest that $`>100`$ white dwarf binaries might be found in massive clusters like 47 Tuc and $`\omega `$ Cen, and several 10’s in more typical clusters. Despite the expectation of large numbers of CV like stars, searches have turned up only a relatively small number. Of course, part of the problem arises because of the difficulty of search for rather faint objects in crowded globular cluster fields. However, exploiting the high resolution of HST, it has become possible to search GGC centers for several of the anticipated CV signatures. Still the number of candidates is small:
* More than 30 low luminosity X-ray sources with $`L_x<10^{34.5}\mathrm{erg}\mathrm{sec}^1`$ (hereafter LLGCXs) have been discovered in 19 GGCs (Johnston & Verbunt 1994). Despite their relatively large numbers there is no consensus model for LLGCXs (see Verbunt et al. 1994) and Hasinger, Johnston & Verbunt 1996). The fainter LLGCXs ($`L_x<10^{32}\mathrm{erg}\mathrm{sec}^1`$) might well be associated with CVs (van Paradijs (1983), Hertz & Grindlay, 1983).
* There are three objects connected with conventionally detected CVs: a dwarf nova in M5 (Margon, Downes & Gunn 1981); HST UV detections of optical counterparts to a dwarf nova in 47 Tuc (Paresce & DeMarchi 1994) and possibly the historical nova in M80 (Shara & Drissen 1995)
* Using HST, H$`\alpha `$ emission has been observed from three objects in NGC 6397 (Cool et al. 1995), and two objects in NGC 6752 (Bailyn et al. 1996).
* Also using HST, a number of candidate CVs have been selected on the basis of UV excess and variability: the Einstein dim source in 47 Tuc (Paresce, De Marchi & Ferraro 1992) and few CVs candidates in NGC 6624 (Sosin & Cool 1995).
We are involved in two long-term HST projects to study in detail the evolved populations in a sample of GGCs, at different wavelengths ranging from the UV to the near IR. Although we were not specifically hunting for CVs, we have used the UV exposures to search for exotic objects in the core of GGCs. This search has been very fruitful: in our data-base (9 clusters) we have found in all GGCs properly observed (with exposures deep enough and in the right UV bands) there is at least one faint object with a strong UV excess with respect to the main stellar population of the cluster). These stars are brighter in the UV than in the visible, so we will refer to them as UV-dominant (UVD). We wish to carefully distinguish between these objects and objects which are called ‘UV-excess’ objects on the basis of their colors in the visible or perhaps near UV. Previously we have reported on three UVD stars in the GGC M13 (Ferraro et al. 1997). Two of these objects have been found to lie within the error boxes of LLGCXs (Fox et al. 1996), and we argue that they are excellent CV candidates. Here we report on the discovery of another faint UVD star in the core of M92, and we suggest that this star is physically connected to the X-ray emission detected in the cluster.
## 2 Observations
HST-WFPC2 frames were obtained on December 1995 (Cycle 5 : GO 5969, PI: F. Fusi Pecci). We report here results obtained using the deep exposures (3600 sec and 2200 sec) through the U (F336W) and mid-UV (F255W) filters, respectively. The CMDs presented here are results of the four WFPC2 chips (namely PC1, WF2, WF3, WF4), obtained with the PC located on the cluster center (some further discussion of these results can be found in Ferraro et al. 1997, 1998).
All the reductions have been carried out using ROMAFOT (Buonanno et al 1983), a package specifically developed to perform accurate photometry in crowded fields. In order to identify the objects in each field we used the median frame, obtained by combining all single exposures in each color. The PSF-fitting procedure was then performed on each individual frame separately and the instrumental magnitudes were then averaged. The instrumental magnitudes have been converted to fixed aperture photometry and where appropriate calibrated to the Johnson system using equation 8 and Table 7 in Holtzmann et al. (1995). F255W magnitudes have been calibrated to the STMAG system using table 9 by Holtzmann et al. (1995).
In this paper we adopt for M92 a distance modulus $`(mM)_0=14.78`$ from Ferraro et al. (1999a) who have recently determined moduli for a sample of 61 GGCs within the framework of an homogeneous re-analysis of the evolved sequences of the CMD.
## 3 Results
Figure 1 shows the $`(m_{255},m_{255}U)`$-CMD for the global sample of stars detected in all the four WFPC2 chips. More than 20,000 stars have been measured in these filters in the global WFPC2 Field of View. Inspection of this diagram shows that a few (5) blue low luminosity objects lie significantly outside the main loci defined by the majority of the cluster stars. Four of them are clumped at $`(m_{255}U)1`$, and one object (namely star $`\mathrm{\#}8203`$) shows a very strong UV color $`(m_{255}U)<1.5`$. This object is at only $`15\stackrel{}{\mathrm{.}}7`$ from the cluster center (assumed at $`\alpha _{2000}=17^\mathrm{h}\mathrm{\hspace{0.17em}17}^\mathrm{m}\mathrm{\hspace{0.17em}07}\stackrel{\mathrm{s}}{\mathrm{.}}3,\delta _{2000}=43\mathrm{°}\mathrm{\hspace{0.17em}08}\mathrm{}\mathrm{\hspace{0.17em}11}\stackrel{}{\mathrm{.}}0`$, Djorgovski & Meylan 1993), just outside the core radius of the cluster ($`r_c=14^{\prime \prime }`$). It is located at the extreme Northern edge of the PC chip. This region is still very crowded, and the UV star is close (less than $`1^{\prime \prime }`$) to an HB star and a very bright giant ($`V12.8`$) which is heavily saturated in the $`V`$ and $`I`$ deep exposures. Hence, only measures in the $`F255W`$ and $`U`$ filters are possible.
The possible variability of the UV source was examined by analyzing each available frame separately. No clear indication of variability was revealed from this analysis, but we cannot strongly exclude this possibility since our observations do not have much time coverage.
Fox et al. (1996) recently presented ROSAT High Resolution Imager (HRI) observations of M92 and identified 7 low luminosity X-ray sources in the field of view of the cluster. In particular, they drew attention to the X-ray source found in the core, M92C (hereafter M92X-C). M92X-C is only $`17^{\prime \prime }`$ from the cluster center and has a high probability ($`99.8\%`$) of being associated with the cluster. The LLGCX is also located at the Northern edge of the PC in our HST field of view just in the region where the UV star has been detected. Figure 2 shows a region of $`30{}_{}{}^{\prime \prime }\times 30^{\prime \prime }`$ centered on the nominal position of the M92X-C. The contours of the X-ray emission (from Figure 2 by Fox et al. ) have been overplotted on a digital map of the $`F255W`$ image. The absolute positions, the observed magnitude and the X-ray flux for the UV star and the X-ray source are listed in Table 1. Note that the X-ray luminosity has been properly scaled in order to take into account the different distance modulus adopted here with respect to that used by Fox et al. (1996).
Because the UV object is $`4.5^{\prime \prime }`$ from the nominal position of the X-ray emission, we strongly suggest a physical connection between the two. Note that the 4 UV objects located in the CMD at $`(m_{255}U)1`$ are much more distant ($`d>37^{\prime \prime }`$) from the X-ray source.
## 4 Discussion
Faint UV stars (similar to UV8203) have been discovered in the core of other GGCs, and some of them have been found to be nearly coincident with X-ray sources. The discovery reported in this paper makes the possibility of a chance coincidence of UV objects with the X-ray sources appear even less likely. With the strengthening evidence for a (physical) connection between UV objects and LLGCXs, it is appropriate to compare the photometric properties of a sample of faint UV stars found in the vicinity of LLGCX in GGCs. In doing this we select some of the most recent findings:
M13— Ferraro et al. (1997) using deep UV-HST observations found three extremely blue, low luminosity objects in the very central region of M13. Two of them are nearly coincident with the two-peaked X-ray emission detected by Fox et al (1996).
M5— V101 in M5 was discovered by Oosterhoff (1941), who first suggested that it could be a CV (dwarf nova). Spectroscopic (Margon, Downes & Gunn 1981 and Naylor et al. 1989) and photometric observations (Shara, Potter & Moffat 1987) confirmed this suggestion. Hakala et al. (1997) detected X-ray emission associated with this object using using the ROSAT-HRI.
M80— Shara & Drissen (1995) have found two UV objects in the globular cluster M80, one of them might be associated with T Sco, a nova observed in 1860. This object might be connected with a LLGCX located at $`8^{\prime \prime }`$ from the UV object, although, as suggested by Hakala et al. (1997), the position from the ROSAT PSPC is not accurate enough for a definitive identification.
NGC 6397— Cool et al. (1998, C98) found a population of 7 UV stars in the core of NGC6397. They divided the UV star sample in two subgroups: 4 stars showing variability and UV excess (in the sense that they appear to be blue in the $`U`$ band but are indistinguishable from MS stars in the $`V,VI`$ plane), and three non-variable UV excess stars which are significantly hotter than the main-sequence in all observed CMD planes. They call this second class ‘nonflickerers’ (NF). The four variable stars are all within the ROSAT HRI X-ray error circle and have been confirmed to be CVs (Cool et al. 1995, Edmonds et al. 1999). However, they noted that “…two out three NFs are outside the error circles of the three central X-ray sources detected with ROSAT by Cool et al. (1993).” While C98 claimed that the NF stars were a new class of faint UV stars, it is worth noting that they are very similar to the three UV stars found one year before in M13 by Ferraro et al. (1997).
47 Tuc— At least two objects with a strong UV excess have been identified in the error box of X-ray source in the center of 47 Tuc: V1 (Paresce, DeMarchi & Ferraro 1992) and V2 (a blue variable discovered by Paresce & DeMarchi 1994 (see also Shara et al. 1996). V1 lies within the error circle of the X-ray source $`X0021.87221`$ detected by Einstein (Bailyn et al. 1988) and in the vicinity of a low luminosity X-ray source (X9 in Table 2 of Verbunt & Hasinger 1998). Verbunt & Hasinger (1998) also identified V2 as a candidate optical counterpart of their source X19.
NGC 6752— Bailyn et al. (1996) report the identification of two candidate CVs in NGC 6752. Both stars fall at the edge of the error circle of the X-ray source identified as B by Grindlay & Cool (1996) (flux given in Grindlay 1993). These stars are plotted as empty circles in Figure 3. Contrary to the behavior shown by all the other objects candidates they have a quite strong rise up toward red wavelengths suggesting that these objects have quite different spectral characteristics with respect to the other objects listed above. Their position in the CMD (see Figure 3 by Bailyn et al) resemble the CVs found by C98 in NGC 6397, however observations at wavelength shorter than $`B`$ are needed in order to better constrain the spectral behavior of these stars. For these reasons we exclude them from the following discussion.
The absolute magnitude in different photometric bands for the 10 UV stars possibly connected with X-ray emission in the 8 GGCs quoted above are listed in Table 2. Also reported are the adopted distance moduli and reddening from Ferraro et al. (1999b). Note that the distance moduli adopted typically differ from earlier papers, for example, for NGC 6397 is $`\mathrm{\Delta }(mM)_V0.2`$ mag larger than that adopted by Cool et al. (1998). All the absolute magnitudes and X-ray luminosities in Table 1 have been corrected accordingly. Note also that the X-ray fluxes were determined over different energy ranges which are given in the last column.
In Figure 3 we plot the absolute magnitude for each star listed in Table 2 as a function of the filter’s effective wavelength. In the figure different symbols refer to different clusters: open triangles for M13; filled circles for M92, asterisks for M80, empty pentagons for 47 Tuc, large $`X`$ for M5 and dashed lines for the three CVs in NGC 6397. For comparison the energy distribution for the field CV U Gem is plotted as an heavy solid line. Figure 3 provides an easy way to make quantitative comparison. In particular, all of the UV selected stars show the same overall spectral trend, and we conclude that some of the photometric properties of faint UVD stars associated with LLGCXs are in reasonable agreement with each other and U Gem. It is interesting to note that the slope of UV8203 in M92 and V101 in M5 appear to be steeper than the others, however far-UV observations are required in order to confirm this impression.
Though small, the sample listed in Table 2 can still be used to derive some average properties of these objects, for example the absolute $`V,B,U`$ magnitude, and $`UV`$, $`UB`$ colors which are used to characterize CVs in the field. From the data in Table 2 these figures turn to be: $`M_V=6.7\pm 0.8`$, $`M_B=7.0\pm 0.7`$, $`M_U=5.9\pm 1.0`$, $`UB=1.1\pm 0.3`$ and $`UV=1.0\pm 0.5`$, respectively. These figures are in good agreement with the typical absolute magnitude and colors for CVs in the field ($`M_V+7`$, for the dwarf novae CVs; see van Paradijs 1983).
We may push our working hypothesis that the UVD objects are physically associated with the X-ray sources further. The X-ray luminosity of each source is listed in the penultimate column of Table 2. As can be seen they have comparable X-ray luminosity, within the range 1–$`8\times 10^{32}\mathrm{erg}\mathrm{sec}^1`$ (with a mean value of $`L_x=4\pm 3\times 10^{32}\mathrm{erg}\mathrm{sec}^1`$).
We can compare the observed photometric characteristics of the UV objects found in GGCs with those obtained for field CVs using the values listed in Table 2 along with data from Table 1 of the recent compilation by Verbunt et al. 1997 (hereafter V97) who presented a catalog with 91 CVs in the field detected during the ROSAT All Sky Survey. The upper panel of Figure 4 shows the absolute $`V`$ magnitude as a function of the X-ray luminosity for all field CVs with known distances listed by V97. These are plotted as small empty triangles. The 6 UVD stars found in GGCs for which the $`V`$ magnitude has been measured are plotted as large filled circles. The $`V`$ magnitude for the 4 stars for which no $`V`$ magnitude was directly measured has been computed from $`M_U`$ assuming a mean color $`UV=1`$ (see above). These are plotted as large asterisks. The lower panel shows the X-ray luminosity distribution for the GGC objects (shaded histogram) compared with the distribution for the field CVs. This figure clearly shows that while the absolute $`V`$ magnitude for the candidates in clusters are fully consistent with the field CVs, the X-ray emission for CVs in GGCs seems systematically higher (as already suggested by Ferraro et al. 1997, see also Figure 4 in Verbunt & Hasinger 1998) indicating that the X-ray luminosity of objects found in GGCs is high relative to the visible compared to similar objects in the field.
There is a strong observational selection effect at work in Figure 4. The depth reached by the X-ray observations is typically only a factor of two or three below the level of the detected sources. The deepest is for NGC 6397 (Cool et al. 1993) which reaches $`L_x=3\times 10^{31}\mathrm{erg}\mathrm{sec}^1`$. Dashed lines in Figure 4 show the maximum depth reached in surveys for UVD stars/LLGCXs in M13, M92, & NGC 6397. Figure 8 of V97 shows that the vast majority of field CVs of all types are less luminous than this. A comparison of the $`M_V`$ for the optically identified LLGCXs with typical values from Warner (1987) shows that current optical surveys would also have missed many field CVs if they were located in GGCs. The relatively high $`L_x`$ and $`L_{\mathrm{opt}}`$ for LLGCXs could arise solely because that is all that surveys to date could detect.
This bias has important consequences. Hakala et al. (1997) used the high $`L_x/L_{\mathrm{opt}}`$ value of T Sco in M80 to argue that it was not associated with the candidate object suggested by Shara & Drissen. However, we now see that $`L_x/L_{\mathrm{opt}}`$ is systematically higher for GGC CVs as compared to their field sisters, and thus, the Hakala et al. argument is not valid.
Only the three field CVs with highest $`L_x`$ are consistent with the GGC CV candidates in Figure 4. These are the DQ Her systems V1223 Sgr, AO Psc and TV Col. These are strongly magnetic CVs of a class referred to as intermediate polars (IP). Grindlay (1999) suggests that IPs might dominate the ROSAT survey since they are expected to have a ratio $`F_x/F_{\mathrm{opt}}`$ greater than non-magnetic CVs. The problem with connecting the UVD stars with magnetic CVs is that the magnetic field might truncate the inner portion of the accretion disk (Grindlay 1999). Thus, MCVs might not be expected to have strong UV-excess.
It is worth noting that the V97 $`L_x`$ distributions for field dwarf novae (SU UMa, Z Cam, U Gem stars) are based on small samples (11, 7, 7) and still span one to two orders of magnitudes. Larger samples could well reach into the range observed in GGCs. Indeed, since we could well be observing the bright end of a sample of several tens of assorted types of CVs in a GGC, the large observed values of $`L_x`$ probably do not provide a significant constraint on CV types.
Perhaps the most solid detection of CVs in a GGC are those in NGC 6397 (Edmonds et al. 1999; C98). There are several indications of mass exchange: H$`\alpha `$ emission, X-ray radiation, spectra with emission lines, and time variability or flickering of the optical radiation. The colors of these objects become redder as the wavelength of filters employed increases. This suggests in, for example, $`V`$ and $`I`$ the light from the cool secondary dominates a rather weak accretion disk. This is one of the factors which led Edmonds et al. (1999) and Grindlay (1999) to associate these objects with the intermediate polar class of magnetic CV. The CVs in NGC 6752 are probably similar objects.
The UVD objects are clearly not this sort of creature. Many have properties one might expect from a generic CV, i.e., white dwarf/main sequence binaries with some mass exchange. They share many properties with some field CVs—X-ray luminosity, UV colors, absolute visual magnitudes. The facts that they tend to lie at the extremes of the distributions and that some do not have detected X-radiation may simply be a consequence of the rather low sensitivity of current surveys. In this scenario the difference between the UV dominant CVs and those in NGC 6397 is simply a matter of the relative importance of light from the secondary and the accretion disk. For the GGC UVD stars, as with most CVs in the field, the accretion disk significantly outshines the secondary star (see for example Figure 2.16–2.17 in Warner 1995).
What other options are there? It might be tempting to extend either the horizontal branch (HB) or white dwarf (WD) sequences into the region of the UVD stars. However, the HB terminates at the helium burning main sequence. Particularly for M13 where the observed HB extends almost to its termination, we see that the UVD stars are significantly fainter (Figure 1 of Ferraro et al. 1997). The ordinary, i.e., carbon/oxygen, WD sequence is well observed, for instance, in M 4 and all the observed WDs are fainter than $`M_U8.5`$ (see Figure 6 by Richer et al. 1997), $`1`$ mag. fainter than the faintest UVD in Table 2. So the UVD stars are not HB stars or carbon/oxygen white dwarfs.
Edmonds et al. (1999) have shown that at least one UVD star is not a CV. Their HST-FOS spectrum of one of the NF-UV objects identified by C98 in NGC 6397 suggests a very hot high gravity object. They argue that this is a low mass ($`0.25M_{}`$) helium white dwarf. Such an object could arise from either mass exchange in a binary system shortly after the primary leaves the main sequence or via stellar collisions. The object has a velocity of $`250\mathrm{km}\mathrm{s}^1`$ relative to the cluster CVs which is most easily explained if it still is in a binary system with a dark companion. (It might be worth noting that NGC 6397 is a post-core-collapse cluster with a very dense core, whereas the UVD objects we have found are in the moderate density clusters M13 & M92.) Could He WDs account for all GGC UVD stars? This is certainly not the case for the objects in M5, M80, & 47 Tuc which have shown variability consistent with known CV types. He WDs would not produce X-rays in significant quantities. The chance association of three LLGCXs with such rare objects as the UVD objects we have found in M13 & M92 would seem unlikely.
## 5 Conclusions
We have now identified three UV dominant objects which appear to be associated with X-ray sources in the GGCs M13 and M92. We argue that these are generically CVs, i.e., white dwarf/main sequence binaries with some mass exchange. They share many properties with some field CVs—X-ray luminosity, UV colors, absolute visual magnitudes—although they tend to lie at the extremes of the distributions. The relatively high optical, UV, & X-ray luminosities are consistent with the notion that current surveys do not reach deep enough to detect most of the CVs in GGCs. On the other hand, cluster CVs are older and live in a dramatically different environment from their sisters in the field. Thus, it seems reasonable that cluster objects might be a new class of CVs with properties which slightly differ from those in the field.
The CVs found in NGC 6397 and NGC 6752 differ significantly from our objects. Some of these have little or no UV excess. This should not be surprising. The searches in NGC 6397 and NGC 6752 relied on H$`\alpha `$ and $`R`$ band, whereas we used the UV. There is considerable variety in the properties of field CVs and no reason to suspect less variety in GGC CVs. Different search techniques operating at the margin of detectability will certainly turn up different kinds of objects.
The UVD objects in GGCs probably come in several varieties. (Our GGC projects do not seem to come up with simple answers.) Edmonds et al. (1999) have shown that one is a hot high gravity object, arguablly a He-WD. Still we suspect that most of these will turn out to be CVs, and that many UVD objects with no X-radiation detected to date will show up as LLGCXs in more sensitive X-ray surveys. This suspicion is fueled by the belief that a significant population of generic CVs must be present in GGCs.
We have yet to make an observation directly showing the hot diffuse gas which would be the definitive evidence that our UVD objects are CVs. HST STIS spectra could give such evidence. It would be extremely valuable to develop a technique to identify GGC CVs using UV photometry. Cluster cores are very congested in the red and H$`\alpha `$/$`R`$ band searches will obviously be incomplete because of the interference by bright red giants (see Figure 2 of Bailyn et al. 1996). On the other hand the core of even the densest clusters are relatively open in the UV (see Figure 1 of Ferraro et al. 1999b)
This research was partially supported by the Agenzia Spaziale Italiana (ASI) and by the MURST as part of the project Dynamics and Stellar Evolution in Globular Clusters. F. R. F. acknowledges the ESO Visiting Program for the hospitality. R. T. R. is supported in part by NASA Long Term Space Astrophysics Grant NAG 5-6403 and STScI/NASA Grant GO-6607.
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# 1 The “setting-sun” diagram. The external lines are not amputated. The fifth coordinate of each point, as indicated on the first row below the diagram, corresponds to the case where 1/"/"p (eq. ()) occurs on the rightmost line, and the leftmost point is far off the boundary. This yields eq. (). The second row corresponds to eq. ().
March 2000 TAUP–2626–00
New Domain-Wall Fermion Actions
Yigal Shamir
School of Physics and Astronomy
Beverly and Raymond Sackler Faculty of Exact Sciences
Tel-Aviv University, Ramat Aviv, 69978 ISRAEL
shamir@post.tau.ac.il
ABSTRACT
> In perturbation theory, the wave function of domain-wall quarks decreases exponentially with the fifth coordinate. We show that, regardless of the quark’s own momentum, the fall-off rate of the one-loop wave function is equal to the slowest rate encountered at tree-level for any lattice four-momentum. We propose new domain-wall actions involving beyond-nearest neighbor couplings in the four physical dimensions, for which the perturbative wave function decreases much faster. It is hoped that the new actions may preserve the good chiral properties of domain-wall fermions up to larger values of the lattice spacing.
1. Introduction
The coupling of the right-handed and the left-handed components of Wilson fermions through the QCD interaction leads in the continuum limit to the chiral anomaly , but for finite lattice spacing $`a`$ it also leads to lattice-artefact violations of chiral symmetries. This results in an additive renormalization of the quark mass, as well as in a severe tuning problem for four-fermion operators which are needed for the computation of weak matrix elements. The mass renormalization is $`O(g_0^2)`$ in lattice units. Since $`g_0^2(\mathrm{log}(a\mathrm{\Lambda }_{\mathrm{QCD}}))^1`$, the mass renormalization diverges like $`(a\mathrm{log}(a\mathrm{\Lambda }_{\mathrm{QCD}}))^1`$ in the continuum limit $`a0`$. In the challenging lattice calculation of non-leptonic kaon decays (e.g. $`ϵ^{}/ϵ`$) the tuning problem is formidable because of the large number of chirality-disallowed mixings.
In the domain-wall formulation of lattice QCD \[2-6\], the two chiral components arise as surface states on opposite boundaries of a five-dimensional lattice, and one expects their coupling to vanish when the size of the fifth dimension tends to infinity. (The fifth coordinate will be denoted $`s`$, and it takes values $`0sN_s1`$.) A precise non-perturbative characterization of chiral symmetry violations can be given in terms of the transfer matrix for hopping in the $`s`$-direction . Being the result of slow decay of correlations in the $`s`$-direction, chiral symmetry violations are associated with near-unity eigenvalues of that transfer matrix. For a given gauge-field configuration the approach to the chiral limit is exponential iff the spectrum of the transfer matrix has a gap, and the fall-off rate is determined by the size of the gap.
In full QCD there are several analytic results concerning the $`N_s\mathrm{}`$ limit. The approach to the chiral limit is exponential in perturbation theory , and the same is true non-perturbatively if a constrained gauge action (believed to be in the same universality class as the standard plaquette action) is used .
For an unconstrained action one can also prove non-perturbatively that chiral symmetry is restored in the limit $`N_s\mathrm{}`$, provided the (finite!) number of sites in each of the four physical lattice dimensions is held fixed . When near-unity eigenvalues start playing a significant role, the chiral limit may be approached as slow as $`1/N_s`$. The proof that certain symmetries are restored in the limit $`N_s\mathrm{}`$ is actually valid for any value of the coupling constant. But the identification of the restored symmetries as chiral ones depends on the fermion spectrum. It was recently shown that within the strong-coupling expansion the massless spectrum of the domain-wall lattice hamiltonian is either doubled or empty. Therefore the restored symmetries are not chiral at strong coupling (For further details see Appendix C.1).
Of major importance is the question of how close to the chiral limit one gets in Monte-Carlo simulations. In trying to answer this question we rely on two sources. The first is the spectrum of the transfer matrix, or of the closely-related hermitian Wilson-Dirac operator. In the latter case, a key finding is that the spectral density of near-zero modes (corresponding to near-unity eigenvalues of the transfer matrix) rises by two orders of magnitude as the (quenched) coupling changes from $`6/g^2\beta =6.3`$ ($`a^14\mathrm{GeV}`$) to $`\beta =5.7`$ ($`a^11\mathrm{GeV}`$). We hope that more results on the eigenvalue spectrum will be available in the future.
More information is available through lattice computations of various correlation functions \[13-20\]. A detailed numerical study of the chiral limit of domain-wall fermions was first carried out in the Schwinger model . In QCD the first domain-wall simulations were promising, and the results for weak matrix elements ($`B_K`$, $`O_{LL}`$) and for the strange-quark mass were in agreement with other methods. As of today more data is available. The pion-mass squared, which should extrapolate linearly to zero with the quark mass, is the most obvious measure of chiral symmetry. Using domain-wall fermions at quenched $`\beta =6.0`$ ($`a^12\mathrm{GeV}`$) the extrapolated pion mass (for $`N_s20`$) does not vanish exactly at zero quark mass but, rather, at a negative value of the order of few times $`10^3`$ in lattice units (see also ref. ). This value, however small, is in the same range as the light quark masses.
In a sense, chiral symmetry violations are worst for the pion mass, because the lattice-artefact term in the PCAC relation is an ensemble average of a positive fermion correlator (see Appendix C.1; we comment that that lattice-artefact term is a better measure of chiral symmetry violation compared to the extrapolated pion mass, since it does not suffers from theoretical uncertainties due to chiral perturbation theory and due to finite-volume effects).
No such positivity is encountered in the calculation of weak matrix elements , so chiral symmetry violations are expected to be smaller in this case. For example, in a recent simulation of four-fermion operators using the non-perturbative renormalization scheme, again at quenched $`\beta =6.0`$, and using $`N_s=16`$, it was found that mixing into wrong-chirality operators was practically zero (in comparison with 10% for Wilson fermions in a typical example).
Going to a smaller value of the inverse lattice spacing, the situation at $`a^11\mathrm{GeV}`$ is unsatisfactory, as deviations from chiral symmetry are significant even for $`N_s`$ as large as 50 or 100 (in both quenched and dynamical simulations) . In the opposite direction, at $`a^1>\mathrm{\hspace{0.33em}3}\mathrm{GeV}`$, no difficulties with the restoration of chiral symmetry have been reported .
We believe that the existing results, especially those for weak matrix elements at $`a^12\mathrm{GeV}`$, do represent a breakthrough compared to the “pre domain-wall era”. On the other hand, the results for the pion mass at $`a^12\mathrm{GeV}`$ are not as good as one would hope for, and the present situation at $`a^11\mathrm{GeV}`$ makes scaling studies with domain-wall fermions very difficult.
Having summarized the situation in Monte-Carlo simulations let us return to the underlying physics. The key question is what mechanism(s) determine the abundance of near-unity eigenvalues of the transfer matrix. It is known that a few exact-unity eigenvalues must occur during the transition from one topological sector to another on a finite lattice with periodic boundary conditions . We believe, however, that the role of topology changing has been over-emphasized, for topological considerations alone do not explain the proliferation of near-unity eigenvalues nor the magnitude of the ensuing chiral symmetry violations.
A simple explanation may be that the observed chiral symmetry violations arise (mainly) from generic fluctuations of the gauge field . The effect of fluctuations need not be small, because the coupling constant used in simulations is not small either. At the same time, as long as the coupling constant has not grown too much, perturbation theory should provide a reliable approximation of the leading quantum effects.
In this paper we calculate the fifth coordinate’s wave function of domain-wall quarks in the one-loop approximation (Sec. 2). The results lead us to consider new classes of domain-wall actions (Sec. 3). A preliminary account of this work was given in ref. . (An alternative/complementary approach, whose relative merits are discussed in Sec. 4, is to employ an improved gauge action .)
We now give an overview of the one-loop calculation. For large $`N_s`$, zero bare quark mass and with the right-handed quark field near the $`s=0`$ boundary, the dressed fermion propagator near that boundary is
$$G_{s,s^{}}(p)=P_+\chi (s)\frac{1}{i/p(1+\mathrm{\Sigma }_K)}\chi (s^{})+\mathrm{Reg}.,s,s^{}N_s,$$
(1.1)
where $`P_\pm =\frac{1}{2}(1\pm \gamma _5)`$, “Reg” stands for a continuous function of the four-momentum $`p`$, and $`\mathrm{\Sigma }_K\mathrm{\Sigma }_K(g^2,g^2\mathrm{log}(p^2))`$. A unique feature of the domain-wall scheme is $`\chi (s)`$, the $`s`$-coordinate wave function for (right-handed) quark modes. At tree level one has
$$\chi _0(s)q_0^s(1M)^s.$$
(1.2)
(The five-dimensional mass term $`M`$ is often referred to as the domain-wall height, and should not be confused with the quark mass .) By choosing $`M=1`$ the free wave function can be completely localized on the boundary
$$\chi _0(s)=\underset{M1}{lim}(1M)^s=\delta _{s,0}.$$
(1.3)
The result of the one-loop calculation of the wave function is
$$\chi _1(s)s^2q_1^s,$$
(1.4)
which contains also a power correction. Like an ordinary Wilson mass, $`M`$ is renormalized additively. Making an optimal choice of $`M`$ we find
$$q_1=0.5.$$
(1.5)
It should be noted that the difference between $`q_0`$ and $`q_1`$ is $`O(1)`$. In the full one-loop result (eqs. (2.3) to (2.9) below) $`g^2`$ occurs as a pre-factor of relatively little importance.
Let us now explain the physical origin of $`q_1`$. Consider the free domain-wall propagator $`G_{s,s^{}}^0(p)`$ for a given four-momentum $`p`$ in the vicinity of the $`s=0`$ boundary at zero quark mass. The $`s`$-correlations described by this propagator are controlled by an exponent $`\alpha (p)`$. Each term in the propagator involves a factor $`\mathrm{exp}(d\alpha (p))`$ where $`d`$ stands for either the separation $`|ss^{}|`$ or the sum of distances from the boundary $`s+s^{}`$. For the standard domain-wall action one has $`\mathrm{max}\{\mathrm{exp}(\alpha )\}=0.5`$ for $`M=1`$ where the maximum over the Brillouin zone is obtained at the “corner” $`p_\pi =(\pi ,0,0,0)`$ and its three permutations. We will denote the set of global maxima by $`𝒫`$.
Now, at tree level, a fermion eigenmode with momentum $`p`$ propagates independently of all other eigenmodes. But for any non-zero gauge coupling the fermions propagate in non-trivial backgrounds, and these backgrounds allow any given momentum eigenmode to couple to all momentum eigenmodes. In particular, small-momentum quark modes couple to modes with $`p𝒫`$.
We arrive at the following physical picture. A four-dimensional fermion mode created on a given $`s`$-layer mixes on that layer with the modes of $`𝒫`$ through the gauge field. As a mode with $`p𝒫`$, the fermion propagates with minimal suppression to some other layer $`s^{}`$, where the action of the gauge field turns it back into the original mode. Propagation in the $`s`$-direction is therefore dominated by the modes of $`𝒫`$ leading to
$$q_1=\mathrm{max}\{\mathrm{exp}(\alpha (p))\}.$$
(1.6)
For $`M=1`$ this reduces to eq. (1.5). If we would momentarily regard the fifth direction as an imaginary-time direction, the above is recognized as the familiar result that propagation is always dominated by the lightest excitation in any given channel. The domain-wall case is particularly simple in that the gauge field is independent of the $`s`$-coordinate.
Under certain conditions (basically that the coupling constant is not too large, see Appendix C.2 for a more detailed discussion) it should be possible to describe the results of numerical simulations too in terms of an effective wave function $`\chi _{\mathrm{eff}}(s)s^{1\delta }q_{\mathrm{eff}}^s`$. This means that every quark’s wave function is assumed to be the product of a four-dimensional wave function and the universal fifth-coordinate wave function $`\chi _{\mathrm{eff}}(s)`$. The exponential fall-off rate is accounted for by $`q_{\mathrm{eff}}`$. At relatively weak coupling ($`a^1>\mathrm{\hspace{0.33em}3}\mathrm{GeV}`$) there seems to be no problem with the restoration of chiral symmetry, suggesting that $`q_{\mathrm{eff}}<1`$. For $`a^12\mathrm{GeV}`$, the rate at which chiral symmetry is restored depends sensitively on the observable. This, as well as other indications, suggest that $`q_{\mathrm{eff}}`$ is very close to one, and the restoration of chiral symmetry really follows a power-law behavior. (In ref. an estimate of $`q_{\mathrm{eff}}`$ was given which, however, is unjustified because the power-law correction was ignored.) Then, at $`a^11\mathrm{GeV}`$ the notion of a universal, localized, $`s`$-coordinate wave function breaks down.
Comparing the perturbative results with the numerical data shows that, not surprisingly, the optimal tree-level value $`q_0=0`$ completely fails to describe that data. In comparison, the one-loop result $`q_1=0.5`$ lies approximately “half-way” between the tree-level value and the close-to-one values of $`q_{\mathrm{eff}}`$ which seem to account for the results of simulations. Thus $`q_1`$ gives at least some indication of the actual behavior of the system.
In this paper we adopt $`q_1`$ as an analytic criterion for the quality of domain-wall actions. In Sec. 3 we consider new families of domain-wall actions involving beyond-nearest neighbor coupling. We compute the resulting $`q_1`$, and find that values much smaller than 0.5 can be achieved. Finally, in Sec. 4 we discuss the relevance of our results to numerical simulations.
Some technicalities of the one-loop calculation are relegated to Appendix A. Higher order corrections are briefly discussed in Appendix B. An expanded discussion of some non-perturbative issues can be found in Appendix C.
2. The one-loop wave function
In this section we calculate the one-loop wave function of domain-wall quarks, relegating some of the technicalities to Appendix A. The finite-$`N_s`$ tree-level propagator and the one-loop self energy for domain-wall fermions were calculated in ref. (see also ref. ). Here we will be interested in the range $`1sN_s`$, therefore we can use the simpler expressions for the tree-level propagator in the limit $`N_s\mathrm{}`$ . Assuming the right-handed quark is localized near the $`s=0`$ boundary, the singular part of the tree-level propagator $`G_{s,t}^0(p)`$ is
$$P_+\frac{M(2M)(1M)^{s+t}}{i/p}.$$
(2.1)
For $`M1`$ this becomes
$$P_+\frac{\delta _{s,0}\delta _{t,0}}{i/p},$$
(2.2)
which means that the massless right-handed fermion field is fully localized on the boundary layer.
The first quantum effect beyond the free theory is an additive correction, $`\delta M`$, to the five-dimensional mass $`M`$. It arises in a mean-field approximation, or in perturbation theory from tadpole diagrams. It is well known that this effect must be treated non-perturbatively . We thus discard the tadpole diagrams, absorbing them into the tree-level action via the replacement $`MM\delta M`$.
The full one-loop wave function is
$$\chi _1(s)=|1+\delta MM|^s+\delta \chi _1(s),$$
(2.3)
where $`\delta \chi _1(s)`$ comes from the “setting sun” diagram only. We will find $`\delta \chi _1(s)`$ by matching the non-analytic piece of the dressed propagators with the r.h.s. of eq. (1.1). We are interested in the behavior of $`\delta \chi _1(s)`$ when $`M`$ is close to its optimal mean-field value (see also Sec. 3). In the calculation below we thus set $`M=1+\delta M`$ in tadpole-improved perturbation theory. Since the tadpole-improved free propagator is a function of $`M\delta M`$ (and not of $`M`$ and $`\delta M`$ separately) the resulting propagator is identical to the ordinary propagator with $`M=1`$. The setting-sun diagram will therefore be computed using the expression for the ordinary tree-level propagator for $`M=1`$.
The setting-sun diagram is depicted in Fig. 1. Notice that we have not amputated the external legs . Except for $`s=t=0`$ (see (2.2)), the tree-level propagator is not singular at $`p=0`$. To obtain a contribution to the r.h.s. of eq. (1.1), at least one of the three fermion lines must coincide with expression (2.2). (The kinetic self-energy correction $`\mathrm{\Sigma }_K`$ in eq. (1.1) arises when all three lines coincide with eq. (2.2), see ref. , and will not be discussed here any further.) Assume first that the rightmost propagator in Fig. 1 coincides with eq. (2.2). Since the leftmost coordinate $`s`$ is by assumption far from the boundary, we may take the limit $`p0`$ in the expressions for the self-energy part and for the leftmost propagator. We thus arrive at
$$\delta \chi _1(s)=g^2C_2\underset{s^{}0}{}G_{s,s^{}}^{0+}\mathrm{\Sigma }_{s^{},0}^+=g^2C_2\mathrm{\Sigma }_{s1,0}^+,$$
(2.4)
where $`C_2`$ is the quadratic Casimir and
$$G_{s,t}^{0\pm }=\frac{1}{2}\mathrm{tr}P_\pm G_{s,t}^0(p=0),\mathrm{\Sigma }_{s,t}^\pm =\frac{1}{2}\mathrm{tr}P_\pm \mathrm{\Sigma }_{s,t}(p=0),$$
(2.5)
and $`\mathrm{\Sigma }_{s,t}(p)`$ is the 1PI self-energy obtained by amputating the external legs in Fig. 1. In the second equality of eq. (2.4) we used the explicit expression for $`G_{s,t}^{0+}`$ far from the boundary (see Appendix A).
When substituting eq. (2.3) into eq. (1.1) we find another term, $`\delta \chi _1(s^{})`$. This term is obtained when the leftmost propagator in Fig. 1 coincides with eq. (2.2). Following similar steps we now find
$$\delta \chi _1(s^{})=g^2C_2\mathrm{\Sigma }_{0,s^{}1}^{}.$$
(2.6)
Thanks to a “parity” symmetry (see Appendix A) one has $`\mathrm{\Sigma }_{s,t}^+=\mathrm{\Sigma }_{t,s}^{}`$. Hence eqs. (2.4) and (2.6) agree.
It remains to compute the diagonal part of the self-energy. One can write
$$\mathrm{\Sigma }_{s,0}^+=_\pi ^{+\pi }\frac{d^4k}{(2\pi )^4}h^+(k)\mathrm{exp}(s\alpha (k)).$$
(2.7)
(Note that the external momentum is zero.) The $`s`$-dependence enters through the exponential. All other factors were lumped into $`h^+(k)`$ (see Appendix A for more details). For large $`s`$, the above integral can be computed using a saddle-point approximation. As mentioned in the introduction, the global maximum of $`\mathrm{exp}(\alpha )`$ corresponds to the lattice momentum $`p_\pi `$ and its three permutations. In the computation we take $`h^+=h^+(p_\pi )`$ outside the integral, and expand the exponent to second order around $`p_\pi `$ where we define $`k=(\pi +k_{},\stackrel{}{k}_{})`$. Including a factor of four to account for the degeneracy of the global maximum we obtain
$`\mathrm{\Sigma }_{s,0}^+`$ $`=`$ $`{\displaystyle \frac{3}{2}}\left({\displaystyle \frac{1}{2}}\right)^s{\displaystyle \frac{dk_{}d^3k_{}}{(2\pi )^4}\mathrm{exp}\left(\frac{s}{24}(7k_{}^2+k_{}^2)\right)}`$ (2.8)
$`=`$ $`{\displaystyle \frac{54}{\pi ^2\mathrm{\hspace{0.17em}7}^{3/2}s^2\mathrm{\hspace{0.17em}2}^s}}.`$
For the fundamental representation of SU(3) one has $`C_2=4/3`$. Substituting in eq. (2.4) we finally find
$$\delta \chi _1(s)=g^2\frac{72}{\pi ^2\mathrm{\hspace{0.17em}7}^{3/2}(s1)^2\mathrm{\hspace{0.17em}2}^{s1}}g^2\frac{0.788}{s^2}\left(\frac{1}{2}\right)^s,s1.$$
(2.9)
Extrapolating eq. (2.9) to smaller values of $`s`$ suggests that as soon as (or shortly after) we move off the boundary layer, $`\delta \chi _1(s)`$ dominates over the tree-level term in eq. (2.3). This is true even if $`g^2`$ is small (or had the prefactor in eq. (2.9) been numerically small). The reason is that the relative magnitude of the two terms is proportional to $`(0.5/(1+\delta MM))^s`$, and since $`1+\delta MM1`$ this grows exponentially fast.
3. New actions
We have found that the $`s`$-coordinate’s wave function of domain-wall quarks is dominated by quantum effects. The arguments of Sec. 2 show that the broadening of the wave-function is controlled, in the one-loop approximation, by the maximum of $`\mathrm{exp}(\alpha (p))`$ over the Brillouin zone. This remains true for other domain-wall actions unless the $`s`$-couplings are drastically changed. So, if a different domain-wall action yields a smaller $`\mathrm{max}\{\mathrm{exp}(\alpha )\}`$, namely a faster fall-off of the wave function at the one-loop level, it is plausible that that new action also performs better non-perturbatively (we return to this issue in Sec. 4).
The standard domain-wall action contains two parameters, the domain-wall height $`M`$ and the Wilson parameter $`r`$ (which is usually set equal to one). The $`M`$-dependence of $`\mathrm{max}\{\mathrm{exp}(\alpha )\}`$ was investigated in ref. . As mentioned in Sec. 2, however, the additive renormalization of $`M`$ must be treated non-perturbatively. The optimal value used in simulations ($`M1.8`$) is nicely consistent with mean-field estimates. We will thus assume that the numerical optimization of $`M`$ corresponds to setting $`M=1+\delta M`$ in tadpole-improved perturbation theory. Again, this means that we should determine $`\mathrm{max}\{\mathrm{exp}(\alpha )\}`$ using the tree-level action with $`M=1`$. As for the Wilson parameter, changing its value in the standard domain-wall action turns out to have little effect (see below).
We will depart from the standard domain-wall action by allowing for couplings not only between nearest neighbors. In view of the obvious increase in computer time needed for the inversion of the fermion matrix, we try to be as economic as possible in our beyond-nearest neighbor excursion. We allow only for coupling between sites $`x`$ and $`x+n\widehat{\mu }`$ (but not e.g. for coupling between $`x`$ and $`x+\widehat{\mu }+\widehat{\nu }`$ for $`\mu \nu `$). In this paper we consider explicitly $`n=2`$ and $`n=3`$, namely next-nearest and next-next-nearest couplings in the same direction. Also the modifications will be restricted to the four-dimensional part of the action, leaving the coupling in the fifth direction intact. (Note that we are interested in achieving a fast fall-off in the $`s`$-direction; any attempt to generate a smoother, continuum-like, behavior in the $`s`$-direction is thus the exact opposite of what we are aiming for.)
The domain-wall operators considered here will have the following general form for zero quark mass
$$D_{s,t}^{\mathrm{d}.\mathrm{w}.}=\delta _{s,t}D+(\delta _{s+1,t}\delta _{s,t})P_++(\delta _{s1,t}\delta _{s,t})P_{},$$
(3.1)
with the understanding that on a finite lattice $`0s,tN_s1`$. The inclusion of a quark mass can be done in the usual way . The four-dimensional part of the action is
$$D(p)=i\underset{\mu }{}\gamma _\mu f(p_\mu )rW(p)+M.$$
(3.2)
This equation gives the tree-level operator in momentum space. The generalized Wilson term $`W(p)`$ is a function of $`\mathrm{cos}(p_\mu )`$. In the kinetic term, $`f(p_\mu )`$ is an odd function of its argument, which we take to be $`\mathrm{sin}(p_\mu )`$ times a polynomial in $`1\mathrm{cos}(p_\mu )`$. Later we will give explicit expressions for $`W(p)`$ and $`f(p_\mu )`$. As explained earlier we set $`M=1`$ in the tree-level action, but we will use the freedom in varying the Wilson parameter $`r`$. Our convention is that $`W(p)`$ and $`r`$ are both positive.
For any domain-wall action of the above form, the exponents $`\alpha (p)`$ are determined by
$$2\mathrm{cosh}(\alpha (p))=\frac{1+B^2(p)+\underset{\mu }{}f^2(p_\mu )}{B(p)},$$
(3.3)
where $`B(p)=1M+rW(p)`$ and $`\alpha (p)0`$ by convention. ($`B(p)=rW(p)`$ for $`M=1`$; for the standard domain-wall action eq. (3.3) reduces to eq. (A.8).) Lowering the global maximum of $`\mathrm{exp}(\alpha (p))`$ corresponds to raising the global minimum of eq. (3.3).
Some insight about the features that control $`\mathrm{min}\{2\mathrm{cosh}(\alpha )\}`$ can be obtained from very general considerations. One has $`W=0`$ for $`p=0`$, and in all cases one aims for $`rW>1`$ at $`p_\pi =(\pi ,0,0,0)`$. As we gradually increase $`p_1`$ for 0 to $`\pi `$ (keeping $`p_2=p_3=p_4=0`$) at some value $`p_c=(p_{1c},0,0,0)`$ we will have $`rW=1`$. Were it not for the $`f(p_1)`$ term in eq. (3.3), at $`p_c`$ we would obtain $`\mathrm{cosh}(\alpha )=\mathrm{exp}(\alpha )=1`$, namely no exponential suppression at all. To avoid this dangerous situation, we would like to have $`f^2(p_1)`$ as large as possible at $`p_1=p_{1c}`$.
Another danger lurks at the (fifteen non-zero) corners of the Brillouin zone. There, by construction, $`f(p_\mu )=0`$, and so $`\mathrm{exp}(\alpha )=(rW)^1`$. We will therefore also be interested in increasing $`rW`$ at the corners of the Brillouin zone.
As a warm-up exercise let us consider the effect of varying $`r`$ in the standard domain-wall action. In this case $`W=_\mu (1\mathrm{cos}(p_\mu ))`$ and $`f(p_\mu )=\mathrm{sin}(p_\mu )`$. For $`r1`$, $`\mathrm{min}\{2\mathrm{cosh}(\alpha )\}`$ occurs at $`p_\pi `$. We can increase $`rW`$ at $`p_\pi `$ by increasing $`r`$. But in that case the value $`rW=1`$ will occur at a smaller $`p_1`$, where $`\mathrm{sin}(p_1)`$ is smaller. As can be seen from Fig. 2a there is a transition region around $`r1.2`$. For larger values of $`r`$, $`\mathrm{min}\{2\mathrm{cosh}(\alpha )\}`$ moves towards the point where $`rW=1`$. As an example, for $`r=2.0`$ one has $`rW=1`$ at $`x=\mathrm{cos}(p_1)=0.5`$, and $`\mathrm{min}\{2\mathrm{cosh}(\alpha )\}`$ is at $`x0.42`$. The largest value of $`\mathrm{min}\{2\mathrm{cosh}(\alpha )\}`$, obtained for $`r1.3`$ – 1.5, is around 2.8. This makes little improvement over the value 2.5 obtained at $`r=1`$.
To avoid this dead-lock we take $`W`$ to be a non-linear function of $`\mathrm{cos}(p_\mu )`$. We define
$$W_n=\underset{\mu }{}(1\mathrm{cos}(p_\mu ))^n.$$
(3.4)
$`W_n`$ requires couplings of sites $`x`$ and $`x+n\widehat{\mu }`$. Once such coupling have been introduced into the generalized Wilson term, we allow them also in the kinetic term. Further raising of the global minimum of $`\mathrm{cosh}(\alpha )`$ will be made possible by choosing $`f(p_\mu )`$ that increases faster than $`\mathrm{sin}(p_\mu )`$, and by adjusting the Wilson parameter $`r`$.
We now turn to the investigation of concrete actions. The minimization problem was solved numerically. Note that the r.h.s. of eq. (3.3) can be expressed as a function of $`x_\mu \mathrm{cos}(p_\mu )`$ only. Using the invariance under permutations of the four components, it is enough to look for the global minimum over the range $`1x_1x_2x_3x_41`$. (One can also study the minimization problem analytically. For any $`\mu `$, $`\mathrm{sin}(p_\mu )=0`$ always satisfies the extremality condition. For all momentum-components where $`\mathrm{sin}(p_\mu )0`$ one finds a coupled algebraic equation in $`x_\mu `$. In all the cases we have studied, it turned out that the global minimum was either of the form $`(p_{\mathrm{min}},0,0,0)`$ or else of the form $`(p_{\mathrm{min}},p_{\mathrm{min}},p_{\mathrm{min}},p_{\mathrm{min}})`$.)
We first consider an action containing next-nearest neighbors
$$D_{23}=i\underset{\mu }{}\gamma _\mu f_3(p_\mu )rW_2+M,$$
(3.5)
$$f_3(p_\mu )=\mathrm{sin}(p_\mu )\left[1+c_3(1\mathrm{cos}(p_\mu ))\right].$$
(3.6)
$`W_2`$ is defined in eq. (3.4). In Table 1 we show the resulting values of $`\mathrm{min}\{2\mathrm{cosh}(\alpha )\}`$ and $`\mathrm{max}\{\mathrm{exp}(\alpha )\}`$ for several values of $`c_3`$. For each $`c_3`$ we looked for the best value of $`r`$ which we denote $`r_{\mathrm{opt}}`$. One sees that values of $`\mathrm{max}\{\mathrm{exp}(\alpha )\}`$ much smaller than 0.5 are feasible. A plot of $`2\mathrm{cosh}(\alpha )`$ for $`c_3=4/3`$ is shown in Fig. 2b. Notice the flatness of $`\mathrm{cosh}(\alpha )`$ for $`1\mathrm{cos}(p_1)0`$ at $`r=r_{\mathrm{opt}}=1.45`$.
We now digress to discuss how the present work relates to the standard “improvement program” (see e.g. the review ). In the study of the hadron spectrum, only a single parameter (the bare quark mass) in the fermion action needs to be tuned. Once the correct continuum limit has been established, attention is focused on eliminating those lattice artifacts that vanish most slowly, that is, linearly with the lattice spacing. However, in the calculation of weak matrix elements one has to first establish the correct continuum limit. This is very problematic with Wilson or staggered fermions because, due to the loss of full chiral and/or flavor symmetry, many subtraction coefficients must be tuned. Controlling those subtractions by having good chiral and flavor properties simultaneously is thus of higher priority than the removal of any other lattice error. Furthermore, in the massless-quark limit $`O(a)`$ lattice artifacts are automatically excluded if chiral symmetry is maintained . In that sense, approaching the chiral limit using domain-wall fermions encompasses the standard improvement program as well.
Coming back to the new domain-wall action, since the new Wilson term $`W_2`$ starts off at order $`p^4`$, the first lattice deviation from a relativistic (tree-level) dispersion relation comes only from the kinetic term. This is shown in the second column of Table 1. We observe that while increasing $`c_3`$ from zero to $`1/3`$ improves the dispersion relation, the opposite is true for $`c_3>1/3`$. Although the error is formally of order $`a^2`$, it might become significant if $`c_3`$ is too large. To gain some idea on the magnitude of the error consider, say, $`p^2(400\text{MeV})^2`$, which is relevant for kaon physics, on a lattice with $`a^12\text{GeV}`$. This means $`a^2p^21/25`$. For the last two rows of Table 1, the effect is 2% and 4% respectively.
If next-next-nearest neighbors in the same direction are also allowed one can decrease $`\mathrm{max}\{\mathrm{exp}(\alpha )\}`$ further while maintaining a vanishing $`p^3`$ term. Let
$$D_{35}=i\underset{\mu }{}\gamma _\mu f_5(p_\mu )rW_3+M,$$
(3.7)
where again $`W_n`$ is defined in eq. (3.4) and where
$$f_5(p_\mu )=\mathrm{sin}(p_\mu )\left[1+\frac{1}{3}(1\mathrm{cos}(p_\mu ))+c_5(1\mathrm{cos}(p_\mu ))^2\right].$$
(3.8)
Some values of $`\mathrm{min}\{2\mathrm{cosh}(\alpha )\}`$ and $`\mathrm{max}\{\mathrm{exp}(\alpha )\}`$ are shown in Table 2. A plot of $`2\mathrm{cosh}(\alpha )`$ for $`c_5=2`$ is shown in Fig. 2c. Even for the last row in Table 2 ($`c_5=50`$), the deviation from Lorentz covariance is at the level of $`(50/4)(a^2p^2)^22\%`$ for $`p^2(400\text{MeV})^2`$. For $`c_5=5`$, the deviation is below 2% up to $`(700\text{MeV})^2`$, and so on.
We conclude with a number of technical comments. Replacing the four-dimensional part of the domain-wall action by $`D_{23}`$ ($`D_{35}`$) approximately doubles (triples) the number of entries in the fermion matrix. Therefore one should expect a corresponding increase in the cost of a single inversion of the fermion matrix at fixed $`N_s`$.
In the continuum limit, both the standard domain-wall action and the new actions discussed above support a single quark (one Weyl field on each boundary) for $`|1M|<1`$. In the case of the standard action there is a four-quark zone (corresponding to the corner $`p_\pi `$ and its permutations) for $`|3M|<1`$. When the Wilson term $`W_n`$ is employed instead, the four-quark zone is at $`|1+r\mathrm{\hspace{0.17em}2}^nM|<1`$. An additional benefit of the new actions is that the four-quark and the single-quark zones are separated by a large gap (as a function of $`M`$) for $`rr_{\mathrm{opt}}`$. We expect that a clear gap should be found in simulations too, even though its precise location will likely be different from the weak-coupling limit.
For the standard domain-wall action, the optimal value of $`M`$ used in simulations agrees well with the mean-field estimate of $`1+\delta M`$. One obtains $`\delta M`$ by substituting a mean value $`u`$ for each link variable in the Wilson term. We will assume that the new actions are gauged in the simplest way, namely using only products of link variables along straight lines. (E.g. the sites $`x`$ and $`x+2\widehat{\mu }`$ are connected via $`U_{x,\mu }U_{x+\widehat{\mu },\mu }`$ etc.) Using a mean link $`u0.8`$ at $`\beta =6.0`$ (see e.g. ref. ), the mean-field estimate is $`2r(34u+u^2)0.9r`$ for $`D_{23}`$ and $`r(1015u+6u^2u^3)1.3r`$ for $`D_{35}`$.
Last, for any domain-wall operator with the form of eqs. (3.1) and (3.2), the $`N_s\mathrm{}`$ limit defines an overlap-Dirac operator obeying the Ginsparg-Wilson relation (for a review see ref. ) given by
$$D_{\mathrm{GW}}=1\gamma _5ϵ(\gamma _5D),$$
(3.9)
where $`ϵ(x)=\pm 1`$ is the sign function acting on each of the eigenvalues of $`\gamma _5D`$.
4. Discussion
In this paper we showed that the one-loop wave function of domain-wall fermions behaves like $`s^2q_1^s`$, where $`q_1`$ is determined by the free-fermion action. This is true for a wide class of domain-wall actions, including those defined in eq. (3.1). For the standard action $`q_1=0.5`$, whereas the addition of beyond-nearest neighbor couplings allows for much smaller values of $`q_1`$.
Only in the weak-coupling limit does $`q_1`$ fully control the wave function. For finite $`g^2`$ up to some critical value $`g_c^2`$ we expect the wave function to be proportional to some $`q_{\mathrm{pt}}^s`$ (up to power corrections) with $`q_{\mathrm{pt}}=q_{\mathrm{pt}}(g^2)`$. The $`g^2`$ dependence can be parametrized in various ways. In Appendix B we consider the role of higher-order diagrams, and the parametrization
$$q_{\mathrm{pt}}^s=q_1^s\mathrm{exp}(s(g^2\eta _1+g^4\eta _2+\mathrm{})),$$
(4.1)
is found to be natural. The Taylor expansion of $`\mathrm{exp}(sg^2\eta _1)`$ corresponds to a family of 1PI diagrams of all orders, where the first $`\eta _1`$-dependent terms are two-loop diagrams. (Analogous statements apply to $`\eta _2`$ etc.) In terms of $`q_{\mathrm{pt}}(g^2)`$, one can define $`g_c^2`$ by the condition $`q_{\mathrm{pt}}(g_c^2)=1`$. The existing numerical results suggest that the (quenched) value of $`6/g_c^2`$ is very close to 6.0 for the standard domain-wall action.
If both $`q_1`$ and $`\eta _1`$ were known for a given action, one could obtain a crude estimate of $`g_c^2`$ via a linear extrapolation. We have computed only the $`q_1`$ values, so we can only conjecture what trends are likely to affect $`g_c^2`$. First, in the lower rows in Table 2, $`q_1`$ is extremely small. Nevertheless, if $`\eta _1`$ is large (and positive), $`g_c^2`$ may end up being approximately the same as (or, for that matter, even smaller than) for the standard domain-wall action. Because of gauge invariance there are vertices that depend linearly on $`c_5`$. Since the lower rows in Table 2 come from actions with numerically large values of $`c_5`$ this should, indeed, lead generically to a large $`\eta _1`$ (and $`\eta _2`$ and so on).
It is therefore safer to focus on the first few rows in Tables 1 and 2, where one is less prone to the above risk. The following heuristic argument suggests that, in that range of parameters, the new domain-wall actions may indeed be superior to the standard one. When the Wilson parameter is equal to $`r_n2^{1n}`$, an action containing the Wilson term $`W_n`$ (eq. (3.4)) gives rise to $`\mathrm{min}\{2\mathrm{cosh}(\alpha )\}=2.5`$ (corresponding to $`q_1=0.5`$) at $`p_\pi =(\pi ,0,0,0)`$ for $`M=1`$. The behavior at $`r=r_n`$ is therefore a common starting point over which we may try to improve by increasing $`r`$. As discussed in Sec. 3, for the standard action $`\mathrm{min}\{2\mathrm{cosh}(\alpha )\}`$ is relatively insensitive to $`r`$. Its largest value (which is 2.8) is obtained around $`r_{\mathrm{opt}}1.3`$ – 1.5, namely $`r_{\mathrm{opt}}`$ is less than 50% above $`r_1`$. In comparison, for $`D_{23}`$ at $`c_3=4/3`$ (Fig. 2b) the largest value of $`\mathrm{min}\{2\mathrm{cosh}(\alpha )\}`$ is achieved at $`r_{\mathrm{opt}}=1.45`$ which is approximately three times $`r_2`$. For $`D_{35}`$ at $`c_5=2`$ (Fig. 2c) the best value is $`r_{\mathrm{opt}}=1.3`$ which is more than five times $`r_3`$.
The ability to reach larger values of $`\mathrm{min}\{2\mathrm{cosh}(\alpha )\}`$ is thus correlated with an enhanced sensitivity to the Wilson parameter, and with a bigger ratio $`r_{\mathrm{opt}}/r_n`$. Now, while $`\eta _1,\eta _2,\mathrm{},`$ might in principle grow as $`r`$ increases from $`r_n`$ to $`r_{\mathrm{opt}}`$, it is clear that as functions of the parameters of the theory their behavior will be very different from $`q_1`$ (see Appendix B). Therefore it is plausible that there exist “windows” of parameters where $`q_{\mathrm{pt}}(g^2)`$ is controlled primarily by the decreasing $`q_1`$, implying that the exponential suppression holds up to a larger value of $`g^2`$.
In Appendix C.2 we discuss how different ways of approaching the chiral limit are related to different forms of the spectral-density function of the (normal-ordered) transfer matrix. At weak coupling one expects to have a gap, namely almost all eigenvalues are smaller than some $`\lambda _0<1`$. The gap region $`\lambda _0\lambda 1`$ is not completely devoid of eigenvalues, but their total number is drastically smaller than just below $`\lambda _0`$. One also expects a big difference between the corresponding eigenfunctions. Those that lie outside the gap should be continuum-like modes that spread all over the lattice, while inside the gap the modes should be highly localized .
As the coupling constant increases a qualitative change takes place. Near-unity eigenvalues of the transfer matrix proliferate. For the standard domain-wall action this is a direct consequence of the proliferation of approximate zero modes of the hermitian, four-dimensional, Wilson-Dirac operator . The change (which seems to take place around quenched $`\beta =6.0`$) shows the key features of the phenomenon known in condensed matter as localization . Due to the randomness of generic gauge-field configurations, in any given part of the lattice there is a finite probability to find a localized (approximate) zero mode.
Viewing the (hermitian) Wilson-Dirac operator as a hamiltonian, under its action the fermions can hop only a single site. But with the new domain-wall actions the relevant hamiltonian is $`\gamma _5D`$ (see eq. (3.2)). Now the fermions may hop also two (or three) sites when the hamiltonian acts just once on a given state. It should be more difficult to trap the new fermions inside a small potential well, as now they have more ways of escaping out of it! This consideration too suggests that the critical coupling, where the exponential suppression is lost, may be larger for the new actions.
The new domain-wall actions considered in this paper carry with them an obvious extra cost for a single inversion of the fermion matrix. One may hope to reduce chiral symmetry violations also by using improved gauge actions, because the latter tend to generate smoother configurations. If this goal is achieved, it may be a numerically much cheaper way to reduce chiral symmetry violations. Using the Iwasaki action it was found that the residual large-$`N_s`$ pion-mass squared (extrapolated to zero quark mass) drops by about a factor of two for $`a^11\mathrm{GeV}`$ . However, the new residual pion mass is still very big. Also, in thermodynamics, the Iwasaki action did not lead to any noticeable reduction in the residual pion mass (for a detailed discussion of various improvements see the first paper of ref. ).
Our analysis suggests a possible explanation why the use of improved gauge actions has had only a limited success. If the (tree-level) domain-wall action is unchanged, the one-loop wave function (1.4) is still controlled by the same value of $`q_1=0.5`$. Only the numerical prefactor may change (cf. eq. (2.9)). Hence, in this approximation, the exponential fall-off rate is not getting any better for improved gauge actions. (The same reasoning applies to the use of “fat links”. One has to be careful, however, because this argument ignores higher-order corrections, cf. eq. (4.1).)
In conclusion, in the one-loop approximation even a small reduction in $`q_1`$ leads to a dramatic suppression of chiral symmetry violations for commonly used values of $`N_s`$. If the actual quark’s wave function is (even partly) correlated with the one-loop one, the new actions could give rise to a significantly better chiral behavior, enough to justify their increased simulation cost.
Acknowledgements
The key results of this paper were worked out during a visit to Brookhaven National Laboratory last summer. Extensive discussions with many members of the lattice groups in Brookhaven Lab. and Columbia Univ. were essential for the development of this work. In forming my view on the current state of domain-wall fermion simulations, as presented in this paper, I benefitted from discussions with the participants of the Dubna workshop on “Lattice Fermions and Structure of the Vacuum.” I thank Tom Blum, Maarten Golterman and Karl Jansen for their useful comments on this paper, and Karl Jansen also for discussing with me some as-yet unpublished results. This research is supported in part by the Israel Science Foundation.
Appendix A. Some technicalities
In this appendix we collect a few useful formulae. The inverse tree-level propagator in momentum space is
$$(G^0)_{s,t}^1=i\delta _{s,t}\underset{\mu }{}\gamma _\mu \mathrm{sin}(p_\mu )+W_{s,t}^+(p)P_++W_{s,t}^{}(p)P_{},$$
(A.1)
where
$$W_{s,t}^+(p)=\delta _{s+1,t}B(p)\delta _{s,t},$$
(A.2)
$$B(p)=1M+\underset{\mu }{}(1\mathrm{cos}(p_\mu )),$$
(A.3)
and $`W_{s,t}^{}(p)=W_{t,s}^+(p)`$. We consider only the case of a zero quark mass. Also, for the calculation of the self-energy we set $`M=1`$ in the tree-level action (see Sec. 2). The tree-level propagator was computed in ref. . At the corners of the Brillouin zone ($`\mathrm{sin}(p_\mu )=0`$, all $`\mu `$) the two chiralities decouple in eq. (A.1). For $`p0`$ the limit is singular. But at the other fifteen corners the limit is regular, leading to
$$G_{s,t}^0=(W^+)_{s,t}^1(p)P_++(W^{})_{s,t}^1(p)P_{}.$$
(A.4)
Specifically at $`p=p_\pi `$ and with a semi-infinite coordinate $`s,t=0,1,2,\mathrm{}`$ one has
$$(W^+)_{s,t}^1=(W^{})_{t,s}^1=\frac{1}{2}\theta (ts)\mathrm{\hspace{0.17em}2}^{st},$$
(A.5)
where $`\theta (ts)=1`$ for $`ts`$ and $`\theta (ts)=0`$ for $`t<s`$.
The explicit expression for the diagonal self-energy $`\mathrm{\Sigma }_{s,t}^+`$ of eq. (2.5) is given by eq. (36) of ref. . With a slight change of notation it reads (recall that the external momentum is zero)
$`\mathrm{\Sigma }_{s,t}^+`$ $`=`$ $`{\displaystyle _\pi ^{+\pi }}{\displaystyle \frac{d^4k}{(2\pi )^4}}\left(4{\displaystyle \underset{\nu }{}}\mathrm{sin}^2(k_\nu /2)\right)^1{\displaystyle \underset{\mu }{}}(\mathrm{cos}^2(k_\mu /2)(W^+G^{})_{s,t}`$ (A.6)
$`\mathrm{sin}^2(k_\mu /2)(W^{}G^+)_{s,t}+{\displaystyle \frac{1}{2}}\mathrm{sin}^2(k_\mu )(G^++G^{})_{s,t}),`$
where
$$(G^\pm )^1(p)=\underset{\mu }{}\mathrm{sin}^2(p_\mu )+W^\pm (p)W^{}(p).$$
(A.7)
Explicit expressions for $`G_{s,t}^\pm `$ can be found in refs. . As explained in Sec. 2 we are interested in $`\mathrm{\Sigma }_{s,0}^+`$ for $`s1`$. In the saddle-point approximation we set the internal momentum on the fermion line to $`p_\pi `$ (or its permutations) in all terms, except in the exponential $`\mathrm{exp}(s^{}\alpha )`$ that occurs inside $`G_{s^{},0}^\pm `$, which is expanded to second order around $`p_\pi `$ using the definition
$$2\mathrm{cosh}(\alpha (p))=\frac{1+B^2(p)+\underset{\mu }{}\mathrm{sin}^2(p_\mu )}{B(p)}.$$
(A.8)
This expansion gives rise to the integrand in eq. (2.8). Since $`\mathrm{sin}(p_\mu )=0`$ for all $`\mu `$ at $`p_\pi `$, the last term in eq. (A.6) is zero. At $`p_\pi `$ one has $`W^\pm G^{}=(W^{})^1`$. Since the matrices $`W^\pm `$ and their inverses are triangular, the second term in eq. (A.6) gives zero too. Only the first term in eq. (A.6) contributes, and only for $`\mu =2,3,4`$, leading to eq. (2.8).
We also need the diagonal tree-level propagator $`G^{0+}`$ (eq. (2.5)) away from the boundaries at $`p=0`$. For $`M1`$ and $`p0`$, the second-order operators $`(W^+W^{})_{s,t}`$ and $`(W^{}W^+)_{s,t}`$ tend to $`\delta _{s,t}`$ (except $`(W^{}W^+)_{s,t}`$ at $`t=s=0`$). As a result $`G_{s,t}^{0+}=\delta _{s1,t}=W_{s,t}^{}`$. (For a semi-infinite $`s`$-coordinate $`W^{}`$ is a right-inverse of $`W^+`$.)
Finally we observe that the interacting domain-wall action has a generalized parity symmetry. Writing the action as
$$\underset{s,s^{};\stackrel{}{x},\stackrel{}{y};x_4,y_4}{}\overline{\psi }_{s,\stackrel{}{x},x_4}D_{s,s^{};\stackrel{}{x},\stackrel{}{y};x_4,y_4}(U)\psi _{s^{},\stackrel{}{y},y_4}$$
(A.9)
one has
$$D_{s,s^{};\stackrel{}{x},\stackrel{}{y};x_4,y_4}(U)=\gamma _4D_{s^{},s;\stackrel{}{x},\stackrel{}{y};x_4,y_4}(U^{})\gamma _4,$$
(A.10)
where
$`U_{\stackrel{}{x},x_4;4}^{}`$ $`=`$ $`U_{\stackrel{}{x},x_4;4}`$
$`U_{\stackrel{}{x},x_4;k}^{}`$ $`=`$ $`U_{\stackrel{}{x}\widehat{\kappa },x_4;k}^{},k=1,2,3.`$ (A.11)
This looks like an ordinary parity transformation, except that we have switched the fifth coordinates of $`\psi `$ and $`\overline{\psi }`$ ($`s`$ and $`s^{}`$). The above discrete symmetry implies
$$\mathrm{\Sigma }_{s,t}(\stackrel{}{p},p_4)=\gamma _4\mathrm{\Sigma }_{t,s}(\stackrel{}{p},p_4)\gamma _4$$
(A.12)
and when $`\mathrm{sin}(p_\mu )=0`$ for all $`\mu `$, one has $`\mathrm{tr}P_+\mathrm{\Sigma }_{s,t}=\mathrm{tr}P_{}\mathrm{\Sigma }_{t,s}`$.
Appendix B. Beyond one loop
In this Appendix we consider the role of higher-order corrections. Specifically the aim is to show how an $`O(g^2)`$ correction to $`q_1`$ is built. The new wave-function is conveniently parametrized as
$`\delta \chi _{\mathrm{pt}}(s)`$ $``$ $`s^2q_{\mathrm{pt}}^s,`$
$`q_{\mathrm{pt}}^s=\left(q_1\mathrm{exp}(g^2\eta _1)\right)^s`$ $`=`$ $`q_1^s\left(1+g^2\eta _1s+{\displaystyle \frac{1}{2}}(g^2\eta _1)^2s^2+\mathrm{}\right).`$ (B.1)
The last expression suggests that, in the wave function $`\delta \chi _{\mathrm{pt}}(s)`$, the $`O(g^2)`$ correction to $`q_1`$ arises from a resummation of perturbation theory. If true, at any finite order we should find contributions to the wave function whose structure is $`\delta \chi _1(s)`$ times an increasing power of $`s`$. (We likewise expect additional $`O(g^4)`$ etc. corrections to $`q_1`$, cf. eq. (4.1); the arguments below are, however, too crude to tell how the power-law part of the wave function depends on $`g^2`$.)
We believe that a term $`g^2\eta s\delta \chi _1(s)`$ arises from the two-loop diagrams in Figs. 3a and 3b (and not from Figs. 3c or 3d, see below). Following Sec. 2, we will consider for definiteness the case where the rightmost fermion line in Fig. 3a corresponds to the singular part of the tree-level propagator, eq. (2.2), and therefore the fifth coordinate of the rightmost vertex is zero. In Fig. 3a the momenta on all three internal fermion lines can be simultaneously equal (or close) to $`p_\pi `$ (recall that the external momentum is (close to) zero). Then, the one-loop exponent $`q_1=\frac{1}{2}`$ comes with a power $`|st^{}|+|t^{}t|+t`$ which, for given $`s`$, is minimal when the points are ordered: $`st^{}t0`$. (Here we have ignored the difference between $`s`$ and $`s1`$ which is negligible for $`s1`$.) When the points are not ordered we obtain an exponentially convergent series in the excessive length of the fermion’s trajectory. For simplicity we will assume that the points are ordered, as we are only interested here in arguing that a contribution proportional to $`g^2s\delta \chi _1(s)`$ exists. (However, the full series has to be summed in order to obtain the numerical value of $`\eta `$.)
Next consider the integration over the loop momentum of the inner loop. As explained in Sec. 2, this integration should give rise to a factor of $`|t^{}t|^2`$. Since the sum $`_mm^2`$ is convergent, as a crude approximation one can say that the points $`t`$ and $`t^{}`$ are forced to be close together. Once this extra constraint has been taken into account, the remaining expression is independent of $`t`$ (and $`t^{}`$). Therefore the $`t`$-summation (approximated by an integral) gives $`_0^s𝑑t=s`$. The gaussian integration over the momentum of the outer loop gives rise roughly to the factor $`s^2`$ (associated with $`\delta \chi _1(s)`$) as before. We conclude that a term proportional to $`g^2s\delta \chi _1(s)`$ may indeed arise from Fig. 3a. (A technical complication is that, if the two loop-momenta are both equal to $`p_\pi `$, the momentum flowing through the inner gauge-boson line is zero. The singularity $`(kp)^2`$ is integrable in four dimensions, but its existence makes the actual calculation quite complicated. We believe that the above considerations are robust enough to grasp the dominant behavior of Fig. 3a.)
The above argument easily generalizes to the higher-loop diagrams of Fig. 4a. At the $`n`$th-order one has $`n1`$ “inner” loops. The fifth coordinates are pair-wise close, but otherwise are constrained only by ordering. Hence one expects a contribution proportional to
$$_0^s𝑑t_1_0^{t_1}𝑑t_2\mathrm{}_0^{t_{n2}}𝑑t_{n1}=\frac{s^{n1}}{(n1)!},$$
(B.2)
in agreement with eq. (B.1). A similar reasoning implies that the two-loop and $`n`$-loop diagrams of Figs. 3b and 4b respectively also contribute to the Taylor series in eq. (B.1).
We now want to explain why the diagrams of Figs. 3c and 3d do not contribute to the r.h.s. of eq. (B.1). Consider first the reducible diagram Fig. 3d. The momentum on the middle fermion line is zero (being equal to the external momentum). Hence the points $`t`$ and $`t^{}`$ are very close (see Sec. 2 and Appendix A). Ignoring the difference between $`t`$ and $`t^{}`$ the two gaussian integrations give rise to the product $`|st|^2t^2`$. The $`t`$-summation is then dominated by $`t`$-values which are either close to zero or to $`s`$ (being the fifth coordinates of the rightmost and leftmost vertices). Therefore the result behaves like $`s^2`$ (and not like $`s^2s=s^1`$ as in the case of Fig. 3a).
While the diagram in Fig. 3c is irreducible, the momenta on the internal fermion lines cannot be all equal to $`p_\pi `$. For example, if the momenta on the first and third internal lines are equal to $`p_\pi `$, then the momentum on middle line is zero. Hence, again, the points $`t`$ and $`t^{}`$ must be very close together, as well as close to either zero or $`s`$, and the final result is proportional to $`s^2`$ as in the case of Fig. 3d.
Appendix C. Some non-perturbative observations
C.1. Chiral-symmetry restoration
We outline here the proof of chiral-symmetry restoration given in ref. in the light of later works (in particular ref. ). For definiteness we focus on the anomalous term in the lattice PCAC relation. (Its vanishing implies that the pion mass will be zero after taking the infinite-volume and massless-quark limits in that order.) The PCAC relation using domain-wall fermions reads
$$\mathrm{\Delta }_\mu A_{5\mu }^a(x)J_5^b(y)=2m_0J_5^a(x)J_5^b(y)+2J_{5q}^a(x)J_5^b(y)+\text{contact term}.$$
(C.1)
Here $`\mathrm{\Delta }_\mu `$ is the backward lattice derivative and $`m_0`$ the bare quark mass. $`A_{5\mu }^a(x)`$ is the Noether current of a lattice transformation that assigns opposite charges to fermions on the half-spaces $`0s<N_s/2`$ and $`N_s/2s<N_s`$. This reduces to a chiral transformation on the quark states, as long as they are localized in their respective half-spaces. The pseudo-scalar density $`J_5^a(x)`$ is composed of fermion variables situated on the two boundaries, and serves as the standard interpolating field for pions. The anomalous (lattice-artefact) term in this relation involves another pseudo-scalar density, $`J_{5q}^a(x)`$, which is localized on two $`s`$-layers exactly half-way in the fifth direction. (We are assuming an undoubled quark spectrum (one quark field for each five-dimensional fermion field, in most applications) which is true at weak coupling. As mentioned in the introduction, the massless spectrum changes at strong coupling in which case the above transformations are no longer chiral.)
Before we proceed with the discussion of the anomalous term we have to address a technical point. Let $`T=_i|it_ii|`$ be the spectral decomposition of the (positive) first-quantized transfer matrix. We define a “normal-ordered” transfer matrix $`Q`$ via its spectral decomposition $`Q=_i|i\lambda _ii|`$ where $`\lambda _i=\mathrm{min}\{t_i,t_i^1\}`$. Physically, the operation of replacing $`t_i>1`$ by its inverse amounts to filling the Dirac sea. The spectrum of $`Q`$ lies in the interval $`0\lambda _i1`$. One can show that the exact domain-wall propagator in a given background field is a sum of terms, each of which involves the matrix $`Q`$ raised to a positive power which is a function of the fifth coordinates.
Let us now assume that, for a given background field, the spectrum of $`Q`$ lies in the interval $`0\lambda _i\lambda _0`$ where $`\lambda _0<1`$. Coming back to the anomalous term in eq. (C.1), for a non-singlet current it involves the propagation of two fermions over an $`s`$-separation equal to $`N_s/2`$ (or $`N_s/2\pm 1`$). The anomalous term is thus bounded by $`(\lambda _0^{N_s/2})^2=\lambda _0^{N_s}`$ times a constant. (Using the second-quantized transfer matrix formalism one can show that the proportionality constant is finite, being the norm of a product of bounded operators .) Moreover, if $`Q`$ has no eigenvalues larger than $`\lambda _0`$ for all gauge fields (a condition which is satisfied for a constrained gauge action ) one finds that the anomalous term falls exponentially after the functional averaging over the gauge field.
In ref. we showed that a very weak bound on the anomalous correlator exists even if there is no gap at all. The point is that exact-unity eigenvalues of the transfer matrix (hence of its normal-ordered version, $`Q`$, too) exist only on a submanifold of the lattice gauge-field space defined by the condition $`\mathrm{det}D_W=0`$ where $`D_W`$ is the hermitian four-dimensional Wilson-Dirac operator. As explained above, when a fermion propagates across an $`s`$-separation $`N_s/2`$, the propagator is bounded by (and, generically, falls like) $`\lambda _0^{N_s/2}`$ where $`\lambda _01`$ is the largest eigenvalue of $`Q`$. But $`\lambda _0^{N_s/2}`$ is negligible unless $`\lambda _0=1O(1/N_s)`$. This condition, in turn, will be satisfied only for gauge-field configurations whose distance from the above submanifold does not exceed $`O(1/N_s)`$. The volume of the (compact) gauge-field subspace contributing to the anomalous correlator is therefore finite, and shrinks like $`1/N_s`$, implying a similar bound on the correlator itself. We comment that the restoration of chiral symmetry is consistent with the fact that the overlap operator defined by the $`N_s\mathrm{}`$ limit (see eq. (3.9)) admits Lüscher’s chiral symmetry (whose generators are functions of the gauge field).
C.2. Spectral density and effective wave function
In the infinite-volume limit one can define a spectral function $`\rho _Q(\lambda )`$ associated with the normal-ordered transfer matrix $`Q`$ introduced in the previous subsection, whose support is (contained in) the interval $`[0,1]`$. Here we will not attempt to compute any spectral function. Instead, we adopt a “phenomenological” point of view. We will assume that a single, continuous, spectral density function $`\rho (\lambda )`$ has been defined as a suitable configuration average of $`\rho _Q(\lambda )`$. Using the considerations of the previous subsection, the aim is to see how different forms of the spectral function lead to different ways of approaching the chiral limit. One can envisage three prototype scenarios which are listed below. (A recent treatment of domain-wall fermions based on spectral integrals can be found in ref. . A numerical study of a spectral quantity which is closely related to $`\rho (1)`$ can be found in ref. . See also ref. .)
1) Exponential suppression. Assume that the support of $`\rho (\lambda )`$ is the interval $`0\lambda \lambda _0`$ where $`\lambda _0<1`$ and that, close to $`\lambda _0`$, $`\rho (\lambda )`$ vanishes like $`(\lambda _0\lambda )^\delta `$ with $`\delta >0`$. Consider the propagation of a single fermion from the boundary layer $`s^{}=0`$ to some other layer $`s`$ (We assume $`1s<N_s/2`$). This involves the integral
$$_0^{\lambda _0}𝑑\lambda \rho (\lambda )\lambda ^s\lambda _0^s_0^{\lambda _0}𝑑\lambda \rho (\lambda )\mathrm{exp}\left(s\frac{\lambda \lambda _0}{\lambda _0}\right)s^{1\delta }\lambda _0^s.$$
(C.2)
In order to obtain the power-law correction we have used the assumed behavior of $`\rho (\lambda )`$ close to $`\lambda _0`$, and wrote $`(\lambda /\lambda _0)^s=\mathrm{exp}(s\mathrm{log}(1+(\lambda \lambda _0)/\lambda _0))`$. (The one-loop result of Sec. 2 corresponds to $`\lambda _0=1/2`$ and $`\delta =1`$.) If $`n`$ fermions propagate across a similar $`s`$-interval, we will obtain the factor $`s^{1\delta }\lambda _0^s`$ for each of them. We may therefore consider $`\chi _{\mathrm{eff}}(s)=s^{1\delta }\lambda _0^s`$ as the effective $`s`$-coordinate wave function for all quark states. Since the anomalous divergence $`J_{5q}^a(x)`$ is a fermion bilinear, chiral symmetry violations should fall like $`\chi _{\mathrm{eff}}^2(N_s/2)N_{s}^{}{}_{}{}^{2(1+\delta )}\lambda _0^{N_s}`$. (The effective wave function describing the propagation of $`n`$ fermions could be somewhat different from $`\chi _{\mathrm{eff}}^n(s)`$ due to interactions between the different particles. Since chiral symmetry violations are related to $`J_{5q}^a(x)`$, the relevant effective wave function is always the one extracted from the sector with one fermion and one antifermion.)
2) Power-law suppression. Assume that $`\lambda _0=1`$ but with a vanishing $`\rho (1)`$, namely $`\rho (\lambda )(1\lambda )^\delta `$ for $`\lambda 1`$ with $`\delta >0`$. In that case the result of the spectral integral (C.2) will be $`s^{1\delta }`$. We might still speak of an effective wave function $`\chi _{\mathrm{eff}}(s)=s^{1\delta }`$ and expect chiral symmetry violations to fall like $`N_{s}^{}{}_{}{}^{2(1+\delta )}`$.
3) (Almost) no suppression. Last assume that $`\lambda _0=1`$ and that $`\rho (1)`$ is non-zero. Remember now the submanifold discussed in the previous subsection of gauge fields supporting an eigenvalue one of $`Q`$. If $`\rho (1)`$ is finite, configurations close to that submanifold must have a non-negligible Boltzmann weight. As we have explained, in this case the only suppression of long-range $`s`$-correlations comes from phase space considerations (the need to pick a configuration located $`O(1/N_s)`$ away from that submanifold). As a result, chiral symmetry violations fall roughly like $`\rho (1)/N_s`$, and the concept of a localized, effective wave function breaks down. ($`n`$-fermion correlations fall like $`1/N_s`$ too, and not like $`1/N_{s}^{}{}_{}{}^{n}`$.)
For clarity, we have presented above the three mathematically distinct scenarios. In reality, however, one is likely to encounter a more complicated behavior, characterized by the existence of a crossover region. Let us reexamine the exponential suppression scenario. As discussed in the previous subsection, in the ensemble of all gauge-field configurations (as opposed to the case where the plaquette is constrained to be everywhere small ) there is a non-zero probability of finding an eigenvalue arbitrarily close to one, for any $`g>0`$. Consider the integrated spectral density $`=_{\lambda _0}^1𝑑\lambda \rho (\lambda )`$, and suppose that $``$ is comparable to $`\lambda _0^{N_0}`$ for some $`N_0`$ (up to power corrections). For $`N_s<N_0`$, chiral symmetry violations will fall like $`N_{s}^{}{}_{}{}^{2(1+\delta )}\lambda _0^{N_s}`$. The point is that for any $`\lambda _0`$ significantly smaller than one, and $`N_0`$ of the order of (few times) ten, $`\lambda _0^{N_0}`$ will be so small that one will never have to use $`N_s>N_0`$ in a simulation. For practical purposes this scenario is therefore the same as the purely-exponential suppression scenario. If, nevertheless, very large values of $`N_s`$ will be tried, then around $`N_sN_0`$ a crossover to some slower fall-off rate will be encountered. (For instance, a crossover to a slower exponential fall-off rate has been observed in the Schwinger model .)
Last consider the relation between the domain-wall actions discussed in Sec. 3 and the associated overlap operators (see eq. (3.9)). The spectral function that controls the approach to the chiral limit of domain-wall fermions also controls the localization range of the overlap operator. As is clear from the above discussion, problems start when there is no gap, namely when it is not possible to identify a range $`\lambda _0\lambda 1`$ (with $`\lambda _0<1`$) where the eigenvalue density is (practically) zero. In that case both the overlap operator and the generators of the associated Lüscher symmetries become non-local. This is the counter-part of the loss of exponential suppression in the domain-wall case.
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# 1 Introduction
## 1 Introduction
The fate of the $`U(1)_A`$ symmetry of QCD at finite temperature is a fascinating problem which could also have interesting consequences for the ongoing heavy ion collisions program and, possibly, for cosmology . Even a partial restoration of $`U(1)_A`$ symmetry in the vicinity of the critical temperature of chiral symmetry breaking ($`T_c200`$ MeV) could dramatically change the mass and mixing pattern of the lightest neutral mesons ($`\pi ^0`$, $`\eta `$ and $`\eta ^{}`$), with signals including enhanced strangeness production or the more speculative possibility of forming parity violating disoriented $`\eta ^{}`$ condensates in heavy ion collisions .
Our aim in the present paper is rather modest: we will study the shift of the mass of the $`\eta `$ and $`\eta ^{}`$ mesons at low temperatures, in a regime in which the hadronic gas is mostly composed of pions. We will work in the framework of $`U(3)_L\times U(3)_R`$ chiral perturbation theory ($`\chi PT`$), in a simultaneous expansion in momenta, quark masses, number of colours $`N_c`$, and temperature $`T`$.
Our motivation for doing this investigation was threefold. First, the predictions of $`\chi PT`$ in a pion thermal bath, although limited in scope to $`T<\text{few}f_\pi `$, are essentially model independent (see for instance, the review of Smilga and references therein). Given the phenomenological success of the large $`N_c`$ expansion in vacuum, one might perhaps hope that the predictions of the present work are as robust. Next, we wanted to see to which extent the results derived in Ref. could be amended. As us, the authors have computed the shift of $`M_\eta `$ and $`M_\eta ^{}`$ at low temperatures using the Di Vecchia-Veneziano-Witten effective lagrangian (DVW) , but only to leading order $`\mathrm{\Delta }M_\eta ^{}^2T^2`$. However, it was not clear to us whether the leading order DVW lagrangian was a good approximation for this problem. Although the parameters of the lagrangian can be fitted to the observed mass and mixing pattern of the $`\eta `$ and $`\eta ^{}`$ mesons to within $`10\%`$ , the decay rates predicted for $`\eta ^{}\eta \pi \pi `$ are off the experimental values by a factor of about $`40`$. This issue, which is obviously relevant in order to determine the shift of $`M_\eta ^{}`$ in a pion bath, is however easily cured at next-to-leading order in the large $`N_c`$ expansion . As we will show, next-to-leading order corrections are also quite important at finite temperature, but not to the point of dramatically changing the conclusion of Ref. : at low temperatures, $`M_\eta ^{}`$ stays essentially constant. Finally, we wanted to see what the large $`N_c`$ expansion could teach us about the fate of the $`U(1)_A`$ symmetry at finite temperature. At zero temperature, in the confined phase of QCD, large $`N_c`$ arguments predict that $`M_\eta ^{}^21/N_c`$. On the other hand, at very high temperatures, in the quark-gluon plasma phase of QCD ($`T\mu _{\mathrm{hadr}}200`$ MeV), instanton calculus is reliable and predicts an effective restoration of $`U(1)_A`$ symmetry. Because of screening, instanton effects are suppressed at very large temperatures $`\mathrm{exp}(8\pi ^2/g(T)^2)`$. At large $`N_c`$, the suppression is more important, as $`1/g^2N_c/\lambda `$ with fixed ’t Hooft coupling $`\lambda =g^2N_c`$, and the exponential tends to vanish $`\mathrm{exp}(N_c/\lambda )0`$ as $`N_c`$ increases. Because of asymptotic freedom, $`\lambda `$ growths at lower temperatures and the instanton argument breaksdown. However, for $`N_c`$ large enough, a natural assumption is that the exponential suppression holds all the way down to the critical temperature of deconfinement $`T_c\mu _{\mathrm{hadr}}`$ . Although we have no proof of this statement, such a behaviour seems natural given the large release of entropy $`N_c^2`$ at $`T_c`$ and is actually known to occur in models in two dimensions . With this assumption, $`M_\eta ^{}`$ can be taken as an order parameter for $`U(1)_A`$ symmetry restoration at $`T_c`$. In Ref. , some information on the behaviour of $`M_\eta ^{}`$ near $`T_c`$ could be extracted assuming that the deconfining phase transition could be of second order at large $`N_c`$ . We will argue here that large $`N_c`$ favours a sharp drop of $`M_\eta ^{}`$ at $`T_c`$, consistent with first order transition to the phase with (approximate) $`U(1)_A`$ symmetry.
Our paper is organized as follows. In the next section, we briefly review $`U(3)_L\times U(3)_R`$ $`\chi PT`$, which extends the framework of the large $`N_c`$ DVW effective lagrangian beyond leading order. For definitiveness, we refer to the recent analysis of Herrera-Siklódy et al. . We then discuss the implications of these corrections at low temperature, in presence of a pion thermal bath. Most formulas are relegated to the appendix. In the last section, we speculate on the effect of higher order corrections in the large $`N_c`$ expansion and draw the conclusions.
## 2 Sketch of $`U(3)_L\times U(3)_R`$ chiral perturbation theory
At low energies and temperatures, the dynamics of QCD is governed by an approximate $`SU(3)_L\times SU(3)_R`$ chiral symmetry which is spontaneously broken to the diagonal $`SU(3)`$ in vacuum. If the mass of the up, down and strange quarks were vanishing, the symmetry would become exact and there would be eight massless Goldstone bosons. Phenomenological lagrangians, which treat the mass of the quarks as small perturbations, provide a powerful framework, known as Chiral Perturbation Theory ($`\chi PT`$), to study the properties of the light $`\pi `$, $`K`$, and $`\eta `$ mesons . The $`\eta ^{}`$ meson doesn’t a priori fit in this frame. It is substantially heavier than the other eight light mesons, and, in vacuum, would stay so even in the chiral limit of zero quark masses, because it receives most of its mass from the $`U(1)_A`$ anomaly through non-perturbative instanton-like effects . The effect of the axial anomaly can however be conveniently turned off by going to the limit of large number of colours $`N_c`$ . At infinite $`N_c`$ and in the chiral limit, the global symmetry becomes $`U(3)_L\times U(3)_R`$, spontaneously broken in vacuum to $`U(3)`$, with nine massless Goldstone bosons. Like chiral symmetry breaking effects by finite quark masses, $`1/N_c`$ suppressed contributions can be systematically introduced as perturbations in an effective lagrangian, an approach which has been quite fruitful . A systematic analysis of next-to-leading corrections, including $`𝒪(p^4)`$ operators, has been initiated in the recent . We refer to these latter works for more details and follow their conventions for ease of reference. We will work in Euclidean spacetime with metric $`g_{\mu \nu }\delta _{\mu \nu }`$ and use the imaginary time formalism to compute the thermal corrections.
### 2.1 Leading order
The leading order effective lagrangian is well-known . In the notation of it is written as
$`_{\mathrm{LO}}`$ $`=`$ $`{\displaystyle \frac{f^2}{4}}\left(v_{02}X^2+_\mu U^{}_\mu UU^{}\chi +\chi ^{}U\right),`$ (1)
where $`U`$ is the $`U(3)`$ matrix
$$Ue^{i\sqrt{2}\mathrm{\Phi }/f},$$
(2)
with $`\mathrm{\Phi }`$ the pseudoscalar meson matrix
$$\mathrm{\Phi }=\left(\begin{array}{ccc}\frac{\pi ^0}{\sqrt{2}}+\frac{\eta _8}{\sqrt{6}}+\frac{\eta _0}{\sqrt{3}}& \pi ^+& K^+\\ \pi ^{}& \frac{\pi ^0}{\sqrt{2}}+\frac{\eta _8}{\sqrt{6}}+\frac{\eta _0}{\sqrt{3}}& K^0\\ K^{}& \overline{K}^0& \frac{2\eta _8}{\sqrt{6}}+\frac{\eta _0}{\sqrt{3}}\end{array}\right),$$
(3)
and $`f=f_\pi =92.4`$ MeV at leading order. The mass matrix is
$$\chi =2B\text{diag}(m_u,m_d,m_s),$$
(4)
but we shall neglect isospin breaking effects ($`m_u=m_dm`$). The constant $`B`$ is related to the value of the $`\overline{q}q`$ condensate, $`M_\pi ^2=2mB2m\overline{q}q/f_\pi ^2`$ at leading order. The combination
$$X(x)\mathrm{log}U(x)+i\theta _{\mathrm{QCD}}=i\frac{\sqrt{6}}{f}\eta _0+i\theta _{\mathrm{QCD}},$$
(5)
is invariant under $`U(3)_L\times U(3)_R`$ transformations, $`\mathrm{log}U\mathrm{log}(g_RUg_L^{})+2i\alpha `$ and $`\theta _{\mathrm{QCD}}\theta _{\mathrm{QCD}}2\alpha `$. Because of this, any arbitrary function of $`X`$ can a priori enter in the construction of the effective lagrangian, with thus little predictive power. This is where the large $`N_c`$ expansion comes to the rescue by limiting the number of operators that can contribute at each level of approximation. In the chiral limit ($`m,m_s0`$), Eq. (1) gives
$$M_\eta ^{}^2=3v_{02},$$
(6)
which is the celebrated Veneziano-Witten relation for three massless flavours , with $`v_{02}2\tau /f^21/N_c`$, where $`\tau `$ is the topological susceptibility of pure Yang-Mills theory. The rationale of $`U(3)_L\times U(3)_R`$ $`\chi PT`$ is to count powers of $`p^2`$, $`m_q`$, and $`1/N_c`$ on the same level $`𝒪(\delta )`$ :
$$𝒪(\delta )p^2m_q1/N_c.$$
(7)
According to this counting rule, the leading order lagrangian (1) is $`𝒪(\delta ^0)`$ because $`f^2𝒪(N_c)`$ <sup>1</sup><sup>1</sup>1Note that the field expansion of $`U`$ brings further powers of $`1/f1/\sqrt{N}_c`$. The $`𝒪(\delta )`$ counting is to be understood to hold at the operator level..
At leading order, there are four unknown parameters in the lagrangian: $`f`$, $`v_{02}`$, and the combinations $`mB`$ and $`m_sB`$ (or $`xm_s/m1`$). On the other hand, we have at our disposal seven observables: $`f_\pi `$, $`f_K`$, the four masses of the light mesons, and the $`\eta `$$`\eta ^{}`$ mixing angle $`\theta `$. Using $`M_\eta ^{}`$ as input and the formulæ given for reference in the appendix, one obtains
$$\begin{array}{ccc}f^2\hfill & =& f_\pi ^2=f_K^2,\hfill \\ 2mB\hfill & =& M_\pi ^2,\hfill \\ x\hfill & & 24.1,\hfill \\ v_{02}\hfill & & 0.22\mathrm{GeV}^2,\hfill \end{array}$$
(8)
which predict that $`\theta 20^{}`$ and
$$M_\eta 494.4\mathrm{MeV}.$$
(9)
Remarkably, the latter number is only 10% off the experimental value $`M_\eta =547.3`$ MeV. It is however known that adjusting the parameters cannot improve the prediction because the ratio $`M_\eta ^2/M_\eta ^{}^2`$ has an upper bound<sup>2</sup><sup>2</sup>2Assuming $`M_\pi =0`$ to simplify, Eq. (1) gives
$$\frac{M_\eta ^2}{M_\eta ^{}^2}=\frac{3y\sqrt{9+2y+y^2}}{3y+\sqrt{9+2y+y^2}},$$
(10) where $`y9v_{02}/2(M_K^2M_\pi ^2)`$. This ratio reaches a maximum at $`y=3`$ (note that $`v_{02}<0`$) corresponding to
$$\frac{M_\eta }{M_\eta ^{}}<0.518,$$
(11) to be compared with the measured ratio $`M_\eta /M_\eta ^{}0.571`$. Taking into account $`M_\pi 0`$ improves things, but not enough. . One has to take into account next-to-leading order corrections to reach agreement .
At leading order, the only coupling between $`\eta ^{}`$ and the pions is from the quark mass term in the lagrangian (1) and is thus chirally suppressed. The amplitude for $`\eta ^{}\eta \pi \pi `$ is then
$$𝒜=\frac{M_\pi ^2}{6f_\pi }\left(2\sqrt{2}\mathrm{cos}(2\theta )\mathrm{sin}(2\theta )\right).$$
(12)
The corresponding decay rates
$$\begin{array}{ccc}\mathrm{\Gamma }(\eta ^{}\eta \pi ^0\pi ^0)& =& 1.0\mathrm{keV},\hfill \\ \mathrm{\Gamma }(\eta ^{}\eta \pi ^+\pi ^{})& =& 1.9\mathrm{keV}2\times \mathrm{\Gamma }(\eta ^{}\eta \pi ^0\pi ^0),\hfill \end{array}$$
(13)
are however much smaller than the experimental ones,
$$\begin{array}{ccc}\mathrm{\Gamma }_{\mathrm{exp}}(\eta ^{}\eta \pi ^0\pi ^0)& =& 42.0\pm 4.2\mathrm{keV},\\ \mathrm{\Gamma }_{\mathrm{exp}}(\eta ^{}\eta \pi ^+\pi ^{})& =& 88.9\pm 7.6\mathrm{keV}.\end{array}$$
(14)
We will not speculate on the reasons for this well-known discrepancy (see Refs. for a more recent discussion), but simply note that within the present framework, this issue can also be resolved at next-to-leading order <sup>3</sup><sup>3</sup>3Note that the amplitude (12) is constant and vanishes in the limit $`m_u=m_d=0`$, for any $`m_s`$. However, general arguments (and a fit to experimental data) indicate that for $`m_u=m_d=0`$, the amplitude should behave like $`𝒜=\text{const}\times p_\pi ^{(1)}p_\pi ^{(2)}`$, where $`p_\pi ^{(1,2)}`$ are the momenta of the outgoing pions, and where the constant is vanishing as the strange quark mass goes to zero. As shown in Ref. , this behaviour can be easily accommodated by introducing higher-order terms, an approach that is systematized by the $`\delta `$ expansion . The smallness of the leading order contribution is then considered as a mere accident, related to the smallness of the ratio $`M_\pi ^2/M_K^2`$..
### 2.2 Next-to-leading order
In our case, at next-to-leading order, $`𝒪(\delta )`$, only a few more terms can be added to the lagrangian (1:
$$\begin{array}{ccc}\hfill _{NLO}& =& _{LO}\hfill \\ & +& \frac{f^2}{4}(v_{31}XU^{}\chi \chi ^{}U+v_{40}U^{}_\mu UU^{}_\mu U\hfill \\ & +& iv_{50}U^{}_\mu U_\mu \theta _{\mathrm{QCD}}+v_{60}_\mu \theta _{\mathrm{QCD}}_\mu \theta _{\mathrm{QCD}})\hfill \\ & & M_0O_0M_3O_3+L_5O_5L_8O_8,\hfill \end{array}$$
(15)
where the $`O_{0,3,5,8}`$ are $`𝒪(p^4)`$ operators whose coupling constants are $`𝒪(N_c)`$:
$$\begin{array}{ccc}\hfill O_0& =& _\mu U_\nu U^{}_\mu U_\nu U^{},\hfill \\ \hfill O_3& =& _\mu U^{}_\mu U_\nu U^{}_\nu U,\hfill \\ \hfill O_5& =& _\mu U^{}_\mu U(U^{}\chi +\chi ^{}U),\hfill \\ \hfill O_8& =& \chi ^{}U\chi ^{}U+U^{}\chi U^{}\chi .\hfill \end{array}$$
(16)
The couplings $`v_{40}`$, $`v_{50}`$, and $`v_{60}`$ are not independent and either one of them can be set to zero by an appropriate change of variables, $`\eta _0/f\eta _0/f+\kappa \theta _{\mathrm{QCD}}`$. We shall choose $`v_{40}=0`$. Moreover, $`v_{50}`$ and $`v_{60}`$ will not appear in our calculations and can be discarded. At $`𝒪(\delta )`$, the only coupling related to the breaking of $`U(1)_A`$ symmetry is thus $`v_{31}𝒪(1/N_c)`$. Note that the corresponding operator is also chirally suppressed $`m_q`$.
At next-to-leading order, seven unknown parameters enter in the definition of the meson mass matrix: $`f`$, $`v_{02}`$, $`v_{31}`$, $`L_{5,8}`$, together with the quark masses $`m`$ and $`m_s`$ (see appendix for details). These can be expressed in terms of seven independent observables: $`f_\pi `$, $`f_K`$, $`M_\pi `$, $`M_K`$, $`M_\eta `$, $`M_\eta ^{}`$, and the $`\eta `$$`\eta ^{}`$ mixing angle $`\theta `$ . At this level, large $`N_c`$ $`\chi PT`$ is thus not predictive. The strategy adopted in Ref. was to impose that $`𝒪(\delta )`$ corrections are not too large so that the large $`N_c`$ expansion makes sense. For mixing angle in the range $`20^{}<\theta <24^{}`$, the fit gives
$$\begin{array}{ccc}\hfill 0.980& 2mB/M_\pi ^2& 0.988,\hfill \\ \hfill 18.3& x& 20.9,\hfill \\ \hfill 0.214\mathrm{GeV}^2& |v_{02}|& 0.239\mathrm{GeV}^2,\hfill \\ \hfill \mathrm{1.35\; 10}^3& L_8& \mathrm{1.57\; 10}^3,\hfill \\ \hfill 0.164& v_{31}& 0.161.\hfill \end{array}$$
(17)
together with $`f=90.8`$ MeV and $`L_5=\mathrm{2.0\; 10}^3`$ which are fixed by $`f_\pi `$ and $`f_K`$. Note that if $`v_{02}`$ does only change by about 10%, the shift in $`m_s`$ is quite large, $``$ 20–25.
Because they have four derivatives, the operators $`O_0`$ and $`O_3`$ do not contribute to the meson mass matrix in vacuum. However, they give the dominant contributions to the decay $`\eta ^{}\eta \pi \pi `$ . This is essentially because the extra derivatives introduce large amplification factors, $`(M_\eta ^{}/M_\pi )^2`$, with respect to the leading order amplitude<sup>4</sup><sup>4</sup>4This may actually casts some doubts on the validity of $`\chi PT`$ for such processes as one could expect higher-order effects to give non-negligible contributions to the decay $`\eta ^{}\eta \pi \pi `$. One may nevertheless hope that the large $`N_c`$ expansion is still reliable and that these corrections are $`1/N_c`$ suppressed. Whether this is true is unfortunately hard to check as we would evidently have too few hadronic data to completely fit the parameters of the effective lagrangian at higher orders in the $`\delta `$ expansion. Of course, this is precisely why the large $`N_c`$ expansion is invoked in the first place.. The observed decay rates are well reproduced with
$$\begin{array}{ccc}\hfill M_0& & \mathrm{1.2\; 10}^3,\hfill \\ \hfill M_3& & \mathrm{0.4\; 10}^3,\hfill \end{array}$$
(18)
values which can be independently inferred from the known $`L_1`$, $`L_2`$, and $`L_3`$ of $`SU(3)\times SU(3)`$ $`\chi PT`$ (in the nomenclature of Gasser and Leutwyler )<sup>5</sup><sup>5</sup>5According to Ref. , $`M_0=\frac{2}{3}(L_1+L_2)+𝒪(N_c^0)`$ and $`M_3=L_3+2M_0`$..
Thus all the parameters of the next-to-leading order effective lagrangian are fixed by low-energy phenomenology.
## 3 $`M_\eta ^{}`$ in a pion thermal bath
In Ref. , the leading order lagrangian (1) has been used to study the shift of $`M_\eta `$ and $`M_\eta ^{}`$ at one-loop in a pion thermal bath. The effect they found is very tiny, as at $`T200`$ MeV the relative mass shifts are only $`𝒪(0.1\%)`$. The reason for this is easy to understand. The $`\eta `$ and $`\eta ^{}`$ mesons receive most of their mass from the topological susceptibility term $`v_{02}`$ and/or from the strange quark mass, while the pion thermal corrections only modifies the tiny contribution from the pion mass term $`M_\pi ^2`$. Thermal kaons could give a larger effect, $`M_K^2`$, but the density of these is exponentially suppressed at low temperatures, $`\mathrm{exp}(M_K/T)`$. One might wonder whether next-to-leading order corrections could directly affect the contribution of the leading order $`U(1)_A`$ breaking term $`v_{02}`$. As we have seen in the previous section, five extra operators appear at next-to-leading order in the large $`N_c`$ expansion and, of these, only the one with coupling $`v_{31}`$ is related to $`U(1)_A`$ symmetry breaking. Unfortunately, this term is also chirally suppressed, $`m`$, and its contribution is only $`𝒪(v_{31}M_\pi ^2T^2/f_\pi ^2)`$. At temperatures of interest, this is small compared to $`v_{02}`$, but of the same magnitude as the leading order thermal correction. The other four operators will also contribute, but in a less interesting way, as they are invariant under $`U(1)_A`$. Furthermore, their effects are also $`M_\pi ^2`$.
We have computed the shift of the mass of $`\eta `$ and $`\eta ^{}`$ at one-loop, at next-to-leading order in the expansion in $`\delta `$. We have not taken into account two-loop corrections from the leading order lagrangian. Although it is not clear whether this is legitimate numerically speaking, neglecting these is however consistent with the rules of large $`N_c`$ chiral perturbation theory. Indeed, the natural extension of $`\delta `$ power-counting to finite temperature is
$$𝒪(\delta )p^2m_q1/N_cT^2.$$
(19)
At leading order, $`M_\eta ^{}^2=𝒪(\delta )1/N_c`$ and the one-loop thermal correction is $`M_\pi ^2T^2/f^2\delta ^3`$. At two-loop, using the leading order lagrangian, the shift is $`M_\pi ^2T^4/f^4\delta ^5`$, while at one-loop using the next-to-leading order lagrangian, the shift is typically $`v_{31}M_\pi ^2T^2/f^2\delta ^4`$ (using $`v_{31}1/N_c`$) and thus dominant. Consistency thus requires to neglect the two-loop contributions. This greatly simplifies the calculations which are a bit cumbersome, but otherwise straightforward.
The relevant diagrams are those of Fig. 1, where the loops contain only pions. At next-to-leading order there are two related thermal loops:
$$I_1(T)=T\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{d^3\stackrel{}{k}}{(2\pi )^3}\frac{1}{K^2+M_\pi ^2},$$
(20)
with $`K^2=k_0^2+\stackrel{}{k}^2`$, where $`k_0=2\pi nT`$, with $`n`$ integer, are the Matsubara frequencies and
$$I_2(T)=T\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{d^3\stackrel{}{k}}{(2\pi )^3}\frac{K^2}{K^2+M_\pi ^2},$$
(21)
with $`I_2(T)=M_\pi ^2I_1(T)`$. As usual, we drop the ultraviolet divergent part of the pion loops as these can in principle be reabsorbed in vacuum parameters, including next-to-next-to-leading order counter-terms. The sum over $`n`$ can then be readily evaluated using standard techniques ,
$$I_1(T)=\frac{d^3\stackrel{}{k}}{(2\pi )^3}\frac{1}{\omega }\frac{1}{\mathrm{exp}(\omega /T)1}=\frac{M_\pi T}{2\pi ^2}\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n}K_1\left(\frac{nM_\pi }{T}\right),$$
(22)
where $`\omega ^2=\stackrel{}{k}^2+M_\pi ^2`$. For instance, for $`T>M_\pi `$
$$I_1(T)\frac{T^2}{12}.$$
(23)
In the sequel we simply compute (22) numerically. As the relevant formula are not particularly transparent, we have relegated them to the appendix. Fig. 2 and Fig. 3 show the shift of $`M_\eta ^{}`$ and $`M_\eta `$ at low temperature both at leading and next-to-leading order. To be definite we have chosen the set of parameters corresponding to $`\theta 20^{}`$. The net thermal effects are not dramatic: both masses decrease, but only slightly. As expected, the shift of the mass $`M_\eta ^{}`$ is more pronounced at next-to-leading order, but the effect is not very significant. Again, this is because, both at leading and next-to-leading orders, the thermal corrections are chirally suppressed, $`M_\pi ^2T^2/f_\pi ^2`$. For completeness, we have also plotted in Fig. 4 the shift of the mixing angle at low temperature. As both the $`\eta `$ and $`\eta ^{}`$ masses diminish, the angle is not very much affected. It decreases a bit (toward ideal mixing?), consistent with the relatively larger shift of $`M_\eta ^{}`$.
## 4 Lessons from large $`N_c`$?
As we have seen in the previous section, the mass of $`\eta ^{}`$ is almost not affected at low temperatures in a pion bath. This is because, at this order, the pion thermal corrections are chirally suppressed, smaller than $`M_\pi ^20.02`$ GeV<sup>2</sup>, and thus essentially negligible compared to the contribution from the $`U(1)_A`$ symmetry breaking term $`v_{02}0.22`$ GeV<sup>2</sup>. In particular, in the chiral limit, $`M_\pi =0`$, all the corrections vanish and $`M_\eta ^{}`$ is temperature independent up to next-to-leading order in $`\chi PT`$. In the chiral limit, the leading contribution from pions to the shift of $`M_\eta ^{}`$ presumably arise from $`𝒪(p^4)`$ operators like<sup>6</sup><sup>6</sup>6The $`O(p^2)`$ operator
$$f^2\frac{X^2}{N_c^2}_\mu U^{}_\mu U,$$
(24) can contribute at one-loop if and only if $`M_\pi ^20`$. It could contribute at two-loop order in the chiral limit $`T^4`$, but doesn’t because pion interactions are too soft. This is a well-known feature of pion thermal corrections which is for instance manifest in the absence of $`T^6`$ terms in the free energy of a pion gas in the chiral limit , or in the fact that massless thermal pions move at the speed of light to order $`T^2`$ .
$$\frac{1}{N_c}X^2_\mu U_\nu U^{}_\mu U_\nu U^{}.$$
(25)
The coupling is $`𝒪(1/N_c)`$ because there is a factor of $`1/N_c^2`$ coming with $`X^2`$ and one of $`N_c`$ from the coupling of the $`𝒪(p^4)`$ operator. The best way to see this is to replace the coupling $`M_0N_c`$ (or $`M_3`$) from the next-to-leading order lagrangian by a function of $`X`$, $`M_0M_0(X)`$ and expand to second order in $`X`$, which brings down a factor of $`1/N_c^2`$. Because there are at least four pions in the expansion of the operator (25), the leading pion thermal correction to $`M_\eta ^{}`$ in the chiral limit is a two-loop effect,
$$\delta M_\eta ^{}^2(T)\frac{1}{N_c}\frac{T^8}{f^6}\frac{1}{N_c^4}\frac{T^8}{\mu _{\mathrm{hadr}}^6}.$$
(26)
We have made the large $`N_c`$ dependence of the pion decay constant $`f^2`$ manifest by defining $`f^2N_c\mu _{\mathrm{hadr}}`$. Of course, the sign of the correction is not known and $`M_\eta ^{}`$ could go up or down. Also, if we compare with $`M_\eta ^{}^2(0)\mu _{\mathrm{hadr}}^2/N_c`$, we infer that $`M_\eta ^{}`$ is quasi-constant for temperatures
$$TT_{}N_c^{3/8}\mu _{\mathrm{hadr}}.$$
(27)
In the large $`N_c`$ framework, the natural scale for deconfinement is $`T_c\mu _{\mathrm{hadr}}`$, which is also the temperature at which the pions from the hadronic gas overlap. It is natural to assume that chiral symmetry restoration takes place at the same temperature, driven by the release of $`𝒪(N_c^2)`$ gluon degrees of freedom . The estimate in Eq. (27) then seems to imply that $`M_\eta ^{}`$ is essentially constant up to the temperature of deconfinement, since $`T_{}T_c\mu _{\mathrm{hadr}}`$ for $`N_c`$ large. This conclusion is however premature because the low momentum expansion breaks down near $`T_c`$ and we must take into account the contribution of operators with arbitrary number of derivatives. We claim that the dominant operators at large $`N_c`$ are of the form
$$\frac{1}{N_c\mu _{\mathrm{hadr}}^{2k4}}X^2_{\mu _1}U_{\mu _2}U^{}\mathrm{}_{\mu _k}U_{\mu _1}U^{}_{\mu _2}U\mathrm{}_{\mu _k}U^{}.$$
(28)
These operators are irrelevant at low energies but become marginal for $`T_c`$. A six-derivative operator, for instance, first contribute at three loops $`\delta M_\eta ^{}^2T^{12}/(N_c^5\mu _{\mathrm{hadr}}^{10})`$. For comparison, the contribution of a three-loop diagram with a four-derivatives (NLO) and a two-derivatives vertices (LO) is $`1/N_cT^{10}/f^81/N_c^5T^{10}/\mu _{\mathrm{hadr}}^8`$ and is subdominant at large $`N_c`$. For generic $`k`$, the operators of Eq. (28) give $`\delta M_\eta ^{}^2T^{4k}/(N_c^{k+2}\mu _{\mathrm{hadr}}^{4k2})`$. The ratio of two consecutive terms $`k`$ and $`k+1`$ becomes $`𝒪(1)`$ at $`T^{}N_c^{1/4}\mu _{\mathrm{hadr}}`$, independent of $`k`$. At large $`N_c`$, $`T_cT^{}T^{}`$ and the perturbative expansion still breaks down above the temperature of deconfinement.
Another set of operators could be relevant at large $`N_c`$ because the $`\eta ^{}`$ is then rather light, $`M_\eta ^{}\mu _{\mathrm{hadr}}/N_c^{1/2}<T_c`$. Thus one should include operators that involve arbitrary powers of the $`\eta ^{}`$ field, like
$$\begin{array}{ccc}\hfill & & f^2F\left(\frac{X}{N_c}\right)_\mu U^{}_\mu U\hfill \\ & =& \left(\frac{\eta _0^2}{\mu ^2N_c^2}+\frac{\eta _0^4}{\mu ^4N_c^5}+\mathrm{}\right)(_\mu \eta ^0)^2+\text{pion terms},\hfill \end{array}$$
(29)
which contributes to the wave-function renormalization of $`\eta ^{}`$,
$$\delta Z_\eta ^{}T^2/(N_c^2\mu _{\mathrm{hadr}}^2)+T^4/(N_c^5\mu _{\mathrm{hadr}}^4)+\mathrm{},$$
(30)
or terms of the form
$$N_c^2\mu _{\mathrm{hadr}}^4G\left(\frac{X}{N_c}\right)\frac{\mu _{\mathrm{hadr}}^2}{N_c}\eta _0^2+\mathrm{\#}_1\frac{1}{N_c^4}\eta _0^4+\mathrm{\#}_2\frac{1}{\mu _{\mathrm{hadr}}^2N_c^7}\eta _0^6+\mathrm{}.$$
(31)
However, a common feature of these operators is that they are very suppressed at large $`N_c`$. They become important only for $`TN_c^{3/2}\mu _{\mathrm{hadr}}`$, much higher than $`T^{}`$ so that their contribution is subleading compared to operators like in Eq. (28).
Can we conclude anything from these considerations? In all the cases discussed above, the leading thermal corrections to $`M_\eta ^{}`$, in the chiral limit and for $`N_c`$ large, become important for temperature which are higher than the critical temperature of deconfinement $`T_c\mu _{\mathrm{hadr}}`$ by a factor of $`N_c^\gamma `$. Although the value of $`\gamma `$ is hard to guess, as various corrections can get mixed up, we believe it is reasonable to conjecture that $`\gamma `$ is strictly positive. This implies that just below $`T_c`$, $`M_\eta ^{}(T)=M_\eta ^{}(0)`$ to a very good approximation. The standard lore is that the deconfining phase transition at $`T_c`$ is of first order for $`N_c`$ large <sup>7</sup><sup>7</sup>7Various arguments, including recent developments in string theory (see Sect. 6.2.2 in Ref ) and lattice simulations of $`N_c=4`$ pure Yang-Mills theory , favour a first order deconfining phase transition. A case for a second order phase transition has been made in Ref. , in light of the structure of the Columbia diagram.. Because the temperature at which hadronic interactions can affect $`M_\eta ^{}`$ is (very much) larger than the temperature of deconfinement, we expect that changes in $`M_\eta ^{}`$ will be instead triggered by the release of the large number of gluons and will thus drop discontinuously at $`T_c`$, i.e. that there is a first order transition to a phase with (approximate) $`U(1)_A`$ symmetry.
This behaviour is not inconsistent with various other expectations. For three light quark flavours, $`N_f3`$, the transition to the chirally symmetric phase is probably first order while for $`N_f=2`$, the phase transition is supposed to be of second order, in the universality class of $`O(4)`$ . It has been argued by Smilga that the latter behaviour is not inconsistent with a first order deconfining phase transition at large $`N_c`$ . The reason is that, unlike for $`M_\eta ^{}`$, there is an infinite subset of thermal corrections that contribute to the same order in $`N_c`$ to the shift of the quark condensate $`\mathrm{\Sigma }\overline{q}q`$,
$$\mathrm{\Sigma }(T)=\mathrm{\Sigma }\left(1\mathrm{\#}\frac{T^2}{N_c^2\mu _{\mathrm{hadr}}^2}F\left(\frac{T}{\mu _{\mathrm{hadr}}}\right)\right).$$
(32)
Even though thermal corrections are suppressed like $`1/N_c^2`$, the (unknown) function $`F(x)`$ may be singular near, but below $`T_c\mu _{\mathrm{hadr}}`$. If, for instance, $`F`$ has a simple pole at $`T=T_0<T_c`$, $`F\mu _{\mathrm{hadr}}^2/(T^2T_0^2)`$ and the chiral phase transition is second order with a critical region near $`T_c`$, that is of order
$$\frac{\mathrm{\Delta }T}{T_c}\frac{1}{N_c^2}.$$
(33)
Alternatively, if $`U(1)_A`$ symmetry is effectively restored at $`T_c`$, the large $`N_c`$ behaviour (32) is also consistent with a fluctuation induced first order phase transition. For $`N_f=1`$ finally, chiral symmetry is broken by the anomaly at all temperature and there is no chiral phase transition. However, if instanton transitions are strongly suppressed just above $`T_c`$, chiral symmetry can be effectively restored and the phase transition is presumably first order.
## 5 Conclusions
We have studied the behaviour of the mass of the $`\eta ^{}`$ pseudoscalar meson at finite temperature using constraints from chiral symmetry and large $`N_c`$ power counting. The main conclusion to be drawn from this work is that $`M_\eta ^{}`$ is essentially unchanged at low temperatures. A tentative analysis of the effect of leading higher order corrections at large $`N_c`$ suggests that $`M_\eta ^{}`$ changes discontinuously at the temperature of deconfinement. The implications of these considerations for the real world, i.e. $`N_c=3`$, are not quite clear as we would expect the suppression of instanton effects only at asymptotically high temperatures. It is however striking that recent lattice simulations, with $`N_f=2`$ staggered and domain wall fermions, both show a strong suppression of $`U(1)_A`$ breaking effects at low temperatures $`T1.2T_c`$. Because this temperature is outside the critical region, the order of the chiral phase transition is probably not affected. It could be of interest to consider doing simulations with $`N_c>3`$, although this would probably be time consuming, or maybe with one flavour and various $`N_c`$. Consider for example a plot of the $`\pi `$ and $`\delta `$ susceptibilities near the critical temperature as computed on the lattice . Large $`N_c`$ arguments suggest that the curves of the susceptibilities would be flatter below $`T_c`$ —because the confined phase is colder— and that splitting between $`\pi `$ and $`\delta `$ (a.k.a. $`a_0`$) should be narrower above $`T_c`$ —because $`U(1)_A`$ breaking is more suppressed, maybe like in Fig. 5—. The Columbia diagram could change accordingly: the critical line around the region of small $`M_u=M_d`$ masses would move as in Fig. 6.
## Acknowledgements
We would like to thank P. Herrera-Siklódy and F. Karsch for useful discussions and S. Peris and J. Taron for a careful reading of the manuscript. Work partly supported by the EEC, TMR-CT98-0169, EURODAPHNE network.
## Appendix
### 5.1 Useful formulæ
#### 5.1.1 Leading order
$`2mB`$ $`=`$ $`M_\pi ^2,`$ (34)
$`x`$ $`=`$ $`2{\displaystyle \frac{M_K^2}{M_\pi ^2}}2,`$ (35)
$`f`$ $`=`$ $`f_\pi ,`$ (36)
$`3v_{02}`$ $`=`$ $`M_\eta ^2{\displaystyle \frac{2M_K^2+M_\pi ^2}{3}}+{\displaystyle \frac{2\sqrt{2}}{3}}(M_K^2M_\pi ^2)\mathrm{tan}\theta .`$ (37)
#### 5.1.2 Next-to-leading order
##### Some definitions:
$`\mathrm{\Delta }_M`$, $`\mathrm{\Delta }_N`$ are defined as
$`\mathrm{\Delta }_M`$ $``$ $`{\displaystyle \frac{8}{f^2}}(M_K^2M_\pi ^2)(2L_8L_5),`$ (38)
$`\mathrm{\Delta }_N`$ $``$ $`3v_{31}12{\displaystyle \frac{L_5}{f^2}}v_{02}.`$ (39)
##### Next-to-leading order parameters:
The next-to-leading order parameters can be expressed in terms of observables through
$`2mB`$ $`=`$ $`M_\pi ^2\left(1{\displaystyle \frac{M_\pi ^2}{M_K^2M_\pi ^2}}\mathrm{\Delta }_M\right),`$ (40)
$`x`$ $`=`$ $`2{\displaystyle \frac{M_K^2}{M_\pi ^2}}(1\mathrm{\Delta }_M)2,`$ (41)
$`f`$ $`=`$ $`f_\pi \left(14{\displaystyle \frac{L_5}{f^2}}M_\pi ^2\right),`$ (42)
$`3v_{02}`$ $`=`$ $`M_\eta ^2{\displaystyle \frac{2M_K^2+M_\pi ^2}{3}}`$ (43)
$`+`$ $`{\displaystyle \frac{2\sqrt{2}}{3}}(M_K^2M_\pi ^2)(1+\mathrm{\Delta }_M\mathrm{\Delta }_N)\mathrm{tan}\theta `$
$``$ $`{\displaystyle \frac{2}{3}}\left[(M_K^2M_\pi ^2)\mathrm{\Delta }_M(2M_K^2+M_\pi ^2)\mathrm{\Delta }_N\right],`$
$`{\displaystyle \frac{L_5}{f^2}}`$ $`=`$ $`{\displaystyle \frac{1}{4(M_K^2M_\pi ^2)}}\mathrm{\Delta }_P,`$ (44)
where
$`\mathrm{\Delta }_M`$ $`=`$ $`{\displaystyle \frac{M_\pi ^2+3M_\eta ^24M_K^2+3(M_\eta ^2M_\eta ^2)\mathrm{sin}^2\theta }{4(M_K^2M_\pi ^2)}},`$ (45)
$`\mathrm{\Delta }_N`$ $`=`$ $`1+{\displaystyle \frac{3}{4\sqrt{2}}}{\displaystyle \frac{(M_\eta ^2M_\eta ^2)\mathrm{sin}2\theta }{M_K^2M_\pi ^2}}+\mathrm{\Delta }_M,`$ (46)
$`\mathrm{\Delta }_P`$ $`=`$ $`{\displaystyle \frac{f_K}{f_\pi }}1.`$ (47)
### 5.2 Results at $`T=0`$
#### 5.2.1 Leading order
##### Mass matrix:
$`m_{88}^2`$ $`=`$ $`{\displaystyle \frac{1}{3}}(4M_K^2M_\pi ^2),`$ (48)
$`m_{80}^2`$ $`=`$ $`{\displaystyle \frac{2\sqrt{2}}{3}}(M_K^2M_\pi ^2),`$ (49)
$`m_{00}^2`$ $`=`$ $`{\displaystyle \frac{1}{3}}(2M_K^2+M_\pi ^2)3v_{02}`$ (50)
$`=`$ $`{\displaystyle \frac{1}{3}}(2M_K^2+M_\pi ^2){\displaystyle \frac{2}{3}}y(M_K^2M_\pi ^2),`$
where
$$y\frac{9v_{02}}{2(M_K^2M_\pi ^2)}.$$
(51)
##### Mixing angle:
$$\mathrm{tan}2\theta \frac{2m_{80}^2}{m_{00}^2m_{88}^2}=\frac{2\sqrt{2}}{1+y}.$$
(52)
##### Physical masses:
$`M_\eta ^2`$ $`=`$ $`M_K^2{\displaystyle \frac{M_K^2M_\pi ^2}{3}}(y+\sqrt{9+2y+y^2}),`$ (53)
$`M_\eta ^2`$ $`=`$ $`M_K^2{\displaystyle \frac{M_K^2M_\pi ^2}{3}}(y\sqrt{9+2y+y^2}),`$ (54)
with
$$M_\eta ^2+M_\eta ^2m_{88}^2+m_{00}^2=2M_K^2\frac{2}{3}y(M_K^2M_\pi ^2).$$
(55)
#### 5.2.2 Next-to-leading order
##### Mass matrix:
$`m_{88}^2`$ $`=`$ $`{\displaystyle \frac{1}{3}}(4M_K^2M_\pi ^2)+{\displaystyle \frac{4}{3}}(M_K^2M_\pi ^2)\mathrm{\Delta }_M,`$ (56)
$`m_{80}^2`$ $`=`$ $`{\displaystyle \frac{2\sqrt{2}}{3}}(M_K^2M_\pi ^2)(1+\mathrm{\Delta }_M\mathrm{\Delta }_N),`$ (57)
$`m_{00}^2`$ $`=`$ $`{\displaystyle \frac{1}{3}}(2M_K^2+M_\pi ^2)3v_{02}`$ (58)
$`+`$ $`{\displaystyle \frac{2}{3}}(M_K^2M_\pi ^2)\mathrm{\Delta }_M{\displaystyle \frac{2}{3}}(2M_K^2+M_\pi ^2)\mathrm{\Delta }_N`$
$`=`$ $`{\displaystyle \frac{1}{3}}(2M_K^2+M_\pi ^2)(12\mathrm{\Delta }_N)+{\displaystyle \frac{2}{3}}(M_K^2M_\pi ^2)\mathrm{\Delta }_M`$
$``$ $`{\displaystyle \frac{2}{3}}y(M_K^2M_\pi ^2).`$
##### Mixing angle:
$$\mathrm{tan}2\theta =\frac{2\sqrt{2}}{1+y}\left(1+\frac{y}{1+y}\mathrm{\Delta }_M\left(1+\frac{1}{1+y}\frac{2M_K^2+M_\pi ^2}{M_K^2M_\pi ^2}\right)\mathrm{\Delta }_N\right).$$
(59)
##### Physical masses:
$`M_\eta ^2`$ $`=`$ $`M_K^2{\displaystyle \frac{M_K^2M_\pi ^2}{3}}(y+\sqrt{9+2y+y^2})`$ (60)
$`+`$ $`\left(1{\displaystyle \frac{9+y}{3\sqrt{9+2y+y^2}}}\right)(M_K^2M_\pi ^2)\mathrm{\Delta }_M`$
$``$ $`{\displaystyle \frac{1}{3}}\left(2M_K^2+M_\pi ^2{\displaystyle \frac{3(2M_K^23M_\pi ^2)y(2M_K^2+M_\pi ^2)}{\sqrt{9+2y+y^2}}}\right)\mathrm{\Delta }_N,`$
$`M_\eta ^2`$ $`=`$ $`M_K^2{\displaystyle \frac{M_K^2M_\pi ^2}{3}}(y\sqrt{9+2y+y^2})`$ (61)
$`+`$ $`\left(1+{\displaystyle \frac{9+y}{3\sqrt{9+2y+y^2}}}\right)(M_K^2M_\pi ^2)\mathrm{\Delta }_M`$
$``$ $`{\displaystyle \frac{1}{3}}\left(2M_K^2+M_\pi ^2+{\displaystyle \frac{3(2M_K^23M_\pi ^2)y(2M_K^2+M_\pi ^2)}{\sqrt{9+2y+y^2}}}\right)\mathrm{\Delta }_N,`$
with
$`M_\eta ^2+M_\eta ^2`$ $`=`$ $`2M_K^2{\displaystyle \frac{2}{3}}y(M_K^2M_\pi ^2)`$ (62)
$`+`$ $`2(M_K^2M_\pi ^2)\mathrm{\Delta }_M{\displaystyle \frac{2}{3}}(2M_K^2+M_\pi ^2)\mathrm{\Delta }_N.`$
### 5.3 Results at $`T0`$
#### 5.3.1 Leading order
##### Mass matrix:
$`m_{88}^2(T)`$ $`=`$ $`m_{88}^2(0){\displaystyle \frac{M_\pi ^2}{2f_\pi ^2}}I(T),`$ (63)
$`m_{80}^2(T)`$ $`=`$ $`m_{80}^2(0){\displaystyle \frac{M_\pi ^2}{\sqrt{2}f_\pi ^2}}I(T),`$ (64)
$`m_{00}^2(T)`$ $`=`$ $`m_{00}^2(0){\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}I(T),`$ (65)
where
$`I(T)`$ $``$ $`{\displaystyle \frac{d^3\stackrel{}{k}}{(2\pi )^3}\frac{1}{\omega }\frac{1}{e^{\beta \omega }1}};\omega =\sqrt{\stackrel{}{k}^2+M_\pi ^2},\beta ={\displaystyle \frac{1}{T}}`$ (66)
$`=`$ $`{\displaystyle \frac{M_\pi T}{2\pi ^2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}K_1\left({\displaystyle \frac{nM_\pi }{T}}\right)\stackrel{M_\pi 0}{}{\displaystyle \frac{T^2}{12}}.`$
##### Mixing angle:
$$\mathrm{tan}2\theta (T)=\mathrm{tan}2\theta (0)+\frac{2\sqrt{2}}{1+y}\frac{y}{1+y}\frac{3}{4}\frac{M_\pi ^2}{M_K^2M_\pi ^2}\frac{1}{f_\pi ^2}I(T).$$
(67)
##### Physical masses:
$`M_\eta ^2(T)`$ $`=`$ $`M_\eta ^2(0){\displaystyle \frac{3}{4}}{\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}I(T)\left(1+{\displaystyle \frac{9+y}{3\sqrt{9+2y+y^2}}}\right),`$ (68)
$`M_\eta ^2(T)`$ $`=`$ $`M_\eta ^2(0){\displaystyle \frac{3}{4}}{\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}I(T)\left(1{\displaystyle \frac{9+y}{3\sqrt{9+2y+y^2}}}\right),`$ (69)
with
$$M_\eta ^2(T)+M_\eta ^2(T)=M_\eta ^2(0)+M_\eta ^2(0)\frac{3}{2}\frac{M_\pi ^2}{f_\pi ^2}I(T).$$
(70)
#### 5.3.2 Next-to-leading order
##### Mass matrix:
$`m_{88}^2(T)`$ $`=`$ $`m_{88}^2(0){\displaystyle \frac{M_\pi ^2}{2f_\pi ^2}}I(T)(1+{\displaystyle \frac{2M_\pi ^2}{M_K^2M_\pi ^2}}(\mathrm{\Delta }_P+{\displaystyle \frac{3}{2}}\mathrm{\Delta }_M)`$ (71)
$`+`$ $`24{\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}(M_0+M_3)),`$
$`m_{80}^2(T)`$ $`=`$ $`m_{80}^2(0){\displaystyle \frac{M_\pi ^2}{\sqrt{2}f_\pi ^2}}I(T)(1+{\displaystyle \frac{2M_\pi ^2}{M_K^2M_\pi ^2}}(\mathrm{\Delta }_P+{\displaystyle \frac{3}{2}}\mathrm{\Delta }_M)`$ (72)
$``$ $`\mathrm{\Delta }_N+24{\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}(M_0+M_3)(1{\displaystyle \frac{y}{3}}{\displaystyle \frac{M_K^2M_\pi ^2}{M_\pi ^2}})),`$
$`m_{00}^2(T)`$ $`=`$ $`m_{00}^2(0){\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}I(T)(1+{\displaystyle \frac{2M_\pi ^2}{M_K^2M_\pi ^2}}(\mathrm{\Delta }_P+{\displaystyle \frac{3}{2}}\mathrm{\Delta }_M)`$ (73)
$``$ $`2\mathrm{\Delta }_N+24{\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}(M_0+M_3)(1{\displaystyle \frac{2y}{3}}{\displaystyle \frac{M_K^2M_\pi ^2}{M_\pi ^2}})).`$
##### Mixing angle:
$`\mathrm{tan}2\theta (T)`$ $`=`$ $`\mathrm{tan}2\theta (0)+{\displaystyle \frac{2\sqrt{2}}{1+y}}{\displaystyle \frac{y}{1+y}}{\displaystyle \frac{3}{4}}{\displaystyle \frac{M_\pi ^2}{M_K^2M_\pi ^2}}{\displaystyle \frac{1}{f_\pi ^2}}I(T)\times `$
$`(1`$ $`+`$ $`{\displaystyle \frac{2M_\pi ^2}{M_K^2M_\pi ^2}}\left(\mathrm{\Delta }_P+{\displaystyle \frac{3}{2}}\mathrm{\Delta }_M\right){\displaystyle \frac{2}{1+y}}\mathrm{\Delta }_M`$ (74)
$`+`$ $`{\displaystyle \frac{1y}{1+y}}\mathrm{\Delta }_N+{\displaystyle \frac{3}{(1+y)y}}{\displaystyle \frac{2M_K^2(1+y)M_\pi ^2}{M_K^2M_\pi ^2}}\mathrm{\Delta }_N`$
$`+`$ $`24{\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}(M_0+M_3)\left({\displaystyle \frac{M_K^2\frac{y}{3}(M_K^2M_\pi ^2)}{M_\pi ^2}}\right)).`$
##### Physical masses:
$`M_\eta ^2(T)`$ $`=`$ $`M_\eta ^2(0){\displaystyle \frac{3}{4}}{\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}I(T)(1+{\displaystyle \frac{9+y}{3\sqrt{9+2y+y^2}}}`$ (75)
$`+`$ $`{\displaystyle \frac{2M_\pi ^2}{M_K^2M_\pi ^2}}((1+{\displaystyle \frac{9+y}{3\sqrt{9+2y+y^2}}})\mathrm{\Delta }_P`$
$`+{\displaystyle \frac{3}{2}}(1+{\displaystyle \frac{27(3+y)+(9+y)y^2}{3(9+2y+y^2)^{3/2}}}`$
$`+{\displaystyle \frac{\frac{2}{3}y^2}{3(9+2y+y^2)^{3/2}}}{\displaystyle \frac{4M_K^2M_\pi ^2}{M_\pi ^2}})\mathrm{\Delta }_M)`$
$``$ $`{\displaystyle \frac{4}{3}}(1+{\displaystyle \frac{(3+y)((9+y)+(3+y)y)}{(9+2y+y^2)^{3/2}}}`$
$`+{\displaystyle \frac{6y}{(9+2y+y^2)^{3/2}}}{\displaystyle \frac{M_\pi ^2}{M_K^2M_\pi ^2}})\mathrm{\Delta }_N`$
$`+`$ $`24{\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}(M_0+M_3)(1+{\displaystyle \frac{9+y}{3\sqrt{9+2y+y^2}}}`$
$`{\displaystyle \frac{4}{9}}y(1+{\displaystyle \frac{3+y}{\sqrt{9+2y+y^2}}}){\displaystyle \frac{M_K^2M_\pi ^2}{M_\pi ^2}})),`$
$`M_\eta ^2(T)`$ $`=`$ $`M_\eta ^2(0){\displaystyle \frac{3}{4}}{\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}I(T)(1{\displaystyle \frac{9+y}{3\sqrt{9+2y+y^2}}}`$ (76)
$`+`$ $`{\displaystyle \frac{2M_\pi ^2}{M_K^2M_\pi ^2}}((1{\displaystyle \frac{9+y}{3\sqrt{9+2y+y^2}}})\mathrm{\Delta }_P`$
$`+{\displaystyle \frac{3}{2}}(1{\displaystyle \frac{27(3+y)+(9+y)y^2}{3(9+2y+y^2)^{3/2}}}`$
$`{\displaystyle \frac{\frac{2}{3}y^2}{3(9+2y+y^2)^{3/2}}}{\displaystyle \frac{4M_K^2M_\pi ^2}{M_\pi ^2}})\mathrm{\Delta }_M)`$
$``$ $`{\displaystyle \frac{4}{3}}(1{\displaystyle \frac{(3+y)((9+y)+(3+y)y)}{(9+2y+y^2)^{3/2}}}`$
$`{\displaystyle \frac{6y}{(9+2y+y^2)^{3/2}}}{\displaystyle \frac{M_\pi ^2}{M_K^2M_\pi ^2}})\mathrm{\Delta }_N`$
$`+`$ $`24{\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}(M_0+M_3)(1{\displaystyle \frac{9+y}{3\sqrt{9+2y+y^2}}}`$
$`{\displaystyle \frac{4}{9}}y(1{\displaystyle \frac{3+y}{\sqrt{9+2y+y^2}}}){\displaystyle \frac{M_K^2M_\pi ^2}{M_\pi ^2}})),`$
with
$`M_\eta ^2(T)+M_\eta ^2(T)`$ $`=`$ $`M_\eta ^2(0)+M_\eta ^2(0){\displaystyle \frac{3}{2}}{\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}I(T)\times `$
$`(1`$ $`+`$ $`{\displaystyle \frac{2M_\pi ^2}{M_K^2M_\pi ^2}}\left(\mathrm{\Delta }_P+{\displaystyle \frac{3}{2}}\mathrm{\Delta }_M\right){\displaystyle \frac{4}{3}}\mathrm{\Delta }_N`$ (77)
$`+`$ $`24{\displaystyle \frac{M_\pi ^2}{f_\pi ^2}}(M_0+M_3)(1{\displaystyle \frac{4}{9}}y{\displaystyle \frac{M_K^2M_\pi ^2}{M_\pi ^2}})).`$
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# 1 Introduction
## 1 Introduction
Heavy quark decays are central to the international effort to test the Standard Model, and the $`b`$ quark has emerged as the focus of this program. These studies include detailed investigations of semileptonic and hadronic decays, as well as increasingly sensitive measurements of rare decays. With major new $`b`$-physics initiatives getting under way at nearly all high energy physics labs, the prospects for definitive tests of the Standard Model, or the discovery of physics beyond it, are excellent.
Flavor physics is interesting because the weak eigenstates of the quarks are mixtures of the mass eigenstates. With three generations the mixing is described by the Cabibbo-Kobayashi-Maskawa matrix (Fig. 1).
Unitarity and the arbitrariness of phases allows the nine complex elements to be reduced to four parameters, as most familiarly parameterized by Wolfenstein . These parameters cannot be predicted, and their determination is one of our most practical needs. Furthermore, redundant measurements provide powerful tests of the validity of the Standard Model.
Specific measurements include direct determinations of the magnitudes of the CKM parameters in a variety of processes, and detailed studies of CP violation, principally in $`s`$ and $`b`$ decays. “Overconstraining” the matrix thus is a matter of measuring the lengths of the sides of the unitarity triangle, as well as its angles, $`\alpha =arg[V_{td}V_{tb}^{}/V_{ud}V_{ub}^{}]`$, $`\beta =arg[V_{cd}V_{cb}^{}/V_{td}V_{tb}^{}]`$ and $`\gamma =arg[V_{ud}V_{ub}^{}/V_{cd}V_{cb}^{}]`$. We already know quite a bit: $`\lambda 0.22`$, $`A0.8`$, and bounds on $`\rho `$ and $`\eta `$ from past measurements. We urgently need precise determinations.
Another powerful probe of the limits of the Standard Model is provided by rare decays, especially rare $`b`$ decays. There are many observables and many challenging measurements. They require very large data samples and mastery of strong-interaction effects that obscure our view of the underlying electroweak physics.
The objective of this review is to report some of the recent developments in heavy-quark decays, hopefully painting a picture of our overall state of knowledge and the pressing open questions. Not included are the crucial topics of lifetimes and mixing, covered elsewhere in these proceedings . In Section 2 I describe the current status of semileptonic $`B`$ decays and the determination of the CKM parameters $`|V_{cb}|`$ and $`|V_{ub}|`$. The focus of Section 3 is rare charmless decays, both two-body hadronic decays and $`bs\gamma `$. Section 4 addresses the interpretation of the various results and implications for the CKM matrix. In Section 5 I mention a few results and near-term prospects in charm physics. This review ends in Section 6 with a brief summary and a survey of the outlook for the not-too-distant future in.
The roster of experimental players in this business is growing with the first operation of several new facilities: KEK-B/BELLE, PEP-II/BaBar, CESR/CLEO III, and HERA-B. Many recent advances in $`b`$ physics have been made by the CLEO experiment working with $`B`$ mesons just above threshold at the $`\mathrm{{\rm Y}}(4S)`$ resonance. CLEO has two distinct data samples: 3.3 million $`B\overline{B}`$ events in the original CLEO II detector and 6.4 million events obtained since 1996 with CLEO II.V, upgraded to include a silicon vertex detector and other improvements. The data sample for CLEO II.V exceeded the project goal as a result of the excellent performance of the CESR storage ring, which reached a luminosity of $`0.8\times 10^{33}`$ cm<sup>-2</sup>s<sup>-1</sup> by the end of the run. The ALEPH, DELPHI, L3, and OPAL experiments at LEP and the SLD experiment at the SLC, have investigated faster-moving $`B`$’s produced in $`Z^0`$ decays. Each LEP experiment collected roughly 0.9 million $`b\overline{b}`$ pairs. With dramatically improved SLC performance toward the end of its run, SLD was able to obtain about 100 thousand $`b\overline{b}`$’s, with the extra advantages of polarized beams and outstanding vertexing. During Run I the Tevatron experiments D0 and, especially, CDF demonstrated that forefront $`b`$ physics can be done in a $`p\overline{p}`$ environment. CDF’s 100 pb<sup>-1</sup> sample, clean lepton triggers and ability to tag displaced $`B`$ vertices produced competitive measurements not just of lifetimes and mixing, but also of some rare $`B`$ decays. There are also a number of current experiments specializing in charm physics, both in $`e^+e^{}`$ (BES) and in fixed-target mode (FOCUS, SELEX, E789, E791). Results from these are beginning to emerge, and the next few years should see many interesting developments.
## 2 $`B`$ Semileptonic Decays
$`B`$ physics is all about Standard Model tests and the determination of CKM parameters, and semileptonic decays are the core of this program. Precise measurements of $`|V_{cb}|`$ and $`|V_{ub}|`$ are the main goals. Since semileptonic decays are our main tool, it is essential that we understand this tool very well. The last few years have seen important developments in both theory and experiment. We have benefited greatly from the increasingly sophisticated application of new theoretical techniques, including Heavy Quark Effective Theory (HQET) and lattice gauge calculations. There has been enhanced coordination between experimentalists and theorists, and the more recent formation of inter-experiment working groups is also proving fruitful. The challenge has been recognized as having two essential components: the extraction of all possible information from the package of measurements, and consistent and realistic assessment of theoretical uncertainties. In this section I address three main topics in $`B`$ semileptonic decays. First I review some long-standing puzzles in the measurements. Following that I assess the state of knowledge of $`|V_{cb}|`$ and of $`|V_{ub}|`$.
### 2.1 Puzzles in Semileptonic $`B`$ Decays
Inclusive semileptonic $`B`$ decay is a beautifully simple process. Inclusive $`bc\mathrm{}\nu `$ provides the most straightforward way to determine $`|V_{cb}|`$, one which is again acknowledged as competitive with exclusive determinations. Inclusive $`bu\mathrm{}\nu `$ gave us the first demonstration that $`V_{ub}`$ is nonzero , and while its interpretation is fraught with model uncertainties, it remains an important measurement. Figure 2 is CLEO’s snapshot of the entire picture of
semileptonic $`B`$ decay in the near-threshold environment of the $`\mathrm{{\rm Y}}(4S)`$. The semileptonic branching fraction and the shape of the lepton momentum spectrum are determined using a lepton-tagged procedure in which charge and angular correlations allow separation of the primary $`b\mathrm{}`$ and secondary $`bc\mathrm{}`$ leptons . Evidence for charmless decays is revealed as an excess in the region of the kinematic end point of the $`bc\mathrm{}\nu `$ lepton spectrum .
The simplicity of the semileptonic decay makes it all the more vexing that it has been the cause of a great deal of anxiety. There are two main puzzles. Why is the $`B`$ semileptonic branching fraction measured at the $`\mathrm{{\rm Y}}(4S)`$ so small? Why is the $`B`$ semileptonic branching fraction measured at the $`\mathrm{{\rm Y}}(4S)`$ smaller than that at the $`Z^0`$? The left-hand graph in Fig. 3 shows an assessment
by Neubert of the problem as of about two years ago . Naive considerations suggest a $`B`$ semileptonic branching fraction of at least 12%, while experiment has consistently given values smaller than this. Mechanisms that enhance hadronic $`B`$ decays can reduce the semileptonic branching fraction, but only by increasing $`n_c`$, the number of charm quarks per $`B`$. The data from the $`\mathrm{{\rm Y}}(4S)`$ did not bear this out. The fact that the branching fraction is smaller at the $`\mathrm{{\rm Y}}(4S)`$ is a separate matter that is also quite perplexing. The dominant $`B`$ mesons at the $`Z^0`$ are the same as those at the $`\mathrm{{\rm Y}}(4S)`$, and the inclusion of higher-mass $`b`$-flavored particles at higher energy would be expected to reduce the average semileptonic branching fraction.
No new data from CLEO have been presented since 1997. There have been new developments on both the semileptonic branching fraction and $`n_c`$ fronts from the LEP experiments.
DELPHI (, L3 and OPAL have all presented new measurements of the $`B`$ semileptonic branching fraction. They use a variety of techniques with second-lepton, $`B`$-vertex and jet-charge tagging, with neural nets employed to separate primary, secondary and background leptons. L3 uses two separate analyses based on double-tag methods to determine simultaneously the $`Z^0`$ $`b`$-quark fraction $`R_b`$ and $`(BX\mathrm{}\nu `$). One analysis uses a displaced-vertex $`b`$ tag, while the other demands a high-$`p_t`$ lepton. The observed lepton $`p_t`$ distributions and the unfolded momentum in the $`b`$ rest frame are shown in Fig. 4. The result of this
measurement of $`(BX\mathrm{}\nu `$), and of the most recent measurements from the other four LEP experiments, are summarized in Table 1.
The new $`Z^0`$ average, $`(BX\mathrm{}\nu )=(10.63\pm 0.17)\%`$, represents a significant decrease from the $`(11.1\pm 0.3)\%`$ average of the high-energy measurement circa 1997. The PDG value for the semileptonic branching fraction at the $`\mathrm{{\rm Y}}(4S)`$ is $`(10.45\pm 0.21)\%`$, quite compatible with the new $`Z^0`$ average. This may still overstate the difference, however, since the PDG average has an aggressively small error considering the experimental and theoretical errors of the input measurements and the spread among them.
New measurements of the multiplicity of charm quarks per $`b`$ decay have been reported by ALEPH and DELPHI . Combining these with an earlier OPAL measurement leads to a new correlated average of $`n_c=1.151\pm 0.022\pm 0.022\pm 0.051`$ , where the errors are statistical, systematic, and that due to input branching fractions. CLEO’s previous number, $`n_c=1.10\pm 0.05`$ becomes $`1.12\pm 0.05`$ when consistent branching fractions are used, again in very good agreement.
The right-hand graph of Fig. 3 is an update of Neubert’s original comparison. It is clear from that graph that the gap between high-energy and low-energy measurements has narrowed considerably. The low-energy data still lies outside the theory comfort zone, but the puzzle seems much less compelling than it did previously.
### 2.2 Determination of $`|V_{cb}|`$
We determine the CKM parameter $`|V_{cb}|`$ by two techniques, both involving semileptonic decays $`bc\mathrm{}\nu `$. The favored method has been to use the rate for the exclusive semileptonic decay $`BD^{}\mathrm{}\nu `$ (or $`BD\mathrm{}\nu `$) at zero recoil. A method that languished in disrepute for some years, but which has been rehabilitated, is to use the inclusive semileptonic decay rate. Both approaches are rooted in HQET, and there is extensive theoretical guidance on extracting $`|V_{cb}|`$ and estimating its uncertainty .
The connection between $`V_{cb}`$ and the semileptonic width $`\mathrm{\Gamma }_{SL}`$ from HQET and the operator product expansion (OPE) is as follows:
$$\mathrm{\Gamma }_{SL}(B)=\frac{G_F^2m_b^5|V(cb)|^2}{192\pi ^3}\left[z_0\left(1\frac{\mu _\pi ^2\mu _G^2}{2m_b^2}\right)2\left(1\frac{m_c^2}{m_b^2}\right)^4\frac{\mu _G^2}{m_b^2}\frac{2\alpha _S}{3\pi }z_0^{(1)}+\mathrm{}\right]$$
(1)
Three HQET parameters appear in the expansion. $`\overline{\mathrm{\Lambda }}`$ connects the quark mass with the meson mass. $`\mu _\pi ^2`$ (or its relative $`\lambda _1`$) relates to the average kinetic energy of the $`b`$ quark. $`\mu _G^2`$ ($`\lambda _2`$) is connected to the hyperfine splitting. Bigi judges that a “prudent” theoretical uncertainty for the extraction of $`|V_{cb}|`$ by this procedure is $`6\%`$ . The contributions of the uncertainties in the experimental inputs, the $`B`$ semileptonic branching fraction ($`(10.5\pm 0.2\pm 0.4)\%`$) and average $`B`$ lifetime ($`1.61\pm 0.02`$ ps), are small in comparison. The result is $`|V_{cb}|=(40.\pm 0.4\pm 2.4)\times 10^3`$.
On the exclusive front, HQ symmetry tells us that a heavy-light meson decaying at rest really isn’t changing at all. A measurement of the decay rate of $`BD^{}\mathrm{}^{}\overline{\nu }`$ at maximum $`q^2`$ ($`w=1`$) gives $`(1)V_{cb}`$:
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dw}}`$ $`=`$ $`{\displaystyle \frac{G_F^2}{48\pi ^3}}|V(cb)|^2m_D^{}(m_Bm_D^{})^2^2(w)𝒢(w)\text{, where}`$ (2)
$`w`$ $`=`$ $`(m_B^2+m_D^{}^2q^2)/(2m_Bm_D^{})`$
Add the form-factor normalization $`(1)`$ from theory and we are done. There has been continuing evolution in thinking about $`(1)`$, and some controversy . Bigi suggests $`(1)=0.88\pm 0.08`$, with a smaller value and a much bigger error than earlier suggestions.
A new measurement of $`BD^{}\mathrm{}^{}\overline{\nu }`$ has been reported by DELPHI , joining ALEPH , OPAL , and CLEO . The new DELPHI measurement (Fig. 5) is based on $`5500`$ tagged decays and has the best precision. (CLEO has so far reported on only one sixth of its total data sample.)
Table 2 summarizes the results on $`|V_{cb}|`$ from $`BD^{}\mathrm{}^{}\overline{\nu }`$, following the LEP $`V_{cb}`$ working group , and Bigi’s
proposal for $`(1)`$. Everything agrees very well. The exclusive $`|V_{cb}|`$ result is consistent with the inclusive, and the overall precision is comparable.
Both extraction procedures rely on the HQET/OPE approach, which is beautiful but largely unvalidated by experiment. Experimental tests are needed, and measurements of the parameters $`\overline{\mathrm{\Lambda }}`$ and $`\lambda _1`$/$`\mu _\pi ^2`$ would be extremely valuable. Measurements of the moments of the hadronic mass and lepton energy in $`B`$ decays have been proposed to do this . CLEO has made a preliminary measurement of this type , the results of which are shown in in Fig. 5. The mass moments and lepton-energy moments do not admit a common solution, and the discrepancy is significant. Perhaps one (or both) of the measurements is flawed, or perhaps there is something wrong with the theoretical approach. Some have suggested that the assumption of quark/hadron duality should be scrutinized. CLEO is updating its measurements with more data and a better understanding of the experimental systematics.
### 2.3 Determination of $`|V_{ub}|`$
Compared to $`|V_{ub}|`$, $`|V_{cb}|`$ was easy. The advantages afforded by heavy-quark symmetry in studying $`bc\mathrm{}\nu `$ do not carry over to the heavy-to-light transition of $`bu\mathrm{}\nu `$. Extraction of $`|V_{ub}|`$ is highly model-dependent, the experiments are tougher, and the achievable precision will likely always be less. The CLEO and ARGUS discovery measurements for $`bu\mathrm{}\nu `$ , and the subsequent confirmation in CLEO II data were based on the nonzero excess of leptons near and above the kinematic limit for $`bc\mathrm{}\nu `$ at the $`\mathrm{{\rm Y}}(4S)`$. The measurement of the yield is straightforward, but since only a tiny corner of the $`bu\mathrm{}\nu `$ phase space is sampled, models must be used to extrapolate to the total rate. It is very difficult to assess the theoretical uncertainty, and my preference is to be very cautious: $`|V_{ub}/V_{cb}|=0.08\pm 0.02`$.
In the past few years, ALEPH , L3 and DELPHI have all presented ambitious analyses that seek to measure the $`bu\mathrm{}\nu `$ component in $`b`$ decays at the $`Z^0`$. The strategy is to reconstruct the charmless hadronic mass $`m_X`$ in $`bX\mathrm{}\nu `$, and to enrich the sample in $`bu\mathrm{}\nu `$ by demanding $`m_X`$ to be less than $`1.6\mathrm{GeV}/c^2`$. Discrimination between $`bu`$-like and $`bc`$-like is based on many event details, including displaced vertices, transverse momentum, presence of kaons, and other features, combined for maximum discrimination with neural nets. This technique exploits the advantages of production at the $`Z^0`$: well-separated jets and fast-moving $`B`$’s, but it requires very detailed understanding of $`bc\mathrm{}\nu `$. The DELPHI $`bu\mathrm{}\nu `$ lepton-energy distribution in the $`B`$ rest frame is shown in Fig. 6. It is fitted to signal and background components to extract $`|V_{ub}/V_{cb}|`$.
The LEP $`|V_{ub}|`$ working group combines the three LEP measurements to obtain an average of $`|V_{ub}|=(4.05_{0.74}^{+0.62})\times 10^3`$, very consistent with CLEO. Because more of the spectrum is measured than in the end-point analysis, the extraction of $`|V_{ub}|`$ should have less theoretical uncertainty in principle. Unfortunately, dealing with the enormous $`bc\mathrm{}\nu `$ component introduces different uncertainties that are also very difficult to quantify.
The first measurement of the exclusive charmless semileptonic decays $`B\pi /\rho \mathrm{}\nu `$ by CLEO was a milestone in the determination of $`|V_{ub}|`$. Conventional wisdom has held that the extraction of $`|V_{ub}|`$ from exclusive decays would be less model-dependent than the earlier end-point measurements. The main reason for this prejudice has been that tools like light-cone sum rules and lattice QCD, along with experimental input from charm decays, would provide necessary form-factor information.
CLEO has presented a new analysis of $`B\rho \mathrm{}\nu `$ with higher efficiency than full reconstruction. Binned maximum-likelihood fits are made of the lepton energy, $`\mathrm{\Delta }E`$ and candidate mass to parameterizations for $`B\rho \mathrm{}\nu `$, $`B\pi \mathrm{}\nu `$, $`B\omega \mathrm{}\nu `$, other $`bu\mathrm{}\nu `$, continuum, and fake leptons. The data sample is divided according to lepton-momentum, with the greatest sensitivity to $`B\rho \mathrm{}\nu `$ in the highest bin ($`>2.3\mathrm{GeV}/c`$). Several models are used to evaluate efficiencies and extract $`|V_{ub}|`$. This measurement is averaged with the previously published CLEO results , to give $`(B^0\rho ^{}\mathrm{}^+\nu )=(2.57\pm 0.29_{0.46}^{+0.33}\pm 0.41)\times 10^4`$ and $`|V_{ub}|=(3.25\pm 0.14_{0.29}^{+0.21}\pm 0.55)\times 10^3`$. The experimental uncertainties in this measurement are smaller than the theoretical uncertainty. Further progress will depend on advances in theory with guidance from experiment. This analysis included a first measurement of the $`q^2`$ distribution for $`B\rho \mathrm{}\nu `$, but the data are not yet sufficiently precise to discriminate among models.
## 3 Rare $`B`$ Decays
Rare decays have provided much of the excitement in $`B`$ physics during the past several years. As data samples have grown, the roster of rare processes that have come within the reach of experiment has lengthened steadily. The discovery of the electroweak penguin decay $`bs\gamma `$, first exclusively and later inclusively , was a major milestone in two ways. First, it excluded a broad range of physics beyond the Standard Model by coming in very close to expectations . Second, it was a first signal of the major role of penguin processes in $`B`$ decays, a feature that has greatly influenced expectations for studies of CP violation.
The principal source for new measurements of rare $`b`$ decays has been the nearly 20 million $`B`$ mesons in the CLEO II/II.V data sample. Contributions have also been made by ALEPH, and by CDF and D0 in searches for modes with dileptons.
### 3.1 Charmless Two-Body $`B`$ Decays
CLEO has made great strides in filling in the table of charmless two-body decays. The implications of these measurements for the future $`B`$ program are significant. The principal contributing processes, $`bs`$ penguins and $`bu`$ trees, are shown in Fig. 7.
Interference between tree and penguin diagrams opens a window on the unitarity-triangle angle $`\gamma `$ in measurements of decay rates. CLEO measurements of $`B\pi \pi `$, $`B\pi \rho `$ and other modes define the strategies for future CP-violation searches, including the determination of $`\alpha `$. Searches for direct CP violation could provide our first glimpse of physics beyond the Standard Model.
CLEO’s two-body charmless decay analyses share a common set of tools that take advantage of the features of $`B\overline{B}`$ production at the $`\mathrm{{\rm Y}}(4S)`$. Candidates are identified based on the beam-constrained mass, $`M_B=\sqrt{E_{beam}^2|\stackrel{}{p}|^2}`$, and the difference between the beam energy and the measured energy of the $`B`$ candidate’s daughters, $`\mathrm{\Delta }E=E_1+E_2E_{beam}`$. For signal events $`M_B`$ must be close to the known $`B`$-meson mass ($`\sigma (M_B)2.5`$ MeV), and $`\mathrm{\Delta }E`$ must be close to zero ($`\sigma (\mathrm{\Delta }E)1525`$ Mev, depending on the mode). Two-body $`B`$ decays have considerable background from continuum $`e^+e^{}q\overline{q}`$, for which the cross section is roughly three times higher than $`B\overline{B}`$. The jet-like continuum background is aggressively suppressed with event-shape cuts based on numerous input variables that are combined into a linear multivariate (Fisher) discriminant. Residual continuum background is estimated with data collected 60 MeV below the $`\mathrm{{\rm Y}}(4S)`$ resonance. In addition to the common selection criteria, there are a number of signal-specific cuts including resonance mass, particle identification and helicity angles.
After imposition of loose cuts, final signals are extracted with unbinned maximum likelihood fits to $`7`$ quantities, including $`M_B`$, $`\mathrm{\Delta }E`$, resonant masses ($`\rho `$, $`K^{}`$, $`\eta `$, $`\eta ^{}`$, $`\omega `$), particle ID, helicity angles, and continuum-suppression variables. In addition to the fits, cut-and-count analyses are performed for confirmation. All of CLEO’s new preliminary results have been obtained using between 5.8 million and the full 9.7 million $`B\overline{B}`$ events in the combined CLEO II and CLEO II.V data sets.
Results from CLEO’s updated search for the decays $`Bh^+h^{}`$ are shown in Fig. 8.
The full likelihood fit of $`B\pi ^+\pi ^{}`$ and $`BK^\pm \pi ^{}`$ gives the confidence-level contours shown, with 11.7$`\sigma `$ and 4.2$`\sigma `$ statistical significance, respectively. (There is no evidence for $`BK^+K^{}`$.) The projected signal for the $`K^\pm \pi ^{}`$ mode is extremely strong, while that for $`\pi ^+\pi ^{}`$ is much less compelling. Projections onto $`\mathrm{\Delta }E`$ for the undifferentiated $`h^+h^{}`$ events analyzed as $`\pi ^+\pi ^{}`$ again show the dominance of $`K^\pm \pi ^{}`$, with no clear indication of a $`\pi ^+\pi ^{}`$ signal. For the events satisfying particle-ID criteria for $`\pi ^+\pi ^{}`$, however, it is clear that a $`\pi ^+\pi ^{}`$ component is needed in addition to the misidentified $`K^\pm \pi ^{}`$ to explain the distribution. The results for the branching fractions of these modes are given in Table 3.
This first measurement of $`B\pi ^+\pi ^{}`$ provides a long-awaited piece of the rare-decay puzzle, and it confirms that studies of this mode, and its future use in CP-violation measurements, are greatly complicated by “penguin pollution.” This study is only one piece of a growing picture, however, and CLEO has also presented new results on the closely related decay modes $`B^+K^0h^+`$ and $`B^+h^+\pi ^0`$, also summarized in Table 3 . In this case there are statistically significant signals for $`K^0\pi ^+`$ ($`7.6\sigma `$) and $`K^+\pi ^0`$ ($`6.1\sigma `$), but not for $`\pi ^+\pi ^0`$, reinforcing the picture of penguin dominance. In addition, there has been a first neasurement of the decay to $`K^0\pi ^0`$, providing a complete set of four $`K\pi `$ branching fractions.
As the available sample of charmless hadronic $`B`$ decays grows, it becomes possible to search for direct CP violation. CP asymmetries are possible when two or more contributing diagrams differ in weak and strong phases. CLEO has presented preliminary measurements the of asymmetry $`𝒜\frac{(bf)(\overline{b}\overline{f})}{(bf)+(\overline{b}\overline{f})}`$ for five charmless two-body final states $`f`$ ($`K^{}\pi ^+`$, $`K^{}\pi ^0`$, $`K_s^0\pi ^+`$, $`K^{}\eta ^{}`$, $`\omega \pi ^+`$) . Within the Standard Model, theoretical expectations for these asymmetries range up to $`0.10`$ . Using the full CLEO II/II.V data sample, the statistical precision on $`𝒜`$ is between $`\pm 0.12`$ and $`\pm 0.25`$ for the modes studied. While these measurements are not yet a powerful test of the Standard Model, increasing event samples could render the larger asymmetries measurable within a few years.
The growing recognition that $`B\pi \pi `$ will not provide an easy route to the unitarity triangle parameter $`\alpha `$ has stimulated the search for alternatives. The most promising avenue was suggested by Snyder and Quinn . They observe that a full Dalitz analysis of $`B\pi ^+\pi ^{}\pi ^0`$ exploits interference among the different $`B\rho \pi `$ modes to remove ambiguities due to unknown phases. This provides a determination of $`\alpha `$ to within about $`6^{}`$ with a sample of $`1000`$ $`B\rho \pi `$ decays, assuming the sample to be essentially background-free.
CLEO has presented preliminary results of searches for $`B`$ decays into final states with a $`K^{}`$, $`\rho `$, $`\omega `$, or $`\varphi `$ meson and a second low-mass meson . The results for all modes investigated are summarized in Table 4.
Fig. 9 shows the likelihood contours and beam-constrained mass distributions for $`B^+\rho ^0h^+`$ candidates.
there is clear evidence of a signal, which translates into a branching fraction measurement $`(B^+\rho ^0\pi ^+)=(1.5\pm 0.5\pm 0.4)\times 10^5`$. There is no significant signal for $`B^+\rho ^0K^+`$, with a 90% confidence-level upper limit of $`(B^+\rho ^0K^+)<2.2\times 10^5`$. The situation is similar for $`B^0\rho ^\pm h^{}`$, for which the significant $`B^0K^+\pi ^{}`$ background demands a cut on the helicity angle. Again there is a measurement for the $`\rho \pi `$ mode ($`(B^0\rho ^\pm \pi ^{})=(3.5_{1.0}^{+1.1}\pm 0.5)\times 10^5`$), and only an upper limit for $`\rho K`$ ($`(B^0\rho ^\pm K^{})<2.5\times 10^5`$ at 90% confidence level).
These measurements allow us to assess the feasibility of measuring $`\alpha `$ with $`B\rho \pi `$. More than 100 fb<sup>-1</sup> will be needed to obtain the specified 1000 events. This sample will require several years of an asymmetric $`B`$ factory to accumulate, and the need to reduce and understand backgrounds will be a major challenge.
Among the other measurements reported in Ref. is the intriguing observation of the decay $`B^+\omega \pi ^+`$. Fig. 10
shows the projected distributions of beam-constrained mass and $`\mathrm{\Delta }E`$. The signal is solid, and the measured branching fraction $`(B^+\omega \pi ^+)=(11.3_{2.9}^{+3.3}\pm 1.5)\times 10^6`$ agrees well with CLEO’s measurement for $`B\rho ^0\pi ^+`$, as expected from isospin. While the branching fraction for $`B^+\omega \pi ^+`$ is consistent with CLEO’s published upper limit on this mode , the new upper limit on $`B\omega K^+`$ conflicts with the previously reported observation . There is no obvious explanation for this change other than a fluctuation in the previous search. The new measurement is an improvement over the first in several ways. The data sample has almost tripled and analysis-procedure improvements have increased the reconstruction efficiencies by between 10% and 20%.
A “poster child” for the challenge of interpreting charmless hadronic $`B`$ decays is the decay $`B\eta ^{}K`$. In 1998 CLEO reported an unexpectedly large branching fraction for this mode , stimulating considerable theoretical interest. An updated search for two-body $`B`$ decays to $`\eta `$ and $`\eta ^{}`$ has now been reported , and the mystery has not gone away. The distributions of beam-constrained mass for $`B\eta K^{}`$ and $`B\eta ^{}K`$ are shown in Fig. 11. The branching fractions for the modes with clear signals are $`(B^+\eta ^{}K^+=(6.5_{1.4}^{+1.5}\pm 0.9)\times 10^5`$ and $`(B^0\eta ^{}K^0=(4.7_{1.4}^{+2.7}\pm 0.9)\times 10^5`$. These exceed all theoretical predictions . No statistically significant signal is observed among the 17 other modes involving $`\eta `$ and $`\eta ^{}`$. A few of these limits are impinging on the predicted Standard Model range.
Fig. 12 shows summary graphs for all CLEO-measured
rare two-body $`B`$-decay processes. Comparisons with theoretical predictions are included. Perhaps the most impressive feature of the work done is the breadth of the set of modes that have been measured. This prepares us for global analyses of rare charmless hadronic decays in which multiple measurements of related modes are used to extract detailed information about the CKM matrix. I return to this question in Sec. 4.
### 3.2 $`bs\gamma `$ and $`bs\mathrm{}^+\mathrm{}^{}`$
Inclusive measurements of $`bs\gamma `$ provide powerful constraints on physics beyond the Standard Model. CLEO has recently presented an updated analysis of 3.3 million $`B\overline{B}`$ events . The technique is an amalgam of continuum suppression through shape variables with a neural net and pseudo-reconstruction of $`BX_s\gamma `$. For the latter, the $`X_s`$ consists of a charged or neutral kaon and up to four pions, one of which can be a $`\pi ^0`$. The photon spectrum is shown in Fig. 13, and the branching fraction measurement is $`=(3.15\pm 0.35\pm 0.32\pm 0.26)\times 10^4`$, where the errors are statistical, systematic and model-dependent, respectively. ALEPH has also presented an inclusive measurement of $`bs\gamma `$ . In their analysis non-$`B`$ backgrounds are suppressed with an opposite-hemisphere lifetime tag. As for CLEO a pseudo-reconstruction approach is employed, in which $`BX_s\gamma `$ is assembled from between one and eight tracks, $`K_s^0`$’s and $`\pi ^0`$’s. The photon spectrum is shown in Fig. 13. ALEPH’s
result is $`=(3.11\pm 0.80\pm 0.72)\times 10^4`$, very consistent with CLEO’s.
It has recently been recognized that additional sensitivity to new physics is provided by the rate asymmetry $`𝒜=\frac{\mathrm{\Gamma }(bs\gamma )\mathrm{\Gamma }(\overline{b}\overline{s}\gamma )}{\mathrm{\Gamma }(bs\gamma )+\mathrm{\Gamma }(\overline{b}\overline{s}\gamma )}`$. Some non-Standard Model predictions give asymmetries as large as 40% . CLEO’s updated study of inclusive $`bs\gamma `$ includes an extension of the pseudo-reconstruction analysis to measure this asymmetry. The strangeness content of the $`X_s`$ system can be used to tag the flavor of the decaying $`B`$, but mistags and untaggable states must be carefully accounted for. CLEO finds $`𝒜=(0.16\pm 0.14\pm 0.05)\times (1.0\pm 0.14)`$, with both additive and multiplicative (mistagging rate) systematic uncertainties. The 90% confidence-level range on $`𝒜`$ is $`0.09<𝒜<0.42`$.
A probe of non-Standard Model physics similar to $`bs\gamma `$ is provided by $`bs\mathrm{}^+\mathrm{}^{}`$. CDF dominates the search for the exclusive decays $`BK/K^{}\mu ^+\mu ^{}`$ , with a very large sample of hadronically produced $`B`$’s and the capability to tag $`B`$ production by displaced vertices. Fig. 14 shows the distributions of $`M(\mu ^+\mu ^{})`$ vs.
$`M(K/K^{}\mu ^+\mu ^{}`$) for the CDF data. The background is largely confined to the easily excluded $`J/\psi `$ and $`\psi ^{}`$ dilepton mass bands, leaving a very clean measurement. CDF obtains the 90% confidence limits $`(B^+K^+\mu ^+\mu ^{})<5.2\times 10^6`$ (Standard Model prediction: $`0.30.7\times 10^6`$), and $`(B^0K^0\mu ^+\mu ^{})<4.0\times 10^6`$ (Standard Model: $`14\times 10^6`$). With a data sample of 2 $`fb^1`$ expected for Run II, the observation of this mode should not be far off.
Since interpretation of the exclusive decays is somewhat problematic (as for $`bs\gamma `$), inclusive measurements would have some advantage. Both D0 and CLEO have reported inclusive analyses. For D0 this is a search for lepton pairs with masses between the charmonium region and the $`B`$ meson, a window that includes only a portion of $`bs\mathrm{}^+\mathrm{}^{}`$, but which is very clean. CLEO employs a pseudo-reconstruction technique similar to the $`bs\gamma `$ procedure. All limits obtained are an order of magnitude or more above Standard Model expectation.
## 4 Interpretation – CKM
Information relevant to the determination of the CKM parameters is being accumulated at an accelerating rate. While principal responsibility for its interpretation in these proceedings falls to Adam Falk , I will briefly comment on the conventional view and then highlight a speculative interpretation of CLEO’s rare-$`B`$-decay data.
A number of authors have incorporated the principal constraints from $`B`$ decay ($`|V_{ub}/V_{cb}|`$, $`\mathrm{\Delta }m_d`$ and the limit on $`\mathrm{\Delta }m_s`$) with input from $`K_L^0`$ CP violation ($`|ϵ_K|`$) in global fits to obtain the Wolfenstein parameters $`\overline{\rho }`$ and $`\overline{\eta }`$, and the angles $`\alpha `$, $`\beta `$ and $`\gamma `$ of the unitarity triangle. Parodi et al. and Mele have performed maximum likelihood fits that assign Gaussian errors to several theoretical inputs. The fits of Parodi et al. give the solution shown on the left-hand side of Fig. 15,
leading to $`\overline{\rho }=0.202_{0.059}^{+0.053}`$ and $`\overline{\eta }=0.340\pm 0.035`$, which in turn give $`\mathrm{sin2}\alpha =0.26_{0.28}^{+0.29}`$, $`\mathrm{sin2}\beta =0.725_{0.060}^{+0.050}`$ and $`\gamma =(59.5_{7.5}^{+8.5})`$ degrees.
Stone has pointed out the danger of underestimating the overall uncertainty when assuming Gaussian errors for theoretical inputs . Plaszczynski has taken a much more cautious approach, considering all theoretical models on an equal basis and presenting the full spread in the resulting parameter values as shown in Fig. 15. He obtains the larger ranges $`0\overline{\rho }0.3`$ and $`0.2\overline{\eta }0.45`$, leading to $`0.50\mathrm{sin2}\beta 0.85`$ and $`0.95\mathrm{sin2}\alpha 0.50`$. Of course everything is consistent with CDF’s first direct measurement, $`\mathrm{sin2}\beta =0.79_{0.44}^{+0.41}`$ .
Because rare hadronic $`B`$ decays incorporate both penguin and $`bu`$ tree processes, their rates and CP asymmetries carry information about weak phases. In particular, it has been suggested by several authors that combinations of measured rates can be used to extract the value of $`\gamma `$, the phase of $`V_{ub}^{}`$. The first suggestions focused on the $`BK\pi `$ branching fractions, but these approaches do not set significant bounds on $`\gamma `$ with current data.
A much more aggressive procedure to extract maximal information from the data has been suggested by Hou, Smith and Würthwein . They assume factorization holds and write the $`B`$-decay amplitudes in terms of five parameters: $`\gamma =Arg(V_{ub}^{})`$, $`|V_{ub}/V_{cb}|`$, $`R_{su}`$ (incorporating information about quark masses), $`F^{B\pi }`$ ($`B\pi `$ form factor), and $`A_0^{B\rho }`$ ($`B\rho `$ form factor). The CP-averaged branching fractions for 14 of CLEO’s measured charmless two-body decays were fitted with this parameterization. (Modes with $`\eta `$ and $`\eta ^{}`$ were excluded based on their anomalies.) The constraint $`|V_{ub}/V_{cb}|=0.08\pm 0.02`$ was imposed. The result of the fit is shown in Fig. 16. The
minimum $`\chi ^2`$ (10.3 for 10 degrees of freedom) occurs for $`\gamma =113_{23}^{+25}`$ degrees. The other fit parameters all give very reasonable values. This result agrees with earlier observations that CLEO data favor $`\mathrm{cos}\gamma <0`$ .
This is an intriguing result. The reasonable values of the fit parameters other than $`\gamma `$ suggest that there may be some validity, in spite of the very model-dependent assumptions. On the other hand, this may be nothing more than a misleading coincidence. Skepticism is appropriate.
## 5 Other Topics
It is impossible to report exhaustively on all of the activity in heavy-quark decays within a single review. While my focus has been on CKM tests and measurements relevant to CP violation, there is other work that is also having impact. Even in the $`B`$ sector I have had to ignore some work, including CLEO and LEP studies of hadronic decays and exclusive semileptonic decays to charm, that are important elements of a comprehensive understanding of $`B`$ decay.
Charm physics remains an extremely valuable complement to $`b`$ physics in our program of Standard Model testing. Because the expected rates for rare FCNC processes are extremely small, the potential to see new physics in $`D\overline{D}`$ mixing, in rare $`D`$-meson decays or in CP-violating processes is great. Additionally, studies of semileptonic and leptonic charm decays are an important adjunct to the CKM measurements, with potential to reduce model uncertainties in the extraction of $`V_{ub}`$ and other parameters. A number of experiments have presented new results on charm decays, with much more on the way. This is a very broad program, the components of which have been the subject of several excellent recent reviews, including that of lifetimes and mixing elsewhere in these proceedings .
As in $`B`$ physics, studies of hadronic charm decays are important for developing a comprehensive understanding of heavy flavors, probing questions of final-state interactions and interference effects. Both meson and baryon decays are useful in this effort, and previously reported results from E791, CLEO and other experiments will be greatly enhanced by FOCUS and SELEX.
A number of new form-factor measurements for semileptonic $`D`$ and $`D_s`$ decays have been presented in the past year by E687 . FOCUS will soon have results from larger samples, and CLEO will also extend previous studies to their full data set. In tandem with HQET these measurements will significantly reduce model uncertainties in the extraction of $`|V_{ub}|`$ from data on semileptonic $`B`$ decays. Measurements of heavy-meson decay constants in leptonic decays of charmed mesons provide input to $`B`$-physics analyses and tests of lattice calculations.
New limits on rare or forbidden charm decays have been presented by E791. A blind search for 24 modes was performed, with no signals observed in any and 90% confidence-level upper limits that range from $`10^3`$ or $`10^4`$ for $`K/\pi \mathrm{}^+\mathrm{}^{}`$ to less than $`10^5`$ for $`\mathrm{}^+\mathrm{}^{}`$. Again FOCUS will benefit from much greater statistics, with improvements in sensitivity for these modes of an order of magnitude or better.
## 6 Summary and Conclusion
The past several years have seen steady progress on a broad program of Standard Model tests in $`B`$ decays, but there remains much to be done.
The embarrassment of the $`Z^0/\mathrm{{\rm Y}}(4S)`$ disagreement on the $`B`$ semileptonic branching fraction has eased. The basic experimental observation that there are too few semileptonic decays for the observed multiplicity of charm quarks is still with us, but it is not of crisis proportions. Theoretical tools for describing semileptonic decays have matured, but underlying assumptions like quark-hadron duality must be scrutinized. Hints of inconsistency between HQET-inspired interpretations of CLEO’s hadronic-mass and lepton-energy moments in semileptonic $`B`$ decays are troubling. A great deal more data and a great deal of work will be required of to reach final conclusion on the values of $`V_{ub}`$ and $`V_{cb}`$. Intensive theory/experiment collaboration is a big plus.
In rare $`B`$ decays we have a number of major developments. The decay $`B\pi \pi `$ has been observed, and the rare hadronic decay picture is filling in with more measurements and tighter limits.
We stand on the verge of truly powerful tests of the Standard Model. First efforts to measure CP asymmetries and CDF’s first measurement of $`\mathrm{sin2}\beta `$ are opening salvoes in the next phase of the campaign to make redundant measurements of the sides and angles of the unitarity triangle. So far fits to the usual experimental constraints show the Standard Model to be holding up well, but this is only the beginning.
The exciting future of heavy flavor physics is well documented elsewhere in these proceedings. The three $`e^+e^{}`$ $`B`$ factories, complemented by the upgraded Tevatron detectors, will produce a wealth of new physics. It is to be hoped that these facilities, their successor $`e^+e^{}`$ machines of still higher luminosity, and specialized detectors at hadron colliders, will carry us well beyond the Standard Model.
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# 1 *
SLAC–PUB–8327
February 2000
BIT-STRING PHYSICS PREDICTION OF $`\eta `$, THE DARK MATTER/BARYON RATIO AND $`\mathrm{\Omega }_M`$ <sup>*</sup><sup>*</sup>*Work supported by Department of Energy contract DE–AC03–76SF00515.
H. Pierre Noyes
Stanford Linear Accelerator Center
Stanford University, Stanford, CA 94309
## Abstract
Using a simple combinatorial algorithm for generating finite and discrete events as our numerical cosmology, we predict that the baryon/photon ratio at the time of nucleogenesis is $`\eta =1/256^4`$, $`\mathrm{\Omega }_{DM}/\mathrm{\Omega }_B=12.7`$ and (for a cosmological constant of $`\mathrm{\Omega }_\mathrm{\Lambda }=0.6\pm 0.1`$ predicted on general grounds by E.D.Jones) that $`0.325>\mathrm{\Omega }_M>0.183`$. The limits are set not by our theory but by the empirical bounds on the renormalized Hubble constant of $`0.6<h_0<0.8`$. If we impose the additional empirical bound of $`t_0<14Gyr`$, the predicted upper bound on $`\mathrm{\Omega }_M`$ falls to $`0.26`$. The predictions of $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$were in excellent agreement with Glanz’ analysis in 1998, and are still in excellent agreement with Lineweaver’s recent analysis despite the reduction of observational uncertainty by close to an order of magnitude.
Contributed paper presented at DM2000
Marina del Rey, California, February 23-25, 2000.
First Afternoon Session, Thursday, February 24
The theory on which I base my predictions is unconventional. Hence it is easier for me to show you first the consequences of the predictions in comparison with observation, in order to establish a presumption that the theory might be interesting, and then show you how these predictions came about.
The predictions are that (a) the ratio of baryons to photons was $`\eta =1/256^4=2.328\mathrm{}\times 10^{10}=10^{10}\eta _{10}`$ at the time of nucleogenesis, (b) $`\mathrm{\Omega }_{DM}/\mathrm{\Omega }_B=127/10=12.7`$ and (c) $`\mathrm{\Omega }_\mathrm{\Lambda }=0.6`$. Comparison of prediction (a) with observation is straightforward, as is illustrated in Figure 1.
Comparison with observation of prediction (b) that the ratio of dark to baryonic matter is not straightforward, as was clear at DM98; I suspect that is matter will remain unresolved at this conference (DM2000). However, according to the standard cosmological model, the baryon-photon ratio remains fixed after nucleogenesis. In the theory I am relying on, the same is true of the of the dark matter to baryon ratio. Consequently, if we know the Hubble constant, and assume that only dark and baryonic matter contribute, the normalized matter parameter $`\mathrm{\Omega }_M`$ can also be predicted, as we now demonstrate.
We know from the currently observed photon density (calculated from the observed $`2.728^oK`$ cosmic background radiation) that the normalized baryon density is given by
$$\mathrm{\Omega }_B=3.67\times 10^3\eta _{10}h_0^2$$
(1)
and hence, from our prediction and assumptions about dark matter, that the total mass density will be 13.7 times as large. Therefore we have that
$$\mathrm{\Omega }_M=0.11706h_0^2.$$
(2)
Hence, for $`0.8h_00.6`$ , $`\mathrm{\Omega }_M`$ runs from $`0.18291`$ to $`0.32517`$. This clearly puts no restriction on $`\mathrm{\Omega }_\mathrm{\Lambda }`$.
Our second constraint comes from integrating the scaled Friedman-Robertson-Walker (FRW) equations from a time after the expansion becomes matter dominated with no pressure to the present. Here we assume that this initial time is close enough to zero on the time scale of the integration so that the lower limit of integration can be approximated by zero . Then the age of the universe as a function of the current values of $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ is given by
$`t_0`$ $`=`$ $`9.77813h_0^1f(\mathrm{\Omega }_M,\mathrm{\Omega }_\mathrm{\Lambda })Gyr`$ (3)
$`=`$ $`9.77813h_0^1f(0.11706h_0^2,\mathrm{\Omega }_\mathrm{\Lambda })Gyr`$
where
$$f(\mathrm{\Omega }_M,\mathrm{\Omega }_\mathrm{\Lambda })=_0^1𝑑x\sqrt{\frac{x}{\mathrm{\Omega }_M+(1\mathrm{\Omega }_M\mathrm{\Omega }_\mathrm{\Lambda })x+\mathrm{\Omega }_\mathrm{\Lambda }x^3}}.$$
(4)
For the two limiting values of $`h_0`$, we see that
$`h_0`$ $`=`$ $`0.8,t_0=12.223f(0.18291,\mathrm{\Omega }_\mathrm{\Lambda })Gyr`$
$`h_0`$ $`=`$ $`0.6,t_0=16.297f(0.32517,\mathrm{\Omega }_\mathrm{\Lambda })Gyr.`$ (5)
The results are plotted in Figure 2. We emphasize that these predictions were made and published over a decade ago when the observational data were vague and the theoretical climate of opinion was very different from what it is now. The figure just given was presented at ANPA20 (Sept. 3-8, 1998, Cambridge, England) and given wider circulation in. The calculation (c) that $`\mathrm{\Omega }_\mathrm{\Lambda }=0.6`$ was made by Jones before there was any observational evidence for a cosmological constant, let alone a positive one. The precision of the relevant observational limits has improved
considerably since DM98. A recent analysis of this new data suitable for our purposes has been made by Lineweaver. His one, two and three $`\sigma `$ contours are plotted in comparison with the previous observational limits and our (unchanged) earlier predictions in Figure 3. Note how dramatically the regions of uncertainty have shrunken in two years. It is gratifying that our prior predictions are still close to the center of the allowed region, indicating that it will take a lot more work to show that they are wrong!
The theory I am using has a long history, starting with the discovery of the combinatorial hierarchy in 1961 and the first publication of the work on this idea by Amson, Bastin, Kilmister and Parker-Rhodes in 1966. The theory is unusual in that it starts from minimal assumptions about what is needed for a physical theory and tries to let the structure of the theory grow out of them. My own preferred choice of basic assumptions are that a physicist must (a) be able to tell something from nothing, (b) be able to tell whether things are the same or different, and (c) must assume a basic arbitrariness in the universe which underlies the stochastic effects exhibited by quantum events. I further assume that we should use the simplest possible mathematical structures to model and develop these concepts. (a) is simply modeled by bit multiplication; (b) is simply modeled by bit addition (addition modulo 2, XOR, symmetric difference,…) or, as it is referred to in the ANPA program, discrimination.
The third requirement, together with the usual scientific assumption that we can keep historical records and examine them at later times, is accomplished by constructing a computer model called program universe which yields a growing universe of ordered strings of the integers “0” and “1”. Here we remind the reader of how we use discrimination (“$``$”) between ordered strings of zeros and ones (bit-strings) defined by
$$(𝐚(W)𝐛(W))_w=(a_wb_w)^2;a_w,b_w0,1;w1,2,\mathrm{}.,W$$
(6)
to generate a growing universe of bit-strings which at each step contains $`P(S)`$ strings of length $`S`$. The algorithm is very simple, as can be seen from the flow diagram in Fig. 4. We start with a rectangular block of rows and columns containing only the bits “0” and “1”. We then pick two rows arbitrarily and if their discriminant is non-null, adjoin it to the table as a new row. If it is null, we simply adjoin an arbitrary column (Bernoulli sequence) to the table and recurse to picking two arbitrary rows. That this model contains arbitrary elements and (if interpretable in terms of known aspects of the practice of physics) an historical record (ordered by the number of TICK’s, or equivalently by the row length) should be clear from the outset. The forging of rules that will indeed connect the model to the actual practice of physics is the primary problem that has engaged me ever since the model was created.
Program universe provides a separation into a conserved set of “labels”, and a growing set of “contents” which can be thought of as the space-time “addresses” to which these labels refer. To see this, think of all the left-hand, finite length $`S`$ portions of the strings which exist when the program TICKs and the string-length goes from $`S`$ to $`S+1`$. Call these labels of length $`L=S`$, and the number of them at the critical tick $`N_0(L)`$. Further PICKs and TICKs can only add to this set of labels those which can be produced from it by pairwise discrimination, with no impact from the (growing in length and number) set of content labels with length $`S_C=SL>0`$. If $`N_IN_0(S_L)`$ of these labels are discriminately independent, then the maximum number of distinct labels they can generate, no matter how long program universe runs, will be $`2^{N_I}1`$, because this is the maximum number of ways we can choose combinations of $`N_I`$ distinct things taking them $`1,2,\mathrm{},N_I`$ times. We will interpret this fixed number of possibilities as a representation of the quantum numbers of systems of “elementary particles” allowed in our bit-string universe and use the growing content-strings to represent their (finite and discrete) locations in an expanding space-time description of the universe.
This label-content schema then allows us to interpret the events which lead to TICK as four-leg Feynman diagrams representing a stationary state scattering process. Note that for us to find out that the two strings found by PICK are the same, we must either pick the same string twice or at some previous step have produced (by discrimination) and adjoined the string which is now the same as the second one picked. Although it is not discussed in bit-string language, a little thought about the solution of a relativistic three body scattering problem Ed Jones and I have found shows that the driving term ($`\genfrac{}{}{0pt}{}{><}{}`$) is always a four-leg Feynman diagram ($`><`$) plus a spectator ($``$) whose quantum numbers are identical with the quantum numbers of the particle in the intermediate state connecting the two vertices. The step we do not take here is to show that the labels do indeed represent quantum number conservation and the contents a finite and discrete version of relativistic energy-momentum conservation. But we hope that enough has been said to show how we could interpret program universe as representing a sequence of contemporaneous scattering processes, and an algorithm which tells us how the space in which they occur expands.
Short-circuiting and reordering the actual route by which my current interpretation of this model was arrived at, we note that the two basic operations in the model which provide locally novel bit-strings (Adjoin and TICK) are isomorphic, respectively, to a three-leg or a four-leg Feynman diagram. This is illustrated in Fig. 5. Note that the internal (exchanged particle) state in the Feynman diagram is necessarily accompanied by an identical (but distinct) “spectator” somewhere else in the (coherent) memory.
We do not have space here to explain how, in the more detailed dynamical interpretation, the three-leg diagrams conserve (relativistic) 3-momentum but not necessarily energy (like vacuum fluctuations) while the four-leg diagrams conserve both 3-momentum and energy and hence are candidates for potentially observable events. We are particularly pleased that the observable events created by Program Universe necessarily provide two locally identical but distinct strings (states) because these are the starting point for a relativistic finite particle number quantum scattering theory which has non-trivial solutions. But we do need to explain how this interpretation of program universe does connect up with the work on the combinatorial hierarchy.
At this point we need a guiding principle to show us how we can “chunk” the growing information content provided by the discriminate closure of the label portion of the strings in such a way as to generate a hierarchical representation of the quantum numbers that these labels represent. Following a suggestion of David McGoveran’s , we note that we can guarantee that the representation has a coordinate basis and supports linear operators by mapping it to square matrices.
The mapping scheme originally used by Amson, Bastin, Kilmister and Parker-Rhodes satisfies this requirement. This scheme requires us to introduce the multiplication operation ($`00=0=01=0=10`$, $`11=1`$), converting our bit-string formalism into the field $`Z_2`$. First note, as mentioned above, that any set of $`n`$ discriminately independent (d.i.) strings will generate exactly $`2^n1`$ discriminately closed subsets (dcss). Start with two d.i. strings $`𝐚`$, $`𝐛`$. These generate three d.i. subsets, namely $`\{𝐚\}`$, $`\{𝐛\}`$, $`\{𝐚,𝐛,𝐚𝐛\}`$. Require each dcss ({ }) to contain only the eigenvector(s), of three $`2\times 2`$ mapping matrices which (1) are non-singular (do not map onto zero) and (2) are d.i. Rearrange these as strings. They will then generate seven dcss. Map these by seven d.i. $`4\times 4`$ matrices, which meet the same criteria (1) and (2) just given. Rearrange these as seven d.i. strings of length 16. These generate $`127=2^71`$ dcss. These can be mapped by 127 $`16\times 16`$ d.i. mapping matrices, which, rearranged as strings of length 256, generate $`2^{127}11.7\times 10^{38}`$ dcss. But these cannot be mapped by $`256\times 256`$ d.i. matrices because there are at most $`256^2`$ such matrices and $`256^22^{127}1`$. Thus this combinatorial hierarchy terminates at the fourth level. The mapping matrices are not unique, but exist, as has been proved by direct construction and an abstract proof . It is easy to see that the four level hierarchy constructed by these rules is unique because starting with d.i. strings of length 3 or 4 generates only two levels and the dcss generated by d.i. strings of length 5 or greater cannot be mapped.
Making physical sense out of these numbers is a long story , and making the case that they give us the quantum numbers of the standard model of quarks and leptons with exactly 3 generations has only been sketched . However we do not require the completely worked out scheme to make interesting cosmological predictions. The ratio of dark to “visible” (i.e. electromagnetically interacting) matter is the easiest to see. The electromagnetic interaction first comes in when we have constructed the first three levels giving 3+7+127 =137 dcss, one of which is identified with electromagnetic interactions because it occurs with probability $`1/137e^2/\mathrm{}c`$. But the construction must first complete the first two levels giving 3+7=10 dcss. Since the construction is “random” and this will happen many, many times as program universe grinds along, we will get the 10 non-electromagnetically interacting labels 127/10 times as often as we get the electromagnetically interacting labels. Our prediction of $`M_{DM}/M_B=12.7`$ is that naive.
The $`1/256^4`$ prediction for $`N_B/N_\gamma `$ is comparably naive. Our partially worked out scheme of relating bit-string events to particle physics , makes it clear that photons, both as labels (which communicate with particle-antiparticle pairs) and as content strings will contain equal numbers of zeros and ones in appropriately specified portions of the strings. Consequently they can be readily identified as the most probable entities in any assemblage of strings generated by whatever pseudo-random
number generator is used to construct the arbitrary actions and bit-strings needed in actually running program universe. This scheme also makes the simplest representation of fermions and anti-fermions contain one more “1” and one less “0” than the photons (or visa versa). (Which we call “fermions” and which “anti-fermions” is, to begin with, an arbitrary choice of nomenclature.) Since our dynamics insures conventional quantum number conservation by construction, the problem — as in conventional theories—is to show how program universe introduces a bias between “0” ’s and “1” ’s once the full interaction scheme is developed.
Since program universe has to start out with two strings, and both of these cannot be null if the evolution is lead anywhere, the first significant PICK and discrimination will necessarily lead to a universe with three strings, two of which are “1” and one of which is “0”. Subsequent PICKs and TICKs are sufficiently “random” to insure that (at least statistically) there will be an equal number of zeros and ones, apart from the initial bias giving an extra one. Once the label length of 256 is reached, and sufficient space-time structure (“content strings”) generated and interacted to achieve thermal equilibrium, this label bias for a 1 compared to equal numbers of zeros and ones will persist for 1 in 256 labels. But to count the equilibrium processes relevant to computing the ratio of baryons to photons, we must compare the labels leading to baryon-photon scattering compared to those leading to photon-photon scattering. This requires the baryon bias of 1 to appear in one and only one of the four initial (or final, since the diagrams are time symmetric) state labels of length 256 involved in that comparison; the two relevant diagrams are illustrated in Fig.6 1, which assumes that the above mentioned interpretation of the strings causing observable TICK’s as four leg Feynman diagrams has been satisfactorily demonstrated. As a trivial example take the baryon-antibaryon-photon vertex to be $`𝐁\overline{𝐁}\gamma =\mathrm{𝟎}`$ with $`𝐁=(1110)`$, $`\overline{𝐁}=(0010)`$ and $`\gamma =(1100)`$. We conclude that, in the absence of further information, $`1/256^4`$ is the program universe prediction for the baryon-photon ratio at the time of big bang nucleosynthesis.
Since Jones’ paper is still in preparation, I am at liberty here only to quote:
> From general operational arguments, Ed Jones has shown how to start from $`N`$ Plancktons and self-generate a universe with $`N^{}`$ baryons which—for appropriate choice of $`N`$—resembles our currently observed universe. In particular it must necessarily have a positive cosmological constant characterized by $`\mathrm{\Omega }_\mathrm{\Lambda }0.6\pm 0.1`$.
We note further that Jones’ general arguments a) are completely compatible with program universe and b) do not in themselves fix the value of $`N`$. Further, the estimate given above, which was made before and independent of the calculations reported in the last section, fell squarely in the middle of the region allowed in 1998 (see Fig. 2), and continues to do so despite the remarkable progress that has been made since DM98 (see Fig.3). Clearly, pursuing the combination of these two lines of reasoning could prove to be very exciting.
APPENDIX
In order to underpin our claim that we can model a finite particle number version of relativistic quantum mechanics with particle creation, etc. using bit-strings we give on the next page the predictions of coupling constants and mass ratios calculated using our theory. As in any mass, length, time theory we are allowed three empirical, dimensional constants which are measured by standard techniques to connect our abstract theory to measurement. These we take to be the mass of the proton $`m_p`$, Planck’s constant $`\mathrm{}`$ and the velocity of light $`c`$. Everything else is calculated. Agreement with observation, given on the next page, is not perfect; we believe it is impressive. For more detail see.
A tentative bit-string representation of the quantum numbers of the (three generation) standard model of quarks and leptons is given on the following page (Fig. 7).
$$G_N^1\frac{\mathrm{}c}{m_p^2}=[2^{127}+136]\times [\mathrm{𝟏}\frac{\mathrm{𝟏}}{\mathrm{𝟑}\mathrm{𝟕}\mathrm{𝟏𝟎}}]=\mathrm{1.693\; 31}\mathrm{}\times 10^{38}$$
$$experiment=1.693\mathbf{58}(21)\times 10^{38}$$
$$\alpha ^1(m_e)=137\times [\mathrm{𝟏}\frac{\mathrm{𝟏}}{\mathrm{𝟑𝟎}\times \mathrm{𝟏𝟐𝟕}}]^\mathrm{𝟏}=137.0359\mathbf{674}\mathrm{}.$$
$$experiment=\mathrm{137.0359\; 895}(61)$$
$$G_Fm_p^2/\mathrm{}c=[256^2\sqrt{2}]^1\times [\mathrm{𝟏}\frac{\mathrm{𝟏}}{\mathrm{𝟑}\mathrm{𝟕}}]=1.02\mathbf{758}\mathrm{}\times 10^5$$
$$experiment=\mathrm{1.02\; 682}(2)\times 10^5$$
$$sin^2\theta _{Weak}=0.25[\mathrm{𝟏}\frac{\mathrm{𝟏}}{\mathrm{𝟑}\mathrm{𝟕}}]^\mathrm{𝟐}=0.2267\mathrm{}$$
$$experiment=0.22\mathrm{𝟓𝟗}(46)]$$
$$\frac{m_p}{m_e}=\frac{137\pi }{<x(1x)><\frac{1}{y}>}=\frac{137\pi }{(\frac{3}{14})[1+\frac{2}{7}+\frac{4}{49}](\frac{4}{5})}=1836.15\mathbf{1497}\mathrm{}$$
(7)
$$experiment=\mathrm{1836.15\; 2701}(37)$$
$$m_\pi ^\pm /m_e=275[\mathrm{𝟏}\frac{\mathrm{𝟐}}{\mathrm{𝟐}\mathrm{𝟑}\mathrm{𝟕}\mathrm{𝟕}}]=273.12\mathbf{92}\mathrm{}$$
$$experiment=\mathrm{273.12\; 67}(4)$$
$$m_{\pi ^0}/m_e=274[\mathrm{𝟏}\frac{\mathrm{𝟑}}{\mathrm{𝟐}\mathrm{𝟑}\mathrm{𝟕}\mathrm{𝟐}}]=264.2\mathbf{143}\mathrm{}$$
$$experiment=\mathrm{264.1\; 373}(6)]$$
$$m_\mu /m_e=3710[1\frac{3}{3710}]=207$$
$$experiment=\mathrm{206.768\; 26}(13)$$
$$G_{\pi N\overline{N}}^2=[(\frac{2M_N}{m_\pi })^21]^{\frac{1}{2}}=[195]^{\frac{1}{2}}=13.96\mathrm{}.$$
$$experiment=13.3(3),orgreaterthan13.9$$
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# Influence of diffraction on the spectrum and wavefunctions of an open system
## I Introduction
In this work, we discuss the transmission spectrum and wavefunctions of an open resonator coupled to a quantum point contact (QPC). The system exhibits both stable and unstable dynamics, depending on the value of a single parameter. The spectral properties of the resonator are determined by the interference of closed (not necessarily periodic) orbits that begin and end at the QPC. Semiclassically, one computes the transmission of such a system with a sum over the classical trajectories. However, it was found experimentally that there were many resonances in the spectrum which did not appear in the theory when classical trajectories involving only specular reflection were considered. We found that the missing resonances were only reproduced when nonclassical closed orbits that undergo diffraction were included into the semiclassical sum for the transmission.
This issue of including diffraction into the semiclassical propagator has been considered by various authors, for both closed systems and open systems . The basis for many of these treatments is the geometric theory of diffraction, originated by Keller . A distinguishing feature of the work presented here is that, in the unstable regime of our resonator, the effect of the diffractive orbits is of the same order as the purely classical orbits. A consequence of this is that there are as many (or more) resonances that are supported by diffractive orbits as are supported by simple classical orbits. This is related to the fact that there is only one classical closed orbit in our system in the unstable regime. This is in stark contrast to the case of a closed, unstable (chaotic) system. Normally, for a closed system, diffraction plays a minor role because of the overwhelming number of nondiffractive periodic orbits present. However, when the system is open, the great majority of long periodic orbits might vanish, if they are allowed to escape the system. In such a case, where only a very small number of classical periodic orbits are present, diffractive orbits may gain in importance and considerably affect the spectrum of the system. This situation is realized by our resonator.
The paper is organized as follows: in Sec. II, we discuss the resonator being studied. In Sec. III, we describe the experimental apparatus. In Sec. IV, we present the experimental results, which are comprised of measured spectra and wavefunctions. In Sec. V we provide a short introduction to the geometric theory of diffraction, which is the theory that describes how diffractive rays contribute to the semiclassical description of a quantum mechanical wavefunction. In Sec. VI we show in detail how the geometric theory of diffraction is incorporated into the semiclassical trace formula. Here theoretical results are seen to be in excellent agreement with measured data. In Sec. VII the physics of the resonator in the time domain is discussed, and again very good agreement between theoretical and experimental observations is found.
In Sec. VIII we discuss the behavior of the system as the transition between stable and unstable dynamics is crossed, as well as analogies of the open resonator with some well-known closed systems, namely the lemon and stadium billiards. Finally, in Sec. IX we discuss the possibility of imaging a pure state quantum mechanical wavefunction in a mesoscopic system with the help of an atomic force microscope, with concluding remarks in Sec. X. A short paper discussing the experimental results presented here has been previously published . See also .
## II The resonator
Recently, Katine studied the transmission behavior of an open quantum billiard in the context of a two dimensional electron gas (2DEG) in a GaAs/AlGaAs heterostructure . Their resonator was formed by a wall with a small aperture (the QPC), and an arc-shaped reflector. A schematic of this resonator is shown in Fig. 1. The voltage on the reflector could be varied, effectively moving the reflector towards or away from the wall. Their measurements showed a series of conductance peaks, analogous to those seen in a Fabry-Perot, as the reflector position was varied.
An interesting property of the resonator considered here is that it is geometrically open, but in the stable regime it is classically closed. In the unstable regime, the resonance properties of the billiard are determined in large part by diffraction.
The resonator shown in Fig. 1 has two distinct modes of operation. When the center of curvature of the reflector is to the left of the wall (the regime studied in ), then almost all classical paths starting from the QPC that hit the reflector remain forever in the region between wall and the reflector: the dynamics is stable and the periodic orbits can be semiclassically quantized. Each quantized mode of the resonator can be characterized by two quantum numbers $`(n,m)`$, which represent the number of radial and angular nodes respectively. As the reflector-wall separation is varied, the conductance exhibits a peak each time one of these quantized modes is allowed. Once an electron is in the resonator, the only way for it to leave is by tunneling back through the QPC or by diffracting around the reflector; since both processes are slow, the resonances have narrow widths. Because the QPC is on the symmetry axis, only modes with even $`m`$ are excited significantly. The states of the resonator in the stable regime bear a strong resemblance to a certain symmetry class of states in a lemon-shaped billiard . This class has even symmetry about the short axis of the lemon, and odd symmetry about the long axis. This connection will be explored more fully in Section VIII.
When the center of curvature is to the right of the wall, however, the dynamics becomes unstable: all classical trajectories beginning at the QPC rapidly bounce out of the resonator, except for a single unstable periodic orbit along the axis of symmetry, which we will call the “horizontal” orbit \[see Fig 1(b)\]. The horizontal orbit is a member or a class of orbits that we call “geometric,” because their paths are governed by specular reflection off the wall and reflector, and do not undergo diffraction. Although the horizontal orbit returns to the QPC, the electron has a low probability of escaping the resonator there because the QPC is much smaller than the de Broglie wavelength of the electron. Because the horizontal orbit is the only periodic orbit in the unstable regime, one might expect resonant buildup only along the symmetry axis. Such a spectrum would be quasi-one-dimensional, with only the half-wavelength periodicity of a Fabry-Perot cavity. However, in numerical simulations it was found that there are other transmission resonances in the unstable regime which did not correspond to any classical periodic orbits . It was proposed that these anomalous peaks are supported by diffraction off the tips of the reflector. Unfortunately, in the mesoscopic experiments, decoherence of the electron wave by impurities in the GaAs/AlGaAs heterostructure shortens the lifetime of the resonances, leaving insufficient energy resolution to resolve the diffractive peaks .
## III Experiment
Because of the problems of dissipation and decoherence in the mesoscopic experiments, we decided to investigate a parallel plate microwave resonator with a similar geometry. In microwave experiments, decoherence and dissipation are not a problem, the geometry of the resonator can be specified much more accurately, and the dynamical range of available wavelengths is much larger. The experimental setup is shown in Fig. 2.
The equation governing the component of the electric field normal to the plates for the TEM mode is identical to the two-dimensional time-independent Schrödinger equation . Therefore, by studying the modes of parallel-plate resonators we can gain insight into the behavior of two-dimensional solutions to the Schrödinger equation.
The resonator consisted of two parallel copper plates, 1 meter square, separated by a distance of 1.25 cm. One side of the resonator consisted of a copper wall. The other three sides were lined with a 11.5 cm thick layer of microwave absorber (C-RAM LF-79, Cuming Microwave Corp.) designed to provide 20 dB of attenuation in the reflected wave intensity in the range 0.6-40 GHz. The absorber prevented outgoing waves from returning to the resonator, thereby simulating an open system in the directions away from the wall. An antenna was inserted normal to the plates, 2 mm from the wall, to simulate the QPC. The curved reflector was formed from a rectangular aluminum rod bent into an arc with radius of curvature $`R=30.5`$ cm. Various opening angles $`\alpha `$ were used: $`115^{}`$, $`112^{}`$, $`109^{}`$, and $`106^{}`$.
Instead of measuring the transmission of the resonator, we measured the reflection back from the antenna; for this we used an HP8720D network analyzer in “reflection” mode (the complex $`S_{11}`$ parameter of the resonator was measured). We inferred the transmission probability $`|T|^2`$ via $`|T|^2=1|R|^2`$, where $`R=S_{11}`$ is the measured reflection coefficient. Because of the proximity of the antenna to the wall, it was only weakly coupled to the resonator; therefore, in the absence of the reflector, the transmission coefficient was close to zero. However, when the reflector was present, the transmission experienced pronounced maxima at certain frequencies.
## IV Results
In Fig. 3 we show a transmission spectrum at fixed frequency, as the distance between the wall and reflector is varied.
In the stable regime, we see that the peaks are narrow and well defined. This is because the dynamics is stable in this regime: nearly all trajectories starting from the QPC that hit the reflector remain forever in the region between the wall and the reflector. In this regime, there exist invariant tori, which may be semiclassically quantized to produce the states of the stable resonator. Such a classical orbit is shown in Fig. 4,
along with its quantum mechanical wavefunction counterpart. We see that the trajectory does not approach the region where the resonator is open. Thus, it behaves as if the cavity were closed—hence the narrow widths of the peaks in the stable regime.
In the unstable regime, the transmission curve is quite different. Here, there are two types of resonances. The first type, labeled $`f`$ in Fig. 3, is related to the horizontal orbit along the axis of symmetry, and bears some resemblance to a Fabry-Perot type resonance between two half-silvered mirrors. The second type, labeled $`d`$, is supported by diffraction off the tips of the reflector.
We verified experimentally that the $`d`$-peaks were indeed supported by diffraction by surrounding the tips of the reflector with microwave absorber and repeating the experiment, as indicated in Fig. 5. When this was done, the Fabry-Perot resonances were unaffected, but the diffractive peaks were entirely eliminated from the spectrum. This makes sense because the gradually thickening absorber smoothly attenuated reflections from rays coming near the tip, leaving no sharp discontinuity from which rays could diffract.
The wavefunctions corresponding to peaks $`f_1`$ and $`d_1`$ were measured using the technique of Maier and Slater . They showed that the frequency shift of a given resonance due to a small sphere of radius $`r_0`$ at a position $`(x,y)`$ is given by
$$\frac{\omega ^2\omega _0^2}{\omega _0^2}=4\pi r_0^3\left(\frac{1}{2}H_0^2(x,y)E_0^2(x,y)\right),$$
(1)
where $`E_0`$ and $`H_0`$ are the unperturbed electric and magnetic fields. Thus, the frequency shift is proportional to the local intensity of the microwave field, and by measuring the shift as a function of the position of the sphere, the field intensity of a particular mode can be mapped out. Note that the frequency shift is positive in regions where the magnetic field is large, and negative where the electric field is large. Also, the factor of $`1/2`$ multiplying the magnetic field in Eq. (1) indicates that the sphere is a stronger perturbation to the electric field than the magnetic field. In our measurements, we found this to be the case: the shifts were predominantly negative. Appreciable positive shifts were found only at the nodes of the electric field, corresponding to maxima of the magnetic field.
Figure 6 shows theoretical quantum wavefunctions compared with experimentally measured frequency shifts for the resonances labeled by $`f_1`$ and $`d_1`$ in Fig. 3. The theoretical wavefunctions were generated using Edwards’ wavelet method presented in .
The measured frequency shift is plotted as a function of sphere position. For these measurements, we used a steel bead of diameter 4.0 mm for the perturbation. The bead was rastered over the inside of the cavity by means of an external magnet. That way, the bead could be moved around inside the cavity without taking the cavity apart. It is important to note that the frequency shift is not proportional to $`E^2`$, but rather to $`H^2/2E^2`$. Therefore we show only negative contour lines below 20% of the maximum negative shift, and thereby emphasize regions of strong electric field. The similarity between theory and experiment is striking.
The wavefunction labeled $`f_1`$ in Fig. 6 is clearly associated with the horizontal orbit along the axis of symmetry. Rays emanating from a point source located on the axis of symmetry next to the wall bounce off the reflector and come to an approximate focus about 10 cm from the source. The focus is approximate because of spherical (or in this case cylindrical) aberration.
Now we turn our attention to the state labeled $`d_1`$ in Fig. 6. As noted above, the only periodic orbit in the unstable regime is the horizontal orbit, along the axis of symmetry. The pictured wavefunction, however, clearly has very little amplitude along this periodic orbit. Instead the wavefunction has a band of higher amplitude running from the region of the tip of the mirror to the QPC, but in the unstable regime there is no classical periodic orbit that does this. Later in the paper, it will be shown that states such as $`d_1`$ are supported by orbits that undergo diffraction off the tips of the reflector. One such orbit is shown in Fig. 1(b). Rays that hit the smooth surfaces of the reflector or wall undergo specular reflection, whereas the rays that hit near the reflector tips can be diffracted. A fraction of the wave amplitude can then return to the QPC from this region, thus setting up a non-classical closed orbit. All peaks labeled with a $`d`$ in Fig. 3 are supported by such diffractive orbits.
Numerical calculations have shown that for energies off resonance, the quantum wavefunction is often intermediate between those shown for $`f_1`$ and $`d_1`$, in the sense that amplitude seems to be running from the QPC to some point between the center and the tip of the reflector .
This can be understood in terms of the interference of paths with each other as they “walk off” the horizontal orbit and escape the resonator. Thus diffraction does not necessarily play a major role in determining the off-resonance wavefunctions. However, diffraction is instrumental in determining the on-resonance wavefunctions underlying the conductance peaks $`d_1`$ and $`d_2`$ in Fig. 3. Figure 7 shows a more global picture of the transmission properties of the resonator.
Here we plot the transmission of the resonator as both the wavelength and the reflector-wall separation are varied. Each vertical slice through this figure is a frequency spectrum with fixed reflector position; the dotted line marks the classical transition from stable to unstable motion that occurs when the reflector’s center of curvature moves to the right of the QPC. The vertical axis indicates how many wavelengths fit along the horizontal orbit between the QPC and the reflector. The repetition of the resonance pattern every half-wavelength in the vertical direction is analogous to the half-wavelength periodicity of a Fabry-Perot cavity.
In the stable regime the peaks have been labeled with their quantum numbers, $`(n,m)`$. Because of the choice of vertical axis, the $`m=0`$ resonance peaks are approximately horizontal in this figure. As the stable/unstable transition is approached, the peaks with high $`m`$ disappear one by one because their large angular sizes allow them to escape around the reflector.
At the stable/unstable transition, all of the resonances in a family would be approximately degenerate, but instead there is an avoided crossing. The level repulsion is caused by a coupling that is partly mediated by diffraction; this subject will be explored more thoroughly in Section VIII.
In Fig. 8 the angular dependence of a diffractive state is shown as the reflector is passed adiabatically through the stable/unstable transition point. In the stable region, most of the amplitude is along the symmetry axis. As the reflector moves further from the wall, the amplitude slowly separates into two lobes, with very little in the center. When the reflector is only slightly in the unstable region, there are bands of amplitude running to a point on the reflector intermediate between the center and the tips, as in the curve for $`(DR)/R=0.016`$ in Fig. 8.
In the unstable regime, the only remaining classical periodic orbit is the horizontal orbit, which itself becomes unstable. The Fabry-Perot peak (labeled $`f`$) is essentially quantized along the horizontal orbit, so its position shows a simple dependence on reflector position. It becomes broad in the unstable regime, with a lifetime given by the classical Lyapunov stability exponent of the horizontal orbit. Two diffractive resonances (labeled by $`d`$), are also visible in each family; they separate from the Fabry-Perot type peak as the reflector is moved away from the wall.
The diffractive peaks labeled by $`d`$ in Fig. 7 cannot be explained by semiclassical theory unless diffraction off the tips of the reflector is included, as will be shown in the following sections.
## V Geometric theory of diffraction
Before we consider the problem of computing semiclassically the transmission properties of our resonator, let us study the simpler problem of diffraction of a plane wave off an infinite half-line in 2D. This problem will serve as a good introduction into the geometrical theory of diffraction, which will be used to include diffraction into the semiclassical propagator.
The problem is illustrated in Fig. 9. A plane wave
$`e^{ikx}`$ is normally incident on the half-line from the left. The half-line extends up from the middle of the figure, indicated by the dark line. We take the tip of the line to be our coordinate origin. Within the geometrical optics approximation, the problem is divided into three separate regions: that of transmission, reflection, and shadow, labelled I, II, and III, respectively. The values of the wavefunction in each region are indicated in the figure as well. In region I, the wave does not hit the wall and thus is unchanged within the geometric optics approximation. In region II, the wave is perfectly reflected and thus a standing sine wave is set up there. In region III, we have a perfect shadow region, which is completely dark. Along the reflection and shadow boundaries indicated, the solution is discontinuous. Of course, these discontinuities are not present in the exact solution; they are an artifact of the geometric approximation. As we shall see, it is the diffraction off the tip of the wall that corrects these discontinuities.
In 1953, Keller showed that one can think of diffraction as originating from a group of “diffracted rays” originating from the edge of the wall . The idea is illustrated in Fig. 10.
Away from the shadow and reflection boundaries, these diffracted rays have the form of an outgoing cylinder wave, multiplied by an angle dependent “diffraction coefficient.” However, Keller’s original theory was shown to be invalid on the reflection and shadow boundaries. A properly uniformized geometric theory of diffraction was developed by Kouyoumjian and Pathak . The diffracted rays are multiplied by a suitable complex number which depends both on the angles of the incident and diffracted rays relative to the wall, as well as the distance from the edge. In this uniformized theory, the solution to the half-line is given by
$$\psi (𝐫)=\psi _\mathrm{g}(𝐫)+D(\theta ,\theta ^{},r,k)e^{ikr},$$
(2)
where $`\psi _\mathrm{g}(𝐫)`$ is the solution given by geometrical optics, shown in Fig. 9. The diffraction coefficient $`D(\theta ,\theta ^{},r,k)`$ is given by ,
$`D(\theta ,\theta ^{},r,k)`$ $`=`$ $`\text{sgn}(a_i)K\left(|a_i|\sqrt{kr}\right)+\text{sgn}(a_r)K\left(|a_r|\sqrt{kr}\right),`$ (3)
$`a_{i,r}`$ $`=`$ $`\sqrt{2}\mathrm{cos}\left({\displaystyle \frac{\theta \theta ^{}}{2}}\right),`$ (4)
and $`K(x)`$ is a modified Fresnel integral:
$$K(x)=\frac{1}{\sqrt{\pi }}e^{ix^2i\pi /4}_x^{\mathrm{}}e^{it^2}𝑑t.$$
(5)
The angles $`\theta ,\theta ^{}`$ are shown in Fig. 11.
In the half-line problem, $`\theta =\pi /2`$, because the incident wave is normal to the wall. In addition, in Eq. (2) we understand that the origin of the coordinate system is at the tip of the half-line.
In Fig. 12 we compare the result of Eq. (2) to the exact solution for the half-line. The prediction of geometric optics is also shown.
There is very good agreement between the exact solution and the uniform geometric theory. Note especially that the discontinuities on the shadow and reflection boundaries are completely removed by the uniform theory.
## VI Semiclassics in the energy domain
Now we turn back to the problem of calculating the transmission properties of our resonator. We need to find an expression for the Green function for the resonator, because the transmission can be easily written in terms of the diagonal part of the energy Green function :
$$T(E)\mathrm{Re}\left[iG(𝐫_{\mathrm{QPC}},𝐫_{\mathrm{QPC}},E)\right],$$
(6)
where $`𝐫_{\mathrm{QPC}}`$ is the center of the QPC. The physical reason that $`G(𝐫_{\mathrm{QPC}},𝐫_{\mathrm{QPC}},E)`$ appears in Eq. (6) is because all waves enter our resonator within a fraction of a wavelength of that point. If a point source of waves is launched at $`𝐫_{\mathrm{QPC}}`$ with energy $`E`$, the complex number $`G(𝐫_{\mathrm{QPC}},𝐫_{\mathrm{QPC}},E)`$ is just the amplitude for returning to $`𝐫_{\mathrm{QPC}}`$. If a significant fraction of the rays emanating from the point source return to $`𝐫_{\mathrm{QPC}}`$ in phase, then $`G(𝐫_{\mathrm{QPC}},𝐫_{\mathrm{QPC}},E)`$ will be appreciable, and the returning waves will beat against the original wave and have a large effect on the transmission.
### A Geometric orbits in the semiclassical propagator
The 2D semiclassical energy Green function $`G_{\mathrm{sc}}(𝐫,𝐫^{},E)`$ can be written as a sum over paths from $`𝐫`$ to $`𝐫^{}`$ thus :
$$G_{\mathrm{sc}}(𝐫,𝐫^{},E)=\frac{2\pi }{(2\pi i)^{3/2}}\underset{\mathrm{paths}}{}\frac{1}{\sqrt{A}}\mathrm{exp}\left[iS(𝐫,𝐫^{})i\pi \mu /2\right],$$
(7)
where $`S(𝐫,𝐫^{})`$ is the action and $`\mu `$ is the Maslov index for the path. The stability coefficient $`A`$ is given by
$$A=\frac{𝐱_f}{𝐩_i}=\underset{\mathrm{\Delta }𝐩_i0}{lim}\frac{\mathrm{\Delta }𝐱_f}{\mathrm{\Delta }𝐩_i},$$
(8)
using the definitions from Fig. 13, and where $`\mathrm{\Delta }𝐱_f`$ is the component of $`\mathrm{\Delta }𝐱_f`$ that is perpendicular to $`𝐩_{1f}`$.
The coefficient $`A`$ describes the stability of trajectories beginning from a particular point in phase space. If $`A`$ is large, then small changes in the initial direction of the trajectory lead to large displacements in the final positions. More precisely, if the distance $`\mathrm{\Delta }𝐱_f`$ grows exponentially with the length of the trajectories, we say that the trajectories are chaotic. If $`\mathrm{\Delta }𝐱_f`$ grows only linearly, then they are stable.
For the case of our resonator in the unstable regime, only one type of orbit enters into the sum in Eq. (7): the horizontal orbit. Therefore, in order to find the contributions of the geometric orbits to the transmission spectrum, we need only the actions ($`S=kl`$, where $`l`$ is the length of the orbit) and stability coefficients $`A`$ for the primary horizontal orbit and its repetitions. In addition, we need to keep track of the Maslov index for each orbit. A series of such orbits, together with the associated Maslov indices, is shown schematically in Fig. 14(a).
In this figure, the QPC/wall is located at the lower part of each diagram, while the reflector is located at the upper part. Each upward (downward) sloping line segment represents part of a trajectory from the QPC (reflector) to the reflector (QPC). Maslov indices of $`\mu =2`$ are indicated for points where the wave is reflected at the wall or arc-reflector, and an index of $`\mu =1`$ is acquired each time the ray passes though the focus on its return from the reflector toward the wall.
Figure 15 shows a transmission spectrum for the resonator in the unstable regime, computed using only the horizontal orbit.
For this calculation, the sum in Eq. (7) was cut off after the 20th term, that is, orbits of up to 20 round trips were included in the sum. The half-wavelength periodicity of the spectrum is clearly seen in the figure. Upon comparison with Fig. 3, we see that the peak positions match very well with the experimentally measured spectrum. Note, however, the absence of the peaks corresponding to diffractive orbits which are present in Fig. 3. It is to the calculation of these diffractive resonances that we now turn.
### B Diffractive orbits in the semiclassical propagator
So far, the Green function in Eq. (7) includes paths that undergo evolution under the free-particle Hamiltonian, including bounces off the wall and mirror. In order to calculate the spectral properties of the resonator semiclassically, diffractive orbits must be included into the sum over paths that forms the semiclassical propagator. This problem has been studied by a number of authors . In the literature, much attention has been focused on finding the effects of diffraction on the spectra of closed systems. However, in closed systems, diffractive orbits generally play a minor role because they are overwhelmed by the huge number of unstable periodic orbits that do not involve diffraction. In this section, these methods will be extended to include open systems.
For our purposes, it is sufficient to include diffraction at the level of a single diffraction event per orbit. Although multiply diffracted orbits strictly belong in the semiclassical sum, in practice they can be safely neglected. This is because of the amplitude of an incident ray on the reflector tip is subsequently sprayed in all directions, so that only a small part returns in a direction that eventually leads it back to the QPC. Therefore, we will consider only singly diffracted orbits. The Green function is the product of three amplitudes: the first for going from the starting point to the point of diffraction, the second for the diffraction event itself, and the third for going from the diffraction point to the final point, as follows:
$`G_{\mathrm{diff}}(𝐫,𝐫^{},E)`$ $`=`$ $`G_{\mathrm{sc}}(𝐫,𝐫_d,E)D(l_1,l_2,\theta _1,\theta _2,E)G_{\mathrm{sc}}(𝐫_d,𝐫^{},E)`$ (9)
$`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle \underset{\mathrm{paths}}{}}{\displaystyle \frac{D(l_1,l_2,\theta _1,\theta _2,E)}{\sqrt{A_1A_2}}}\mathrm{cos}\varphi _1\mathrm{cos}\varphi _2\mathrm{exp}\left[i(S_1+S_1)i\pi (\mu _1+\mu _2)/2\right],`$ (10)
where the diffraction event occurs at position $`𝐫_d`$, and the total path from $`𝐫`$ to $`𝐫^{}`$ is made up of two legs, one of length $`l_1`$, with stability coefficient $`A_1`$, and the other of length $`l_2`$, with stability coefficient $`A_2`$. Each of these legs has an action $`S_1=kl_1`$, $`S_2=kl_2`$ and Maslov index $`\mu _1`$, $`\mu _2`$ associated with it. The factors $`\mathrm{cos}\varphi _1`$, $`\mathrm{cos}\varphi _2`$ represent the coupling of each leg to the QPC, and will be discussed in the next section. The various parameters are illustrated in Fig. 16 for one of the shorter orbits, corresponding to the third term in Fig. 14(b). In that figure, the first few terms entering into Eq. (9) are shown, where diffraction events are represented by filled circles.
The diffraction coefficient $`D(l_1,l_2,\theta _1,\theta _2,E)`$ depends on the lengths of each leg as well as the angles that the incident and diffracted ray make with the surface at the tip of the obstacle. These various lengths and angles are shown in Fig. 16. The diffraction coefficient in the sum above is similar to that appearing in our study of the half-line, differing only in the argument of the Fresnel integral. It is given by,
$`D(\theta _1,\theta _2,l_1,l_2,k)`$ $`=`$ $`\text{sgn}(a_i)K\left(|a_i|\sqrt{{\displaystyle \frac{kl_1l_2}{l_1+l_2}}}\right)+\text{sgn}(a_r)K\left(|a_r|\sqrt{{\displaystyle \frac{kl_1l_2}{l_1+l_2}}}\right),`$ (11)
$`a_{i,r}`$ $`=`$ $`\sqrt{2}\mathrm{cos}\left({\displaystyle \frac{\theta _1\theta _2}{2}}\right),`$ (12)
and $`K(x)`$ is the modified Fresnel integral defined in Eq. (5).
The effect of the diffractive terms in the sum is shown in Fig. 17. For this calculation, all orbits with up to 20 round trips between the wall and mirror and zero or one diffractive events were included in the sum.
We see that the effect of the diffractive terms is to modulate the geometric result, with new peaks appearing to the right of the geometric peaks. The theoretical curve appearing in Fig. 17 is overlaid with the experimental data in Fig. 3; the agreement between theory and experiment is quite good, both in the peak positions and widths of the geometric and diffractive peaks. We emphasize here that the semiclassical prediction breaks down for reflector/wall separations near $`D=R`$, because there the focus approaches the point where the Green function is evaluated, and the semiclassical prediction diverges. This is the reason for the incorrect, large transmission calculated at $`D=R`$.
In Fig. 18, we plot the semiclassically calculated transmission of the resonator in the unstable regime versus both reflector/wall separation and wavelength. The parameters are identical to those of the experimental data shown in Fig. 7. The separation of the diffractive peaks from the Fabry-Perot peaks with increasing reflector/wall separation is quite clear. Two diffractive peaks per geometric peak are visible. The half-wavelength periodicity is also apparent. The similarity between theory and experiment in these two figures is striking.
### C Semiclassical wavefunctions
In this section, we describe the procedure for including diffraction into a semiclassical calculation of the resonator wavefunctions. As in the previous sections, we can split up the sum over paths into a geometric part and a diffractive part. The contribution to $`\psi `$ from the geometric orbits is proportional to the semiclassical amplitude for getting from the QPC to the point of interest:
$$\psi _{\mathrm{geo}}(𝐫)G_{\mathrm{sc}}(𝐫_{\mathrm{QPC}},𝐫,E),$$
(13)
where $`𝐫_{\mathrm{QPC}}`$ is the location of the QPC, and $`𝐫`$ is the location of interest. This Green function is the same as we have already encountered in the semiclassical expression for the transmission neglecting diffraction, appearing in Eqs. 6 and 7. The only difference is that now we are looking at an off-diagonal element of the Green function, rather than a diagonal element. Some of the shorter trajectories that are included in the sum are shown in Fig. 19(a). All such trajectories begin at the QPC and end at the point $`𝐫`$, undergoing specular reflection at the wall and reflector.
To include diffraction, we simply add diffractive terms to the sum, much as we did to incorporate the effects of diffraction in the formula for the Green function in Eq. 9. We have
$$\psi (𝐫)G_{\mathrm{sc}}(𝐫_{\mathrm{QPC}},𝐫,E)+G_{\mathrm{sc}}(𝐫_{\mathrm{QPC}},𝐫_{\mathrm{tip}},E)D(l_1,l_2,\theta _1,\theta _2,E)G_{\mathrm{sc}}(𝐫_{\mathrm{tip}},𝐫,E),$$
(14)
where the diffraction coefficient $`D(l_1,l_2,\theta _1,\theta _2,E)`$ is identical to that appearing in Eq. 11. The lengths $`l_1,l_2`$ and angles $`\theta _1,\theta _2`$ are defined the same way as in Fig. 16, with the obvious difference that the final point of the trajectory is no longer the location of the QPC, but the location of interest, $`𝐫`$. In Fig. 20
we show the result of a semiclassical calculation for the wavefunctions pictured in Fig. 6. For comparison, we also show the result when diffractive paths are left out of the sum. In the calculation, all singly diffracted paths involving up to 20 bounces were included.
We saw earlier that diffraction has a large effect on the transmission spectrum, which in turn could be derived from the value of the wavefunction near the QPC. By contrast, we see now that the inclusion of diffraction has a relatively minor effect on the overall appearance of the wavefunctions. The explanation for this apparent paradox is that the QPC is near many reflection boundaries, where the diffractive corrections are especially large.
The transmission through a small aperture (here the QPC) is extremely sensitive to small returning bits of amplitude, which interfere coherently with the wave entering through the aperture to modulate the transmission. In our case, a major source of returning amplitude is provided by the diffractive orbits. This is analogous to scanning tunneling microscope “quantum corral” imaging , where the tunneling from the tip plays the role of the quantum point contact, and reflections from atoms and impurities represent the diffraction and geometric scattering off the ends of the reflector. The modulation of the transmission through the QPC by the presence of the reflector can be understood in terms of a small returning wave amplitude beating against a much stronger “nascent” amplitude coming out of the QPC: if the amplitude from the QPC in absence of the reflector is $`A`$, and the returning amplitude with the reflector present is $`a`$, then the total amplitude at the QPC is simply $`A+a`$. The transmission will be proportional to the square of this amplitude,
$$T|A+a|^2=|A|^2+aA^{}+a^{}A+|a|^2.$$
(15)
The two interference terms above, linear in $`a`$, are responsible for all the structure in the transmission when the reflector is present in the cavity.
## VII Semiclassics in the time domain
Further evidence of diffractive orbits in the transmission spectrum can be obtained by analyzing the spectrum in the time domain. The two representations are related by the Fourier transform:
$$g(t)=_{\mathrm{}}^{\mathrm{}}S_{11}(\omega )e^{i\omega t}𝑑t.$$
(16)
Here $`g(t)`$ represents the amplitude for a pulse launched from the QPC at time $`t=0`$ to return at time $`t`$. That is, if a short pulse were emitted from the antenna at time $`t=0`$, echos would return to the antenna at certain later times. These echos are indicated by peaks in the return spectrum. In Fig. 21, we plot the amplitude $`|g(t)|`$ as a function of time for the resonator in the stable and unstable regimes. In these plots, time has been normalized so that one unit is the time it takes for light to travel a distance equal to the radius of curvature of the reflector.
In the stable regime, the echos persist for hundreds of bounces, indicating that indeed the dynamics is stable in this regime. In this regime, the lifetime of the states is limited by resistive losses in the copper plates of the resonator, which was quite small: typical quality factors of the resonances in this regime were $`Q3000`$. However, in the unstable regime, the echos are significantly reduced in amplitude after only a few returns. For the first few return peaks, $`|g(t)|^2`$ decays exponentially with a decay constant given approximately by the Lyapunov exponent for the horizontal orbit. Later peaks, where diffractive orbits are more important, also decay exponentially but with a different decay constant.
In Fig. 22, we show an expanded view of some of the return peaks shown in Fig. 21(b). Of importance here is the splitting of the return peaks which is visible on echos 5-9. This splitting is due to the coexistence of orbits with slightly different periods. The longer of these orbits is just the horizontal orbit, which appears on the right of each group. The left peak of each group is made up of a family of diffractive orbits of nearly the same length. We have done a quantitative study of the lengths of the closed orbits and find excellent agreement with the observed splitting. The calculated lengths of the orbits appear in the plot as vertical bars above the peaks. The horizontal orbit length is marked with a longer bar. The lengths of all the orbits in units of the radius of curvature appear in Table. I. The presence of this splitting in the return spectrum is strong evidence in support of the claim that diffraction off the edges of the reflector supports other closed orbits, which lead to resonances in the transmission spectra. Note that for the long orbits, the diffractive peaks are even stronger than the peaks from the geometric orbit. This is because the number of diffractive orbits increases linearly with the length of the orbit, whereas there is always only one geometric orbit, regardless of length.
## VIII The avoided crossing
We now turn our attention to the avoided crossing that appears in Fig. 7 near $`(DR)/R=0`$. We mentioned earlier that the level repulsion at this point is in part mediated by diffraction. To show this, we compare our open resonator with three closely-related closed systems.
In the stable regime, imagine closing the QPC aperture and increasing the open angle of the arc-shaped reflector until it touches the straight wall. The resulting shape is one-half of a lemon billiard (see Fig. 23(a)). The shape of the lemon billiard is determined by the parameter $`x_0(DR)/R<0`$, which is the $`x`$-value of the center of curvature of the arc on the right in units of the radius of curvature. The lemon billiard has been studied before ; for our purposes it is important that the classical dynamics in the lemon billiard is dominated by a large regular region.
As $`x_0`$ is increased to $`0`$, the area enclosed by the arc becomes half of a perfect circle, and the closed system becomes completely integrable (see Fig. 23(b)). The eigenstates of a circular billiard are $`J`$-type Bessel functions.
Finally, as $`x_0`$ is made positive, the semicircular arc can be extended with horizontal straight segments to form half of a stadium billiard. The classical dynamics in the stadium billiard is completely chaotic. See Fig. 23(c).
The eigenstates of these three systems share certain properties with the scattering states of our open system. In these closed systems, there are no wall ends so we expect diffraction to play a smaller role in the energies and wavefunctions.
The similarity between these closed systems and our resonator is greatest in the stable regime. In this regime, the states of the resonator are essentially the same as the states of the lemon billiard that are even about the $`x`$-axis, and odd about the $`y`$-axis. (The straight wall in the resonator lies on the $`y`$-axis, thus enforcing a node there, while the symmetric position of the antenna means that it only excites states that are even about the $`x`$-axis.) As the reflector position is varied, we correspondingly change the parameter $`x_0<0`$ of the lemon billiard, and compare the states of each system. In the experiment, we see an avoided crossing at $`x_0=0`$. To investigate this, we can follow the states of the lemon/stadium billiard as the parameter $`x_0`$ is swept slowly through zero.
In Fig. 24(a)
we plot the (even-about-$`x`$, odd-about-$`y`$) levels of the lemon/stadium billiard as a function of the parameter $`x_0=(DR)/R`$ and wavelength. On the right, the corresponding experimental transmission spectra for the open system are shown. Notice the exact matching of the transmission peaks with the lemon billiard states in the stable regime. This excellent agreement is an artifact of the fact that, while our resonator is a geometrically *open* system, it is classically *closed*, in the sense that almost all trajectories beginning at the QPC that hit the reflector are doomed to forever remain in the region between the wall and the reflector. We only begin to notice the “openness” of our resonator for the peaks corresponding to large numbers of angular nodes. As we shall see, such states have a large angular spread, and it grows with decreasing $`x_0`$, so that eventually the caustics of the classical orbits supporting these states approach the tip of the reflector. When this happens, the states broaden and disappear. For example, this is apparent for the peak with quantum numbers (13,8) at around $`x_0=0.1`$. The wavefunction corresponding to this peak has 13 radial nodes between the $`y`$-axis and the reflector, and 8 angular nodes. The states with fewer angular nodes do not have such a wide angular extent, and therefore they do not disappear from the spectra until the reflector is farther from the wall.
It is apparent that there are many states in the closed lemon billiard that do not appear at all in the analogous open system. All such states have considerable amplitude near the walls of the lemon billiard that are ‘missing’ from the open resonator—therefore in the open system the analogous states escape the resonator immediately. Thus they are absent in the experimental transmission spectra.
At $`x_0=0`$, there is an avoided crossing in both the open and closed systems. However, we see that the level repulsion is stronger in the open system. For example, observe the level repulsion of the states (17,0) and (16,2). The distance of closest approach of these two levels is four times greater in the open system than in the closed system. (For the open system, we judge the distance between the “levels” as the distance between the peak maxima.) Now, the major difference between the two systems as far as the eigenstates are concerned is the presence of diffractive orbits in the open case. Semiclassically speaking, there are diffractive paths that go from one state to the other and introduce a coupling that increases the level splitting. Therefore, we believe that the avoided crossing is in large part mediated by *diffraction* in the open system.
For $`x_0>0`$, the correspondence between the closed and open systems comes to a somewhat abrupt halt. This is apparent in Fig. 24. The reason is as follows. In the closed system, orbits with times up to about the Heisenberg time, $`t_H1/\mathrm{\Delta }E`$, affect the properties of the quantum states. In the open system, on the other hand, only a few orbits from a very specific part of phase space stay in the system for more than a few bounces. Therefore, we can expect states of the closed and open systems to correspond only if a significant subset of trajectories remains in the non-escaping region of phase space for roughly the Heisenberg time. Trajectories that leave that region of phase space explore parts of phase space where the systems differ and the correspondence breaks down. For the circle billiard at our experimental parameters, the particle makes approximately $`7`$ horizontal bounces within the Heisenberg time of the quarter-stadium.
Using the 2D density of states, we can estimate the range of $`x_0`$ for which we may expect rough correspondence between the eigenstates of the closed system with the resonances of the open system in the unstable regime. There is another characteristic time, $`t_\lambda `$, which we will call the Lyapunov time, which is the characteristic time that it takes for a trajectory to “fall off” the horizontal periodic orbit. We expect correspondence when $`t_Ht_\lambda `$. The Lyapunov time is
$$t_\lambda =\frac{L}{k\lambda _{\mathrm{Lyap}}},$$
(17)
where $`L`$ is the length of the periodic orbit, $`\lambda _{\mathrm{Lyap}}`$ is the Lyapunov exponent of the orbit, and $`k`$ is the wavenumber ($`L/k`$ is the period of the orbit). We will focus on the horizontal orbit for the case where the system is just barely unstable, so that $`L2R`$. The Lyapunov exponent for the horizontal orbit is given by the logarithm of the largest eigenvalue of the monodromy matrix, linearized about the horizontal orbit,
$$M=\left(\begin{array}{cc}1+2x_0& 2x_0(1+x_0)\\ 2& 1+2x_0\end{array}\right)$$
(18)
. The largest eigenvalue of this matrix is given by
$`m_+`$ $`=`$ $`1+2x_0+2\sqrt{x_0(1+x_0)}`$ (19)
$``$ $`1+2\sqrt{x_0},`$ (20)
for small $`x_0`$. Thus the Lyapunov exponent is
$`\lambda _{\mathrm{Lyap}}`$ $``$ $`\mathrm{ln}(1+2\sqrt{x_0})`$ (21)
$``$ $`2\sqrt{x_0}.`$ (22)
Therefore, the Lyapunov time for the horizontal orbit is approximately
$$t_\lambda \frac{R}{k\sqrt{x_0}}.$$
(23)
This must be compared with the Heisenberg time, which is $`t_H=A/2\pi =R^2/8`$ (the effective area of the billiard is $`A=\pi R^2/4`$ because of symmetry—we only take one quarter of the area of a full circle). We expect the correspondence between the open and closed systems to break down when these times are of the same order. That is, we look for the value of $`x_0`$ for which $`t_Ht_\lambda `$, which is
$$x_0\left(\frac{8}{kR}\right)^2.$$
(24)
Now, in our energy range $`kR18\pi `$. Thus we have for the value of $`x_0`$ at which correspondence ceases to hold,
$$x_0\left(\frac{8}{18\pi }\right)^20.02\frac{\lambda }{R}.$$
(25)
As indicated, this value of $`x_0`$ is much less than a wavelength, scaled to the radius of the billiard. This means that the states of the closed billiard already have mixed through a number of avoided crossings by the time the circle has been “stretched” by one wavelength. Indeed in Fig. 24 we see that the first avoided crossings are happening already near $`x_0=0.02`$, consistent with our rough estimate.
A sampling of states for the closed system of the lemon/stadium billiard is shown in Fig. 25. In this figure, the states corresponding to the series of quantum numbers $`(17,0)`$ to $`(13,8)`$ are shown, as the transition from a lemon to a stadium is made. The corresponding values of $`k`$ for each state shown are given in Table II. The value of $`x_0`$ is given at the top of each column. We see that the angular extent of the lemon billiard states increases with the number of angular nodes, as mentioned previously. As $`x_0`$ approaches zero, the angular extent of each state is fully developed and covers the entire billiard, as is required by the rotational symmetry of the circle billiard states. The transition from the first to the second column is diabatic, that is, from left to right we track the states with the same character.
In the third column, we have plotted the states after $`x_0`$ has been adiabatically increased to the small value $`x_0=0.02`$, such that they still bear some resemblance to the states of the circle. This parameter value was chosen to be just before the first avoided crossing in the closed system, in accordance with the discussion above. This effectively means that the periodic orbits are not too different from those of a circle billiard; i.e., the periodic orbits are only weakly unstable. This means that one may still draw analogies between the closed and open systems in this regime, although the character of a particular resonance in the open system may be shared by several states in the corresponding closed system. For example, note that the first two states in the third column have angular lobes which are qualitatively similar to the diffractive state shown in Fig. 6. On the other hand, the third and fourth states in that column have amplitude running from the center out along the $`x`$-axis, qualitatively similar to the Fabry-Perot type state in Fig. 6. There is even the appearance of a focus just to the right and left of the center of the billiard in these two states, analogous to the focus seen in the Fabry-Perot state for the open system.
In the last column, we plot the five wave functions after the parameter $`x_0`$ has been adiabatically increased to $`x_0=0.2`$, or 20% of the radius of curvature. Now we see that the states no longer have anything to do with the states of the open system, exactly for the reasons given above.
## IX Imaging wavefunctions with a coarse probe
We end this paper with an interesting discovery which was made in the course of carrying out the measurements described above: namely, the possibility of measuring an electronic wavefunction in a 2DEG. The viability of measuring pure state wavefunctions in the context of microwave billiard systems has been demonstrated by many authors . However, the imaging of a wavefunction in a real quantum system has not yet been achieved.
We believe a technique similar to the method used here could be applied in clean mesoscopic systems to obtain images of two dimensional wavefunctions in a 2DEG. In such systems, the fermi wavelength $`\lambda `$ of the electrons is on the order of 50 Å. In analogy to the steel ball used in the microwave experiments, an AFM tip held close to the surface of the heterostructure could serve as a perturbation to measure the frequency shifts. However, the perturbation due to a nearby AFM tip on the electron wavefunction is a smooth potential disturbance 50-100 Å in size, comparable to or greater than $`\lambda `$. One might expect that it would be difficult to measure the nodal lines of a wavefunction with such a coarse probe. To investigate this problem, we tried measuring a microwave mode with a coarse probe. Instead of the small steel ball, we used a probe 8 cm in size, corresponding to 1.5 wavelengths, for these measurements. The probe is shown in Fig. 26. It approximately a “chord” of a sphere. The probe could be moved around the upper surface of the top plate of the cavity by means of an external magnet, in a manner similar to that used for the steel ball perturbation. For a discussion of the form of the perturbation caused by this probe, see Appendix B.
The results of a wavefunction measurement of the peak d1 with the coarse probe are displayed in Fig. 27.
Only half of the wavefunction was measured. For comparison, the same state as measured by the 4 mm ball is shown as well. Note that far from the QPC, the nodal lines in the two measurements match very well. For the coarse measurement, data near the wall and mirror were unavailable because of the large size of the probe. The level of detail obtained with the coarse probe is surprising, considering that it was over three half-wavelengths in diameter. This result suggests it might be possible to directly image an electron wavefunction in a 2DEG. A perturbative analysis of the effect of the coarse probe is presented in Appendix B.
## X Conclusion
In summary, we have demonstrated the existence of diffractive orbits in an open microwave billiard, which give rise to resonances and wavefunctions that would not be predicted by a simple semiclassical theory. Such orbits are of importance in open, unstable systems where the number of unstable classical periodic orbits is small. In such systems, diffraction can play a major role in determining the spectrum of the system. Furthermore, we have shown that it may be possible to measure the pure state wavefunction of an electron in a 2DEG, by using a a coarse AFM tip as a probe.
This work was supported by NSF Grant CHE9610501. We are grateful to the Hewlett Packard Corporation for the loan of a network analyzer that was used in these experiments. We thank J. D. Edwards for the computer program that generated the exact quantum mechanical wavefunctions.
## A Maslov indices
In the experiment, the antenna is placed very close to, but not exactly at, the wall. This means there are in fact four orbits associated with each single orbit in the billiard when the source is placed exactly at the wall. The four orbits associated with the shortest horizontal orbit are shown in Fig. 28.
Each orbit in this family has a different total Maslov index and length, summarized in Table III.
We can find the effect of grouping these four orbits into a single orbit by computing the following sum,
$$Ae^{iS_{\mathrm{eff}}i\pi \mu _{\mathrm{eff}}/2}=\underset{n=\alpha ,\beta ,\gamma ,\delta }{}e^{iS_ni\pi \mu _n/2},$$
(A1)
where the sum is over the four orbits as $`d0`$. The coefficient $`A`$ will be found by doing the sum. We have
$`Ae^{iS_{\mathrm{eff}}i\pi \mu _{\mathrm{eff}}/2}`$ $`=`$ $`e^{2ik(Dd)3\pi i/2}+2e^{2ikD5\pi i/2}`$ (A3)
$`+e^{2ik(D+d)7\pi i/2}`$
$`=`$ $`e^{2ikD5\pi i/2}\left(2e^{2ikd}e^{2ikd}\right)`$ (A4)
$`=`$ $`e^{2ikD5\pi i/2}\left(22\mathrm{cos}(2kd)\right)`$ (A5)
$``$ $`\left(2kd\right)^2e^{2ikD5\pi i/2}.`$ (A6)
Thus we have that $`A=(2kd)^2`$, $`\mu _{\mathrm{eff}}=5`$, and $`S_{\mathrm{eff}}=2kD`$, as expected.
A similar analysis on orbits coming to the antenna at an angle shows that the coupling to the antenna varies as $`\mathrm{cos}\varphi `$, as stated in Eq. (9). Referring to Fig. 29,
a source of rays begins at the point $`P`$, and two of the rays find their way to the antenna. The direct path has length $`L_1`$ and Maslov index $`\mu _1=0`$. The second path first bounces off the wall before arriving at the antenna, and has length $`L_2`$ and Maslov index $`\mu _2=1`$. We want to combine these two paths into a single trajectory, by doing a sum over the two trajectories as was done above, in order to find the effective action and Maslov index for the trajectory. We further define the length $`L_0`$, which is the distance between $`P`$ and the intersection of the wall with the axis of symmetry. The two lengths $`L_1`$ and $`L_2`$ are given by
$`L_{1,2}`$ $`=`$ $`\sqrt{\left(L_0\mathrm{cos}\varphi d\right)^2+\left(L_0\mathrm{sin}\varphi \right)^2}`$ (A7)
$``$ $`L_0d\mathrm{cos}\varphi ,`$ (A8)
where $`\varphi `$ is the angle between the trajectory $`L`$ and the $`x`$-axis, and $``$ and $`+`$ correspond to trajectories 1 and 2, respectively. The sum over orbits is then
$`e^{iS_{\mathrm{eff}}i\pi \mu _{\mathrm{eff}}/2}`$ $`=`$ $`e^{iS_1}+e^{iS_2i\pi }`$ (A9)
$`=`$ $`2\mathrm{cos}\left(kd\mathrm{cos}\varphi i\pi /2\right)e^{ikL_0i\pi /2}`$ (A10)
$``$ $`2kd\mathrm{cos}\varphi e^{ikL_0i\pi /2}.`$ (A11)
As expected, the effective action for the two paths is just $`S_{\mathrm{eff}}=kL_0`$, and the Maslov index is $`\mu _{\mathrm{eff}}=1`$. The factor $`\mathrm{cos}\varphi `$ appearing above is exactly the angular dependent coupling coefficient appearing in Eq. (9). This is just the simple angular dependence of a $`p`$-wave point source.
## B Effect of a coarse probe
In this section we derive an effective potential describing the perturbation to a TEM microwave resonance due to a coarse probe. For simplicity we consider the case of a closed cavity. The electric field $`(𝐫,t)=Re𝐄(𝐫)e^{i\omega t}`$ and magnetic field $`(𝐫,t)=Re𝐁(𝐫)e^{i\omega t}`$ of such a state can be expressed in terms of the quantum analogue state, $`\psi (x,y)`$, as follows:
$`𝐄(x,y,z)`$ $`=`$ $`{\displaystyle \frac{1}{h^{1/2}}}\psi (x,y)\widehat{𝐳}`$ (B1)
$`𝐁(x,y,z)`$ $`=`$ $`{\displaystyle \frac{1}{h^{1/2}}}{\displaystyle \frac{1}{ik}}\left[{\displaystyle \frac{\psi }{y}}\widehat{𝐱}{\displaystyle \frac{\psi }{x}}\widehat{𝐲}\right].`$ (B2)
The height of the cavity in the $`z`$-direction is denoted by $`h`$. (We use the normalization conventions $`𝑑x𝑑y\psi ^2=d^3𝐫|𝐄|^2=d^3𝐫|𝐁|^2=1`$.)
The perturbation to the frequency of a microwave resonance by a conducting probe is
$`{\displaystyle \frac{\omega ^2\omega _0^2}{\omega ^2}}`$ $`=`$ $`{\displaystyle _{\mathrm{\Delta }V}}d^3𝐫\left(|𝐁(𝐫)|^2|𝐄(𝐫)|^2\right)`$ (B3)
$`=`$ $`{\displaystyle \frac{1}{h}}{\displaystyle 𝑑x𝑑y\delta h\left[\frac{1}{k^2}(\psi )^2\psi ^2\right]},`$ (B4)
where $`\omega _0`$ is the unperturbed frequency and $`\delta h(x,y)`$ is the thickness of the probe. The volume integral is over the volume excluded by the probe; the second line follows because the fields have no $`z`$-dependence. From the chain rule and the Helmholtz equation,
$`\delta h(\psi )^2`$ $`=`$ $`\left[\delta h\psi \psi \frac{1}{2}(\delta h)\psi ^2\right]`$ (B6)
$`+k^2\delta h\psi ^2+\frac{1}{2}^2(\delta h)\psi ^2.`$
Substitute Eq. (B6) into Eq. (B4). Since $`\delta h`$ is localized away from the lateral boundaries of the cavity, the surface terms of Eq. (B4) vanish and we are left with
$$\frac{\omega ^2\omega _0^2}{\omega ^2}=\frac{1}{k^2}𝑑x𝑑y\frac{1}{2}\frac{^2(\delta h)}{h}\psi ^2.$$
(B7)
Therefore, to first order in perturbation theory, the coarse probe can be thought of as an effective perturbing potential of the form
$$V_{\text{eff}}(x,y)^2\frac{\delta h(x,y)}{h}.$$
(B8)
Note that Eq. (B8) implies that the area integral of $`V_{\text{eff}}`$ vanishes, which is a significant constraint on the type of potential that can be modeled by a conducting probe in a microwave cavity. In particular, the perturbation of a 2DEG by an AFM tip would not have this property.
In the case of our probe, $`\delta h/hf[1(\rho /\rho _0)^2]`$ where $`\rho `$ is the distance from the center of the probe, $`\rho _0=40\text{ mm}`$, and $`f=0.4`$, so
$$V_{\text{eff}}\frac{4f}{\rho _0^2}+\frac{2f}{\rho _0}\delta (\rho \rho _0);$$
(B9)
in other words, the probe’s effective potential is a flat well surrounded by a repulsive ring at the probe’s circumference.
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# 1 Introduction
## 1 Introduction
The space-time structure at Planck scales $`L_P(G\mathrm{}/c^3)^{1/2}`$ will be drastically affected by the quantum gravitational effects and it is generally believed that $`L_P`$ acts as a physical cutoff for space-time intervals. The two main approaches to quantum gravity, Superstring theory and Loop quantum gravity, incorporates the fundamental length scale by considering extended structures, rather than point particles, as fundamental blocks.
The existence of a fundamental length implies that processes involving energies higher than Planck energies will be suppressed, and the ultraviolet behavior of the theory will be improved. This could arise naturally in String theories; several other models also incorporate a Planck length cut-off in a suitable manner to improve the ultra violet behaviour of the theory <sup>?</sup>. One direct consequence of such improved behavior will be that the Feynman propagator(in momentum space) will acquire damping factor for energies larger than Planck energy. However, this propagator – which arises in the standard formulation of quantum field theory – does not take into account of the existence of any fundamental length in the space-time. On the other hand, such a fundamental length scale was introduced into the Feynman propagator in a Lorentz invariant manner by invoking the “principle of path-integral duality” <sup>?</sup>. According to this postulate, the weightage given for a path in the path integral should be invariant under the transformation $`L_P^2/`$, where $``$ is the length of the path and the fundamental length scale $`L_P`$ is assumed to be of the order of the Planck length $`(G\mathrm{}/c^3)^{1/2}`$.\[In this paper, when we say that the fundamental length is $`L_P`$, we actually mean that it is of $`𝒪(L_P)`$.\] Padmanabhan <sup>?</sup> has shown by rigorous evaluation of the path integral by lattice techniques that the effect of the duality principle is to modify the weightage given to a path of proper time $`s`$ from $`\mathrm{exp}(im^2s)`$ to $`\mathrm{exp}[i(m^2sL_P^2/s)]`$, where $`m`$ is associated with the mass of the particle. For example, the Feynman propagator for a free scalar field of mass $`m`$, propagating in the flat $`(3+1)`$ space-time, in Schwinger’s proper time formalism, is described by the integral
$$G_F(x,x^{})=\frac{1}{(4\pi i)^2}_0^{\mathrm{}}\frac{ds}{s^2}\mathrm{exp}(im^2s)\mathrm{exp}[i(xx^{})^2/4s].$$
(1)
The duality principle modifies the Feynman propagator to the form
$`G_F^P(x,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi i)^2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s^2}}\mathrm{exp}(im^2s)\mathrm{exp}[i((xx^{})^2+L_P^2)/4s]`$ (2)
$`=`$ $`{\displaystyle \frac{m}{4\pi ^2}}{\displaystyle \frac{K_1(im\sqrt{x^2L_P^2iϵ})}{\sqrt{x^2L_P^2iϵ}}},`$
where $`K_1(z)`$ is the modified Bessel function of order 1. \[The metric signature we follow is $`(+,,,)`$.\] In momentum space, the modified propagator is
$$G_F^P(p)=i_0^{\mathrm{}}𝑑z\mathrm{exp}(iL_P^2/4z+i(p^2m^2+iϵ)z).$$
(3)
The presence of a fundamental length scale is a feature that is expected to arise in a quantum theory of gravity. Hence, the modification of the weightage factor as mentioned above can be interpreted as being equivalent to introducing quantum gravitational corrections into standard field theory. The ultraviolet divergences in quantum field theory arise from the singularities of the propagator functions on the light cone, and a smearing out of the light cone due to the quantum gravitational corrections, using the principle of “path integral duality”, will lead to the suppression of these divergences.
In an earlier work <sup>?</sup>, the implications of the modified propagator to certain conventional non-perturbative quantum field theoretic results were discussed in detail. It was found that the essential feature of this prescription of path integral duality is to provide an ultra-violet cutoff at the Planck energy scales, thereby obtaining a Lorentz invariant finite results. Encouraged by this fact, we would like to estimate the renormalization factors in QED, and other radiative correction terms which cannot be obtained in the conventional QED calculations.
The standard definition of Feynman propagator for Dirac particles is
$$S(x)=(i\gamma ^\mu _\mu +m)G_F(x),$$
(4)
where, $`G_F(x)`$ is the usual Feynman propagator for the scalar particles. In our analysis of evaluating quantum gravitational corrections to the standard QED calculations, we assume that the effect of “principle of path integral duality” on the Dirac propagator is defined as
$$S^P(x)=(i\gamma ^\mu _\mu +m)G_F^P(x).$$
(5)
Hence, the effect of summing over the quantum fluctuations of the space-time structure, in the low energy scales, can be realized in any field theoretic calculations where the propagator of the scalar particles appear explicitly.
There is a similar approach in the literature by Ohanian <sup>?</sup>, who had used a smeared propagator which is Poincare invariant, to calculate the radiative corrections in QED. The modified Feynman propagator in our case is obtained by rigorous evaluation and the duality principle introduces the fundamental length scale in a Lorentz invariant manner. We also comment on the differences between our approaches and results at the appropriate sections below.
Before proceeding to the technical aspects, it is necessary to outline certain conceptual issues related to this approach. In the conventional approaches to quantizing a field based on some classical Lagrangian, one will invariably obtain quantum corrections to the classical Lagrangian. When these corrections are local in the configuration space and contain terms which are of the same form as those in the classical Lagrangian, it is necessary to absorb them into the parameters of the original Lagrangian. This process of renormalization, has — a priori — nothing to do with divergences. Usually, however, the quantum corrections to the theory — calculated by perturbative methods — lead to divergent expressions. If the new (divergent) terms are not of the form of the original terms, then the theory cannot be interpreted perturbatively. However, if the divergent terms have the same structure as the terms in the original Lagrangian, it is possible to use the procedure of renormalization (which, a priori, has nothing to do with divergences) to give meaning to the theory. To do this properly, it is necessary to first evolve a procedure (called regularization) which allows the divergent expressions to be recast as the limit of some finite quantities. After performing the renormalization subtractions, one is free to take the required limit, leading to finite corrections.
The above approach gets modified in two essential aspects when the modified propagator is used. Firstly, when quantum gravitational corrections lead to a cutoff, the quantum corrections will not have any divergent terms; that is, the regularization is now built into the theory with Planck length acting as regulator. But renormalization of the theory is still needed and the physical and bare coupling constants will differ by a finite amount. Secondly, the regularization procedure is now fixed by our ansatz. This is important because, it is well known from standard work in quantum field theory that different regularization procedures are not equivalent. For example, dimensional regularization and momentum space techniques are not completely equivalent as regards their treatment of the symmetries of the system. Since we have no freedom in choosing a regularization, it is necessary to accept and investigate the final results arising from the ansatz. The use of modified propagator to the field theoretic calculation is a more realistic approach towards the removal of the divergences in field theory based on the relevant physics, rather than a formalistic approach based on improper mathematical manipulations. In an earlier work <sup>?</sup>, it was shown that the physical parameters in this system like mass and coupling constant have additive finite terms, that are proportional to the powers of $`L_P`$ by calculating the effective potential for a self-interacting scalar field theory using the modified propagator.
In this paper, we evaluate the second order radiative corrections in QED using the modified propagator. The three corrections are to the vacuum polarization, electron self-energy, and vertex function. In the conventional QED calculations the divergent terms are absorbed into the physical parameters like mass, charge and spin. The modified propagator here again acts as a regulator to the ultra-violet divergences in the theory. The modified propagator introduces two kinds of regulators in the radiative correction calculations, which are logarithmic and power law. In the three radiative corrections the leading order power law regulator is $`𝒪(L_P^2m^2)`$ and this is very small when compared to the other non-divergent terms in the conventional QED results. The three renormalization factors has the logarithmic corrections which is of the $`𝒪(\mathrm{ln}(L_Pm))`$.
The following point needs to be stressed regarding the actual values of the corrections. In pure QED, the perturbative expansion is in a series in $`\alpha `$ and corrections are of the order of $`\alpha 10^2`$, $`\alpha ^210^4`$ etc in successive orders. The lowest order quantum gravitational corrections are, by and large, of order $`L_P^2m^210^{45}`$. Since $`\alpha ^{22}(L_Pm)^2`$, the first 22 order corrections of QED will dominate over the quantum gravitational corrections computed here! Much before this, electroweak corrections will start modifying the results. Thus the finite corrections computed here are only of conceptual significance — providing the lowest order corrections from quantum gravity — rather than of any operational significance. The only exception to this general situation is when quantum gravitational effects break a symmetry originally present in the theory, which — as we shall see — does happen.
In sections II and III, we evaluate the second order radiative corrections in QED, and we discuss the results and present other applications in Sec. IV.
## 2 Vacuum Polarization and Electron Self Energy
### 2.1 Vacuum Polarization
The interaction of the photon field with electron field modifies the free photon propagator. The photon propagator of momentum $`q`$, with the one loop radiation correction included, is given by
$$iD_{\mu \nu }^F(q)=iD_{\mu \nu }^F+iD_{\mu \rho }^F\frac{i\mathrm{\Pi }^{\rho \sigma }(q)}{4\pi }iD_{\sigma \nu }^F,$$
where $`\mathrm{\Pi }^{\mu \nu }`$ is the vacuum polarization tensor, and $`D_{\mu \nu }^F`$ is the free photon propagator. The free photon propagator, $`D_{\mu \nu }^F`$, in the above relation is the conventional QED propagator. We concentrate here on the corrections to $`D_{\mu \nu }`$ arising from the modification of $`\mathrm{\Pi }_{\mu \nu }`$ rather than from direct modification of $`D_{\mu \nu }`$ due to our ansatz. Using the Feynman rules of QED in the momentum space, we can write the vacuum polarization tensor as
$`\mathrm{\Pi }_{\mu \nu }(q)`$ $`=`$ $`16\pi ie^2{\displaystyle \frac{d^4k}{(2\pi )^3}\left(k_\mu (kq)_\nu +(kq)_\mu k_\nu g_{\mu \nu }(k^2qkm^2)\right)}`$ (6)
$`\times `$ $`G_F^P(k)G_F^P(kq).`$
Substituting for the propagator from Eqn.(3), and following similar calculations as in conventional QED, the vacuum polarization tensor can be separated into gauge invariant and gauge non-invariant part <sup>?</sup>. The resultant gauge invariant part is given by
$$\mathrm{\Pi }_{\mu \nu }^1(q^2)=\frac{4e^2}{\pi }\left(q_\nu q_\mu g_{\mu \nu }q^2\right)_0^1𝑑zz(1z)K_0(\xi ),$$
(7)
where
$$\xi ^2=L_P^2\frac{m^2q^2z(1z)}{z(1z)}.$$
The series expansion of $`K_0(z)`$, about the origin, is given by
$$K_0(z)=\gamma \mathrm{ln}(z/2)\frac{z^2}{4}\left[1\gamma \mathrm{ln}(z/2)\right]+\mathrm{}.$$
(8)
Hence, Eqn.(7) takes the form
$`\mathrm{\Pi }_{\mu \nu }^1(q^2)`$ $`=`$ $`(q_\nu q_\mu g_{\mu \nu }q^2)[(Z_31){\displaystyle \frac{2e^2}{\pi }}A_1{\displaystyle \frac{e^2L_P^2}{\pi }}({\displaystyle \frac{q^2}{2}}A_1`$ (9)
$`+`$ $`{\displaystyle \frac{1}{2}}(1\gamma \mathrm{ln}(L_Pm/2))(m^2{\displaystyle \frac{q^2}{6}}))],`$
where
$$A_1=_0^1𝑑zz(1z)\mathrm{ln}\left(\frac{m^2}{m^2q^2z(1z)}\right)$$
(10)
to the lowest order of $`K_0(\xi )`$. The term $`A_1`$ is the familiar conventional QED non-divergent term and the remaining terms (of the order $`𝒪(L_P^2)`$) are the leading order quantum gravitational power law corrections/regulators to the conventional QED terms. The contribution of the quantum gravitational corrections to the vacuum polarization is extremely small as compared to the conventional non-divergent QED term, i.e. they are of the order $`10^{45}`$ which is much smaller compared to the next order radiative corrections of QED. But as we can notice the quantum gravitational corrections to the conventional QED non-divergent terms become important when $`L_Pq1`$. In this high momentum transfer limit, the propagation effects of the virtual photons will probe the small scale quantum gravitational effects and the corrections will be of the same order as the conventional QED non-divergent terms. The charge renormalization factor, $`Z_3`$, which is divergent in the usual QED calculations is now finite and is given by
$$Z_31=\frac{2e^2\mathrm{ln}(L_Pm/2)}{\pi }\left[\frac{1}{3}+\frac{6\gamma 5}{18\mathrm{ln}(L_Pm/2)}+𝒪(L_P^2m^2)\right].$$
(11)
The estimate of this factor comes out to be
$$Z_31=\frac{2e^2}{3\pi }\mathrm{ln}(L_Pm)0.1.$$
(12)
In the earlier work <sup>?</sup>, authors have calculated charge renormalization factor using Effective action approach, which is a non-perturbative technique. In this approach, they calculated the effect of the classical electro-magnetic background on the quantum charged scalar fields, propagating in the flat space-time, using the above modified propagator. The charge renormalization factor they obtain using the effective action approach is of the form $`Z^P=\frac{q^2}{6\pi }K_0(2L_Pm).`$ On expanding $`K_0(z)`$ using Eqn.(8), their estimate of the charge renormalization factor was also of the order of $`0.1`$.
In the standard QED calculations, there does arise a gauge non-invariant part of the vacuum polarization tensor which is divergent and is renormalized to zero(for instance, see Hatfield <sup>?</sup>). Even though, this is a standard result and can be found in textbooks we have given the relevant steps in the Appendix for the sake of completeness. By retracing the steps given in Appendix, we now obtain a finite quantity to the gauge non-invariant part of the vacuum polarization tensor using the modified propagator. The gauge non-invariant part comes out to be
$`\mathrm{\Pi }_{\mu \nu }^2(q)`$ $`=`$ $`{\displaystyle \frac{iL_P^2e_0^2}{(4\pi )^2}}g_{\mu \nu }{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dz}{z}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dz^{}}{z^{}}}{\displaystyle \frac{1}{(z+z^{})^2}}\mathrm{exp}\left({\displaystyle \frac{iL_P^2}{4}}(z^1+z^1)\right)`$ (13)
$`\times `$ $`\mathrm{exp}i\left({\displaystyle \frac{q^2zz^{}}{z+z^{}}}m_0(z+z^{})\right).`$
In Eqn.(13), transforming the variables $`z`$ and $`z^{}`$ to a new set of variables by the relation $`z=1/t`$ and $`z^{}=1/t^{}`$ leads to
$`\mathrm{\Pi }_{\mu \nu }^2(q)`$ $`=`$ $`{\displaystyle \frac{iL_P^2e_0^2}{(4\pi )^2}}g_{\mu \nu }{\displaystyle _0^{\mathrm{}}}t𝑑t{\displaystyle _0^{\mathrm{}}}t^{}𝑑t^{}{\displaystyle \frac{1}{(t+t^{})^2}}\mathrm{exp}\left({\displaystyle \frac{iL_P^2}{4}}(t+t^{})\right)`$ (14)
$`\times `$ $`\mathrm{exp}i\left({\displaystyle \frac{q^2}{t+t^{}}}m_0(t^1+t^1)\right).`$
Using the standard identity,
$$1=_0^{\mathrm{}}\frac{d\beta }{\beta }\delta \left(1\beta (z+z^{})\right)$$
and scaling $`t_it_i/\beta `$, as in the conventional QED calculations, we get
$$\mathrm{\Pi }_{\mu \nu }^2(q)=\frac{2ie_0^2}{(4\pi )^2}g_{\mu \nu }L_P^6_0^1t(1t)\frac{K_2(\xi )}{\xi ^2}𝑑t,$$
(15)
where
$$\xi ^2=L_P^2(\frac{m_0^2}{t(1t)}q^2)$$
The expansion of $`K_2(z)`$ near the origin is given by
$$K_2(z)=\frac{2}{z^2}\frac{1}{2}+\left(\frac{34\gamma }{32}+\frac{\mathrm{ln}(2)\mathrm{ln}(z)}{8}\right)z^2+\mathrm{}.$$
Substituting the series expansion of $`K_2(\xi )`$ in Eqn.(15), we obtain
$`\mathrm{\Pi }_{\mu \nu }^2(q)`$ $`=`$ $`g_{\mu \nu }{\displaystyle \frac{2ie_0^2}{(4\pi )^2}}[2L_P^2{\displaystyle _0^1}dt{\displaystyle \frac{(tt^2)^3}{(m_0^2q^2t(1t))^2}}`$ (16)
$``$ $`{\displaystyle \frac{L_P^4}{2}}{\displaystyle _0^1}dt{\displaystyle \frac{(tt^2)^2}{(m_0^2q^2t(1t))}}+\mathrm{}].`$
This clearly shows that the the gauge non-invariant part is nonzero for $`L_P0`$ and vanishes as $`L_P0`$. While the term which breaks the gauge invariance is small to be of operational significance, it does have certain conceptual importance. The following points need to be noted regarding this result: (i) Mathematically speaking, this result arises from the fact that our ansatz is equivalent to a momentum space regularization procedure in conventional QED. It is known that, momentum space regularization, in contrast to dimensional regularization, can lead to gauge breaking terms. Usually, this is considered as an argument in favor of dimensional regularization. In our approach, of course, we have no choice and the result arises automatically. In fact, it is very likely that any quantum gravitational cutoff will appear like a momentum space regulator and will break the gauge symmetry. (ii) Previously, Ohanian <sup>?</sup> had obtained a gauge breaking term using a gravitationally smeared propagator. There is however one vital difference between our result and the one obtained by him. Note that, in our approach, the propagator reduces to that of conventional field theory when $`L_P0`$. If the procedure is to be consistent, the gauge breaking term should vanish when the limit of $`L_P0`$ is taken. This is true as regards our result in Eqn.(16) showing that this is indeed a quantum gravitational effect. However, Ohanian <sup>?</sup> obtains a gauge breaking term which does not vanish in the corresponding limit. This suggests that, our approach does allow a consistent interpretation of the results. (iii) It is the extra term $`L_P^2(z+z^{})^2/(4zz^{})`$ which breaks the gauge invariance of the electro-magnetic field. This extra factor can be associated to the current in the charge conservation relation and hence, implying that the charge conservation is no more valid. There have been – more drastic! – attempts in the literature to break even the Lorentz invariance by introducing a coupling of the photon to charged scalar field, through gravitational couplings to the photon, etc <sup>?</sup>. The basic aim in the process of breaking the conformal symmetry of the electro-magnetic field is allowing for the possibility of generating large scale magnetic fields within inflationary scenarios. Most of these studies have used ad-hoc interaction potential to break the conformal invariance. The breaking of gauge invariance of the electro-magnetic field in our case is from a much more deeper “principle of path integral duality”. The connection to the cosmological seed magnetic field is still under investigation.
### 2.2 Electron Self Energy
The interaction of the electron field with photon field modifies the free electron propagator. The electron propagator of momentum $`p`$, with the one-loop correction included, is given by
$$iS_F^{}(p)=iS_F(p)+iS_F(p)(i\mathrm{\Sigma }(p))iS_F(p),$$
where $`\mathrm{\Sigma }(p)`$ (a 4-spinor) is the self-energy function, and $`S_F(p)`$ is the free electron propagator. Using the Feynman rules in the momentum space, we get
$$\mathrm{\Sigma }(p)=4\pi ie^2\frac{d^4k}{(2\pi )^4}\gamma ^\mu (\gamma ^\nu p_\nu \gamma ^\nu k_\nu +m)\gamma _\mu G_F^P(k)G_F^P(pk).$$
(17)
Substituting the propagator from the Eqn.(3), the above equation reduces to
$$\mathrm{\Sigma }(p)=\frac{e^2}{4\pi ^2}_0^1𝑑z\left(2m_0\gamma ^\mu p_\mu z\right)K_0(\xi ),$$
(18)
where
$$\xi ^2=\frac{L_p^2}{z(1z)}\left[m_0^2zp^2z(1z)\right].$$
On expanding $`K_0(\xi )`$ using Eqn.(8) we obtain, the lower order terms corresponding to the conventional QED results and the higher order terms as contributions of the quantum fluctuations of the space-time to the self energy of electron. The electron wave function renormalization, and the shift in the mass are obtained by recasting the interacting field propagator to look like the free field propagator i.e. we set
$$Z_2(\gamma ^\mu p_\mu m_0\mathrm{\Sigma }(p))=\gamma ^\mu p_\mu m+\mathrm{finite}\mathrm{terms}.$$
The electron wave function renormalization factor is obtained by equating terms proportional to $`p`$ in the above equation. This results in
$$Z_2^11=\frac{e^2}{8\pi ^2}\left[\gamma +\mathrm{ln}2\mathrm{ln}(L_pm_0)+𝒪(1)\right].$$
(19)
\[Note that we have neglected the higher order terms of $`K_0(\xi )`$ as these have the dependence as $`L_P^2m_0^2`$ whose contributions are negligible\]. The last term in the above expression is a finite quantity and is of the order one and it has the dependence of the electron momentum. The shift in the mass is given by
$$\frac{\delta m}{m_0}=\frac{e^2}{8\pi ^2}\left[3\mathrm{ln}(L_Pm_0/2)+3\gamma +𝒪(1)\right]$$
(20)
The estimate of the scale factor $`Z_2`$ comes out to be
$$Z_2^11=\frac{\alpha }{\pi }\mathrm{ln}m_0L_P0.1,$$
(21)
and the fractional shift in the mass is
$$\frac{\delta m}{m}\frac{\alpha }{\pi }\mathrm{ln}m_0L_P0.1.$$
(22)
The usual infrared divergences is ignored in the calculation of the renormalization factor $`Z_2`$. Here again, we see that both the mass renormalization factor and the mass shift are also of the same order as the charge renormalization factor i.e. $`0.1`$.
## 3 Vertex Correction and Anomalous magnetic moment
### 3.1 Vertex Correction
For a free Dirac field, the current density is defined as
$$J^\mu =\overline{\psi }\gamma ^0\gamma ^\mu \psi .$$
Thus, in the low energy QED processes the current transfer is related by $`\gamma ^\mu `$, which is the low energy vertex function. The radiative corrections will modify the vertex, to the one-loop correction, as
$$ie\mathrm{\Lambda }_\mu =ie\gamma _\mu ie\mathrm{\Gamma }_\mu ,$$
(23)
where, $`\mathrm{\Gamma }_\mu `$ is the vertex function. Using the Feynman rules, we obtain
$`\mathrm{\Gamma }_\mu (p^{},p)`$ $`=`$ $`(ie_0)^2{\displaystyle \frac{d^4k}{(2\pi )^4}\left[\gamma _\mu (\gamma _\rho p^\rho \gamma _\rho k^\rho +m_0)\gamma _\mu (\gamma _\sigma p^\sigma \gamma _\sigma k^\sigma +m_0)\right]}`$ (24)
$`\times `$ $`G_F^P(k)G_F^P(p^{}k)G_F^P(pk).`$
Substituting for the propagator from the Eqn. (3), the above equation reduces to
$`\mathrm{\Gamma }_\mu (p^{},p)`$ $`=`$ $`(Z_1^11)\gamma _\mu +{\displaystyle \frac{(ie_0)^2}{(4\pi )^2}}{\displaystyle _0^1}𝑑z{\displaystyle _0^{1z}}𝑑z^{}{\displaystyle \frac{\xi }{T_1}}`$
$`\times `$ $`\left[2\gamma _\mu (p^{}p)^2(1z^{})(1z)4imz(1zz^{})(p^{}p)^\mu \sigma _{\mu \nu }\right],`$
where
$$\xi ^2=L_P^2T_1T_2,T_1=(zp^{}+z^{}p)^2zp^2z^{}p^2+(z+z^{})m_0^2,$$
and
$$T_2=(1zz^{})^1+z^1+z^1.$$
The vertex renormalization factor $`Z_1`$ is given by
$`Z_1^11`$ $`=`$ $`i{\displaystyle \frac{(ie_0)^2}{(4\pi )^2}}{\displaystyle _0^1}dz{\displaystyle _0^{1z}}dz^{}[4K_0(\xi )+2m_0^2{\displaystyle \frac{\xi }{T_1}}K_1(\xi )`$ (26)
$`\times `$ $`(2+2(z+z^{})+(z+z^{})^2)].`$
Substituting for $`K_0(\xi )`$ using Eqn.(8), and using the series expansion of $`K_1(\xi )`$ as
$$K_1(z)=\frac{1}{z}+\frac{z}{2}\left(\mathrm{ln}(z/2)+\frac{2\gamma 1}{2}\right)+\frac{z^3}{16}\left(\mathrm{ln}(z/2)\frac{54\gamma }{4}\right),$$
(27)
in the Eqn. (26), the vertex normalization factor comes out to be
$`Z_1^11`$ $`=`$ $`i{\displaystyle \frac{(ie_0)^2}{(4\pi )^2}}{\displaystyle _0^1}dz{\displaystyle _0^{1z}}dz^{}[4\mathrm{ln}(\xi /2)4\gamma +{\displaystyle \frac{2m_0^2}{T_1}}`$ (28)
$`\times `$ $`(2+2(z+z^{})+(z+z^{})^2)]`$
$``$ $`{\displaystyle \frac{\alpha }{\pi }}\mathrm{ln}(m_0L_P),`$ (29)
which is roughly of the order of $`0.1`$. The estimated renormalization factors $`Z_1`$ and $`Z_3`$ using the modified propagator are not equal unlike in conventional QED.
### 3.2 Anomalous Magnetic Moment
The triumph of QED has been the precision test of electron anomalous magnetic moment. The experimental value of the anomalous magnetic moment of an electron is in excellent agreement with the predicted perturbative calculations up to the $`4^{th}`$ order to the $`15^{th}`$ decimal place. It is therefore of interest to compute the quantum gravitational corrections to the magnetic moment of an electron using the modified propagator.
The vertex correction contribution to scattering of an electron in an external field is given by $`\overline{u}(p^{})\mathrm{\Gamma }^\mu (p^{},p)u(p)A_\mu ^c(p^{}p)`$, where $`\overline{u}`$ and $`u`$ are the spinor wave-functions. Since, our interest is in calculating the radiative and the quantum gravitational corrections to the gyro-magnetic ratio of an electron, the term involving $`\sigma _{\mu \nu }`$ in the Eqn. (3.1) is of our concern. The corresponding $``$ matrix for this process is given by <sup>?</sup>,
$``$ $`=`$ $`\overline{u}(p^{})\gamma _\mu ^Mu(p)`$
$`=`$ $`\overline{u}(p^{})4m_0{\displaystyle \frac{(ie_0)^2}{(4\pi )^2}}{\displaystyle _0^1}𝑑z{\displaystyle _0^{1z}}𝑑z^{}z(1zz^{})(p^{}p)^\nu {\displaystyle \frac{\sigma _{\mu \nu }\xi }{T_1}}K_1(\xi )u(p).`$
In the limit of small momentum transfer, $`(p^{}p)^2m_0^2`$, one obtains by expanding $`K_1(\xi )`$ from the Eqn. (27), we get
$``$ $`=`$ $`\overline{u}(p^{}){\displaystyle \frac{\alpha }{m\pi }}(p^{}p)^\mu \sigma _{\mu \nu }{\displaystyle _0^1}𝑑z{\displaystyle _0^{1z}}𝑑z^{}{\displaystyle \frac{z(1zz^{})}{(z+z^{})^2}}u(p)`$ (31)
$``$ $`\overline{u}(p^{}){\displaystyle \frac{\alpha }{m\pi }}{\displaystyle \frac{L_P^2m_0^2}{24}}(p^{}p)^\mu \sigma _{\mu \nu }{\displaystyle _0^1}𝑑z{\displaystyle _0^{1z}}𝑑z^{}z(1zz^{})T_2`$
$`\times `$ $`\mathrm{ln}(L_Pm_0^2(z+z^{})^2T_2/4)u(p).`$
The first term in the above equation corresponds to the usual QED radiative vertex correction of the gyro-magnetic ratio (g) to the order $`e_0^2`$. The second term in the above equation is the quantum gravitational corrections to the gyro-magnetic ratio of an electron. The integral in the second term of the Eqn. (31) is convergent. The contribution of the quantum gravitational correction to the gyro-magnetic ratio to the first order in $`\alpha `$ is of the order $`L_P^2m^210^{45}`$. Obviously, this is not of practical significance.
## 4 Conclusions and Discussion
In this paper, we evaluated the quantum gravitational corrections(QGC) to three radiative corrections, in the first order of $`\alpha `$, in QED using the “principle of path integral duality”. The modified propagator is able to remove all the divergences which usually crop up in the conventional QED calculations. The main features of the modified propagator in QED are as follows: (a) The three renormalization factors($`Z_1`$, $`Z_2`$, and $`Z_3`$), and the mass shift are all of the same order, $`𝒪(\mathrm{ln}(mL_P))0.1`$. The charge renormalization factor $`Z_3`$ using the non-perturbative methods in scalar QED is also of order $`0.1`$ <sup>?</sup>. The renormalization factors $`Z_1`$ and $`Z_3`$ are different in our case as opposed to the conventional QED calculations. (b) The modified propagator makes the gauge non-invariant part of the vacuum polarization tensor to be non-zero for $`L_P0`$ vanishes and in the limit $`L_P0`$. This breaking of the gauge symmetry is also related to also the difference between the two renormalization factors. We have briefly indicated the possible effect of this breaking of gauge invariance, to the generation of large scale magnetic fields within inflationary scenarios (c) The contribution of the QGC to the gyro-magnetic ratio is very small i.e. of the order of $`L_P^2m^2`$ $`10^{45}`$. The contribution of the quantum gravitational correction to the vacuum polarization and the electron self energy is also of the same order.
Acknowledgments
S.S. is being supported by the Council of Scientific and Industrial Research, India.
Appendix A: Evaluation of Gauge Non-invariant part
For the sake of completeness, we outline the essential steps leading to the gauge non-invariant part of the vacuum polarization tensor in standard QED and how it is regularized to zero. The vacuum polarization tensor in the conventional QED calculations in $`2n`$ dimensions is given by
$`\mathrm{\Pi }_{reg}^{\mu \nu }(q)`$ $`=`$ $`(e\mu ^{2n})^2{\displaystyle \frac{d^{2n}k}{(2\pi )^{2n}}\frac{i}{k^2m^2+iϵ}\frac{i}{(kq)^2m^2+iϵ}}`$
$`\times `$ $`2^n\left(k^\mu (kq)^\nu +k^\nu (kq)^\mu ((k^2m^2)kq)g^{\mu \nu }\right),`$
where $`\mu `$ has the dimension of mass. Using the integral representation of the propagator, i.e.
$$\frac{i}{k^2m^2+iϵ}=_0^{\mathrm{}}𝑑z\mathrm{exp}\left(iz\left(k^2m^2+iϵ\right)\right)$$
(A.2)
in Eqn.(4) and completing the squares in the exponential, we obtain
$`\mathrm{\Pi }_{reg}^{\mu \nu }(q)`$ $`=`$ $`(e\mu ^{2n})^2{\displaystyle _0^{\mathrm{}}}𝑑z_1𝑑z_2{\displaystyle \frac{d^{2n}k}{(2\pi )^{2n}}}`$
$`\times `$ $`2^n\left(k^\mu (kq)^\nu +k^\nu (kq)^\mu ((k^2m^2)kq)g^{\mu \nu }\right)`$
$`\times `$ $`\mathrm{exp}\left(i(z_1+z_2)\left(k{\displaystyle \frac{z_2q}{z_1+z_2}}\right)^2+i{\displaystyle \frac{z_1z_2q^2}{z_1+z_2}}i(m^2iϵ)(z_1+z_2)\right).`$
Shifting the variable of integration and using the relations,
$$\frac{d^{2n}p}{(2\pi )^{2n}}p^2\mathrm{exp}(iap^2)=\frac{n}{(4\pi a)^n}\frac{1}{a}\mathrm{exp}(in\pi /2),$$
$$\frac{d^{2n}p}{(2\pi )^{2n}}p^\mu p^\nu \mathrm{exp}(iap^2)=\frac{g^{\mu \nu }}{(4\pi a)^n}\frac{1}{2a}\mathrm{exp}(in\pi /2),$$
we get
$$\frac{\mathrm{\Pi }_{reg}^{\mu \nu }(q)}{(e\mu ^{2n})^2}=\frac{\mathrm{exp}(in\pi /2)}{(4\pi )^n}\left[g_{\mu \nu }Q(q^2,m^2)+(q^\mu q^\nu g^{\mu \nu }q^2)P(q^2,m^2)\right],$$
(A.4)
where
$`Q(q^2,m^2)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑z_1{\displaystyle _0^{\mathrm{}}}𝑑z_2{\displaystyle \frac{2^n}{(z_1+z_2)^n}}\left({\displaystyle \frac{n1}{(z_1+z_2)}}+im^2iq^2{\displaystyle \frac{z_1z_2}{(z_1+z_2)^2}}\right)`$ (A.5)
$`\times `$ $`\mathrm{exp}\left(iq^2{\displaystyle \frac{z_1z_2}{z_1+z_2}}i(m^2iϵ)(z_1+z_2)\right),`$
and
$`P(q^2,m^2)`$ $`=`$ $`i2^{n+1}{\displaystyle _0^{\mathrm{}}}𝑑z_1{\displaystyle _0^{\mathrm{}}}𝑑z_2{\displaystyle \frac{z_1z_2}{(z_1+z_2)^{n+2}}}`$ (A.6)
$`\times `$ $`\mathrm{exp}\left(iq^2{\displaystyle \frac{z_1z_2}{z_1+z_2}}i(m^2iϵ)(z_1+z_2)\right).`$
Global gauge invariance of the action leads to current conservation, $`_\mu j^\mu =0`$. This implies $`q_\mu \mathrm{\Pi }^{\mu \nu }(q)=0`$. The factor proportional to $`g^{\mu \nu }`$ is the gauge non-invariant part of the vacuum polarization tensor and does not seem to satisfy the gauge invariance condition.
We summarize here the standard argument \[see Hatfield <sup>?</sup>\] to show that the first term, $`Q(q^2,m^2)`$, in Eqn.(A.4) can be shown to vanish and hence the regularization preserves the gauge symmetry. Consider only the first term of the gauge non-invariant part and define
$$I(q^2,m^2)_0^{\mathrm{}}𝑑z_1𝑑z_2\frac{1}{(z_1+z_2)^{n+1}}\mathrm{exp}\left(iq^2\frac{z_1z_2}{z_1+z_2}im^2(z_1+z_2)\right).$$
(A.7)
Rescaling the integration variables, $`z_1`$ and $`z_2`$, in the above equation by $`\beta `$, i.e. $`z_1\beta z_1`$ and $`z_2\beta z_2`$, we obtain
$`I(\beta ,q^2,m^2)`$ $`=`$ $`{\displaystyle \frac{1}{\beta ^{n1}}}{\displaystyle _0^{\mathrm{}}}𝑑z_1𝑑z_2{\displaystyle \frac{1}{(z_1+z_2)^{n+1}}}`$ (A.8)
$`\times `$ $`\mathrm{exp}\left(i\beta \left(q^2{\displaystyle \frac{z_1z_2}{z_1+z_2}}m^2(z_1+z_2)\right)\right).`$
The quantity $`I(\beta ,q^2,m^2)`$ is of course independent of the integration variables $`z_1`$ and $`z_2`$, and hence is also independent of the parameter($`\beta `$), i.e. $`I(\beta ,q^2,m^2)/\beta `$ $`=0`$. Differentiating the above expression w.r.t $`\beta `$ and regrouping terms, we get
$`0=\beta {\displaystyle \frac{I}{\beta }}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑z_1𝑑z_2{\displaystyle \frac{\beta ^{2n}}{(z_1+z_2)^n}}\left({\displaystyle \frac{n1}{(z_1+z_2)}}+im^2iq^2{\displaystyle \frac{z_1z_2}{(z_1+z_2)^2}}\right)`$ (A.9)
$`\times `$ $`\mathrm{exp}\left(i\beta \left(q^2{\displaystyle \frac{z_1z_2}{z_1+z_2}}m^2(z_1+z_2)\right)\right).`$
Rescaling the variables $`z_1`$ and $`z_2`$ by $`\beta ^1`$, i.e. $`z_iz_i/\beta `$, in the above expression, we get
$`0=\beta {\displaystyle \frac{I}{\beta }}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑z_1𝑑z_2{\displaystyle \frac{2^n}{(z_1+z_2)^n}}\left({\displaystyle \frac{n1}{(z_1+z_2)}}+im^2iq^2{\displaystyle \frac{z_1z_2}{(z_1+z_2)^2}}\right)`$ (A.10)
$`\times `$ $`\mathrm{exp}\left(i\left(q^2{\displaystyle \frac{z_1z_2}{z_1+z_2}}m^2(z_1+z_2)\right)\right)=Q(q^2,m^2).`$
Thus, the gauge non-invariant part of the vacuum polarization tensor in the standard QED vanishes. Hence, $`q_\mu \mathrm{\Pi }^{\mu \nu }(q)=0`$.
In the rest of this appendix, we retrace the above steps for the modified propagator. The gauge non-invariant part of the vacuum polarization tensor gets modified to the form,
$$\mathrm{\Pi }_{\mu \nu }^{(2)}(q^2,L_P)=(e\mu ^{2n})^2\frac{\mathrm{exp}(in\pi /2)}{(4\pi )^n}g_{\mu \nu }\mathrm{\Pi }_0^{(2)}(q^2,L_P),$$
(A.11)
where
$`\mathrm{\Pi }_0^{(2)}(q^2,L_P)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑z_1𝑑z_2{\displaystyle \frac{2^n}{(z_1+z_2)^n}}\left[{\displaystyle \frac{n1}{(z_1+z_2)}}+im^2iq^2{\displaystyle \frac{z_1z_2}{(z_1+z_2)^2}}\right]`$
$`\times `$ $`\mathrm{exp}\left[iq^2{\displaystyle \frac{z_1z_2}{z_1+z_2}}i(m^2iϵ)(z_1+z_2)+{\displaystyle \frac{iL_P^2}{4}}(z_1^1+z_2^1)\right],`$
when the modified propagator is used. We now show that the gauge non-invariant part is a finite quantity and is of the order $`L_P^2`$. Here again, we consider the first term of the gauge non-invariant part and define
$`I(q^2,m^2,L_P)`$ $``$ $`{\displaystyle _0^{\mathrm{}}}𝑑z_1𝑑z_2{\displaystyle \frac{1}{(z_1+z_2)^{n+1}}}\mathrm{exp}\left(iq^2{\displaystyle \frac{z_1z_2}{z_1+z_2}}im^2(z_1+z_2)\right)`$ (A.13)
$`\times `$ $`\mathrm{exp}\left({\displaystyle \frac{iL_P^2}{4}}(z_1^1+z_2^1)\right).`$
Rescaling the variables $`z_1`$ and $`z_2`$ by $`\beta `$, i.e $`z_i\beta z_i`$, and differentiating the resultant of the above expression w.r.t $`\beta `$ and regrouping terms, we obtain
$$\beta \frac{I(\beta ,q^2,m^2,L_P)}{\beta }Q(\beta ,q^2,m^2,L_P)R(\beta ,q^2,m^2,L_P),$$
(A.14)
where
$`Q(\beta ,q^2,m^2,L_P)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑z_1𝑑z_2{\displaystyle \frac{\beta ^{2n}}{(z_1+z_2)^n}}`$
$`\times `$ $`\left({\displaystyle \frac{n1}{(z_1+z_2)}}+im^2iq^2{\displaystyle \frac{z_1z_2}{(z_1+z_2)^2}}\right)`$
$`\times `$ $`\mathrm{exp}\left[i\beta \left(q^2{\displaystyle \frac{z_1z_2}{z_1+z_2}}m^2(z_1+z_2)\right)+{\displaystyle \frac{iL_P^2}{4\beta }}(z_1^1+z_2^1)\right],`$
and
$`R(\beta ,q^2,m^2,L_P)`$ $``$ $`{\displaystyle _0^{\mathrm{}}}𝑑z_1𝑑z_2{\displaystyle \frac{iL_P^2}{4z_1z_2}}{\displaystyle \frac{\beta ^n}{(z_1+z_2)^n}}`$
$`\times `$ $`\mathrm{exp}\left(i\beta \left(q^2{\displaystyle \frac{z_1z_2}{z_1+z_2}}m^2(z_1+z_2)\right)+{\displaystyle \frac{iL_P^2}{4\beta }}(z_1^1+z_2^1)\right).`$
The term $`I(\beta ,q^2,m^2,L_P)`$ is independent of $`\beta `$ (by the arguments stated earlier in this appendix) and hence, the partial differential $`I/\beta `$ vanishes. Hence,
$$Q(\beta ,q^2,m^2,L_P)=R(\beta ,q^2,m^2,L_P).$$
Rescaling the variables $`z_1`$ and $`z_2`$ by $`\beta ^1`$, i.e. $`z_iz_i/\beta `$, in Eqn.(A.14), we get
$$Q_{rescaled}(q^2,m^2,L_P)=R_{rescaled}(q^2,m^2,L_P)=\mathrm{\Pi }_0^{(2)}(q,L_P),$$
(A.17)
where
$`R_{rescaled}`$ $`=`$ $`{\displaystyle \frac{iL_P^2}{4}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dz_1}{z_1}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dz_2}{z_2}}{\displaystyle \frac{1}{(z_1+z_2)^n}}\mathrm{exp}\left({\displaystyle \frac{iL_P^2}{4}}(z_1^1+z_2^1)\right)`$ (A.18)
$`\mathrm{exp}i\left({\displaystyle \frac{q^2z_1z_2}{z_1+z_2}}m_0(z_1+z_2)\right),`$
and $`\mathrm{\Pi }_0^{(2)}(q,L_P)`$ is the quantity proportional to the gauge non-invariant part of the vacuum polarization tensor defined in Eqn.(A.11). Using our ansatz, we have shown that the gauge non-invariant part of the vacuum polarization tensor is a finite quantity and vanishes as $`L_P0`$.
References
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# Effect of boundary conditions on diffusion in two-dimensional granular gases
## I Introduction
The interactions among grains and between grains and the boundaries influence profoundly the macroscopic behavior of granular systems. To study such complex many body systems, numerical simulations are frequently used where one of the most important ingredients are the collision laws introduced to treat interactions . For dilute assemblies of grains one can use molecular dynamics algorithms where periodic boundary conditions are usually used. If the system is initially in a square box, a particle going out on the left re-enters the system on the right. We will show below that this kind of boundary condition modifies the general dynamics of the grains and introduces large correlations in time. This changes the diffusive behavior of the grains. In this paper, we propose an alternative approach to calculate numerically the coefficient of diffusion, accurately and with only very small finite size effects. To validate our methods, we compare our results for an elastic gas with the Langevin theory.
As granular gases are dissipative it is necessary to feed energy into the system to keep the particles agitated. To thermalize the system, we choose a random acceleration added to each grain at regular time step intervals $`dt`$. Our final goal in this paper is to study the dependence of the dynamic properties of the granular gas on the mode used to force the system. This work is a first step towards understanding the diffusion process in a binary system composed of two grain sizes. The system considered here is composed of one particle, s, of radius $`R_s`$ in a sea of particles of radius $`R_b`$. The particles are spheres constrained to move in a plane and which interact along their equators so that the system is two dimensional. The system considered here is dilute with a packing fraction of 30 %. The simulations are done with the molecular dynamics algorithms (time step driven and event driven ).
To characterize the diffusive behavior, we focus on the mean square displacement of the s particle. It is well known that for a 2D gas, the integral of the auto-correlation function does not converge . This means that the mean square displacement does not vary linearly with time. Therefore, strictly speaking, we cannot define a diffusion coefficient in 2D. However, we show that in a limited range of time, in the stationary state, the mean square displacement can be approximated by the linear function:
$$<(\stackrel{}{r}(t+t_0)\stackrel{}{r}(t_0))^2>4Dt$$
(1)
where $`D`$ can be interpreted as a diffusion coefficient. All quantities are expressed in arbitrary units.
## II Choice of boundary conditions
In this section we show that periodic boundary conditions introduce strong correlations and therefore alter the diffusion process.
### A Periodic boundary conditions
Consistent with common practice, we have used periodic boundary conditions to simulate a system of identical spheres $`R_s=R_b=0.5`$. Initially the particles are placed randomly in a square box of length $`L`$. The number of particles is calculated for each system depending on $`L`$, $`R_s`$ and the packing fraction. Periodic boundary conditions are applied in both directions. In this case, for elastic or forced gases (section IV), we have observed a strong dependence of $`D`$ (or of the mean square displacement) on the system size.
In Fig. 1, we have plotted the mean square displacement $`<(\stackrel{}{r}(t+t_0)\stackrel{}{r}(t_0))^2>`$ (Fig. 1a) and $`C(t)𝑑t`$ (Fig. 1b), both calculated in the stationary state, as function of $`t`$. $`C(t)`$ is the normalized autocorrelation function:
$$C(t)=\frac{<\stackrel{}{\mathrm{v}}(t_0+t)\stackrel{}{\mathrm{v}}(t_0)><\stackrel{}{\mathrm{v}}(t_0)>^2}{<\stackrel{}{\mathrm{v}}(t_0)^2><\stackrel{}{\mathrm{v}}(t_0)>^2}.$$
(2)
First, we note that the mean square displacement, at large time, varies linearly with time as expected but the slope of the curve, i.e. the diffusion coefficient, increases with system size. We show, in the inset to Fig. 1a, that this dependence on $`L`$ appears already at short time, when $`<(r^2(t+t_0)\stackrel{}{r}(t_0))^2>L^2`$. This feature can be also observed in $`\underset{0}{\overset{t\mathrm{}}{}}C(t)𝑑t`$, which is proportional to the diffusion coefficient. Similarly, we observe that the relaxation time $`\tau _r`$ (i.e. $`C(\tau _r)0`$) increases with size. In summary, the bigger the system is, the longer the characteristic time $`\tau _r`$ and the larger the diffusion coefficient $`D`$ are. We recall that such dependence has been observed by Alder et al . They proposed the following law for the dependence of $`D`$ on the number of particles, $`N`$,
$$D(N)=D(\mathrm{})(12/N).$$
(3)
However their numerical simulations do not support this conjecture since they fail to observe any saturation of $`D`$ for large systems. In addition, they found strong correlations in the velocity field characterized by the presence of vortex flow pattern at the microscopic scale. Our results confirm the lack of convergence for $`D`$ with system size. In addition, this variation of $`D`$ with $`L`$ is also observed in the case of inelastic collisions.
Another important remark is in order. If the system size is, for example, 60 (with 1400 particles of radius $`R=0.5`$), the characteristic time $`\tau _r`$ is found to be around $`20`$ which represents about 200 collisions for a particle. This means that a particle needs to undergo 200 collisions to lose completely the memory of its past. According to the Boltzmann theory this time should be limited to only a few collisions. Therefore we cannot accept this result as a valid macroscopic description of a gas. It is worth noting that the same results are found for both, the time step driven and the event driven algorithms.
We now discuss some points helpful for understanding the problem. Initially, each particle has a random velocity drawn from a Maxwellian distribution. We shift the linear and angular momenta so that the system has zero center of mass momentum and zero angular momentum relative to the center of mass. We find, however, that, although the system keeps its center of mass at rest throughout the simulation, the system is no longer isotropic, its moment of inertia becoming that of an ellipsoid. Let $`I(t)`$ be the inertia matrix of the system. Its two eigenvalues $`\lambda _n`$ and $`\lambda _p`$ are related via
$$\lambda _n+\lambda _p=m\underset{k=1}{\overset{N}{}}r_k^2(t),$$
(4)
where the sum is over all $`N`$ particles each of mass $`m`$. Following $`\lambda _n`$ and $`\lambda _p`$ in time shows that the system takes an ellipsoidal form ($`\lambda _n<\lambda _p`$). We have found, as well, an anisotropy in the diffusion tensor $`\widehat{D}(t)`$ defined from $`I(t)`$ as:
$$\widehat{D}(t)=\frac{1}{Nm}\frac{I(t+\delta t)I(t)}{\delta t}$$
(5)
As example we show, in Fig. 2 the two eigenvalues $`D_1`$ and $`D_2`$ of $`\widehat{D}`$ as functions of $`t`$ for a particular periodic system. Clearly, $`D_1`$ and $`D_2`$ are very different for all $`t`$. For all systems we studied, we found two different diffusion coefficients which depend strongly the system size. We were not able to find how these values scale with $`L`$.
In addition, we have found that, contrary to its initial condition, the system starts to rotate. This fact is put in evidence by calculating the two eigenvectors $`\stackrel{}{u}_n`$ and $`\stackrel{}{u}_p`$ of $`I(t)`$. These two (perpendicular) vectors rotate in space and, most importantly, they keep the same direction of rotation for a long time ($`\tau _r`$). We suspect that this rotation induces an anomalous temporal correlation of velocities. One should point out that this rotation phenomenon seems similar to that observed by Alder et al in their simulations with similar periodic boundary conditions.
We strongly believe that the use of periodic boundary conditions is responsible for this anomalous correlation. These boundary conditions present another inconvenience which is connected with the rotation of the system: The square geometry of the system does not permit the conservation of distances between two particles when the system is rotating. In Fig. 3 we show that after a rotation of $`\theta `$ the distance $`d_{ij}`$ between particles $`i`$ and $`j`$ can be drastically changed by the rotation if one of them goes through the boundary. Because the distance between particles is not conserved by rotation, the interaction potential used in the algorithm, which depends only on the relative positions of particles $`d_{ij}`$, is itself not invariant under rotation and so the angular momentum of the system is not conserved. Effectively the total angular momentum is fluctuating as one can see in Fig. 4. Every time a particle goes through the boundary, its angular momentum, $`l_z^i`$, changes sign. Consequently, the change in angular momentum is $`\mathrm{\Delta }L_z=2l_z^i`$. $`\mathrm{\Delta }L_z`$ is always proportional to $`L`$ (the system size) and the total number of particles, $`N`$, is proportional to $`L^2`$. However it appears that the fluctuations of $`L_z`$ get bigger with system size (see Fig. 4).
The use of periodic boundary conditions amounts to replicating the system on a square lattice. There are, therefore, several identical systems which interact through the boundaries. The rotation observed in our system is then extended to all these systems and can create some shear stress, due to frustration of rotation, between neighboring systems.
These boundary conditions can have other consequences on the dynamics of granular systems. For example, during the simulation of a cooling state the system evolves towards clusters whose orientation depends on the type of boundary conditions . In other similar simulations it was shown that there exist large spatial correlations between particles where the velocities stay correlated over a distance of about $`L/2`$.
### B Reflecting boundaries
In the case of reflecting boundaries the system is rotationally invariant leading to better behavior of the mean square displacement. However, in this case, the mean square displacement is limited at long time by the system size. To circumvent this problem, we proceed as follows. The test particle s is initially put at the center of system at $`t=0`$. The evolution of the position and velocity of this test particle are then followed until it reaches the boundary of the system in time $`t_w`$. We then repeat this many times collecting statistics for many test particles with different initial velocities.
The mean square displacement is calculated over 500 such trajectories and limited to time smaller than the smallest $`t_w`$. In this case, as one can see in Fig. 5, there is no dependence of the mean square displacement on the system size. Therefore, we can now trust the results of our numerical simulations.
We recall that the integral of the velocity correlation function does not converge in 2D. However, in a limited range of time (see Fig. 5), the quantity $`<(\stackrel{}{r}(t+t_0)\stackrel{}{r}(t_0))^2>`$ can be approximated by a straight line and $`D`$ calculated according to Eq. (1). Therefore, the estimate of $`D`$ with this method is an approximation.
## III Diffusion in an elastic gas
We first validate our algorithm using reflecting boundaries for an elastic gas (i.e., where the collision between particles are elastic). Then, we compare the numerical results with those given by the Langevin equation. Indeed, near equilibrium, the dynamics of s can be described approximately, by a Langevin equation:
$$\begin{array}{c}\frac{d\mathrm{v}_\mathrm{i}(t)}{dt}=\gamma \mathrm{v}_\mathrm{i}(t)+\mathrm{\Gamma }_i(t),\\ <\mathrm{\Gamma }_i(t)\mathrm{\Gamma }_j(t^{})>=q\delta _{i,j}\delta (tt^{}).\end{array}$$
(6)
where $`i`$ denotes the two direction $`x`$ and $`y`$. Integrating Eq. (6), the dependence of the mean square velocity on time is simply given by:
$$v^2(t)=v^2(0)e^{2\gamma t}+\frac{q}{\gamma }(1e^{2\gamma t}).$$
(7)
In this paper $`\mathrm{v}`$ denotes the instantaneous velocity of one particle and $`v^2`$ the mean square velocity averaged over the different s trajectories. We can rewrite Eq. (7) using the mean square velocity in the equilibrium state $`v^2(\mathrm{})=q/\gamma `$:
$$v^2(t)=v^2(\mathrm{})+(v^2(0)v^2(\mathrm{}))e^{2\gamma t}.$$
(8)
where $`1/\gamma `$ corresponds to the relaxation time.
For example, if $`R_s>>R_b`$ or, equivalently, $`m_s>>m_b`$ ($`m_{s,b}`$ is the mass of the particle of radius $`R_{s,b}`$), $`1/\gamma `$ is very large: the collision of s with a light particle $`b`$ will not affect strongly the velocity of s. So a great number of collisions is needed before s reaches its equilibrium state. Knowing the total kinetic energy of the system $`E_k^{tot}`$ (which is given by the initial velocity of each particle), we can easily calculate the square velocity in the equilibrium state $`v^2(\mathrm{})=q/\gamma `$ (using the Boltzmann distribution law for elastic gases). Using the simulations to calculate $`v^2(t)`$ for the s particles for initial conditions which are very different from the stationary state, i.e. $`v^2(0)v^2(\mathrm{})`$ and comparing with Eq. (8), we obtain the relaxation time for a big (heavy) particle (Fig. 7). Note that we are looking for agreement near the equilibrium state, where Eq. (6) is valid. Indeed the dissipation term $`\gamma `$ must depend on both velocities $`v_s^2`$ and $`v_b^2`$, as we will show.
Equation (6) also gives the mean square displacement as a function of time,
$$<(\stackrel{}{r}(t)\stackrel{}{r}(0))^2>=(v_0^2\frac{q}{\gamma })\frac{(1e^{\gamma t})^2}{\gamma ^2}+\frac{2q}{\gamma ^2}t\frac{2q}{\gamma ^3}(1e^{\gamma t}).$$
(9)
Comparing Eq. (9) and Eq. (1) at large time, the coefficient of diffusion is seen to be $`D=\frac{v^2(\mathrm{})}{2\gamma }`$. In Fig. 7, we compare the theoretical mean square displacement, Eq. (9), with the numerical one, obtained for the case $`R_s=3R_b`$. Note that in Eq. (9) all the parameters are known. Clearly, the agreement is very good. This confirms that even for a large test particle, the motion is well described by the simple Langevin equation. This observation, while reasonable, is not trivial since the limited size of our system and the radius of the particles are comparable to the mean free path.
We can now present a theoretical calculation of $`\gamma `$ which describes dissipation in the Langevin equation for all pairs ($`R_s`$,$`R_b`$). This will allow us to compare the theoretical values with the numerical ones as a function of $`R_s`$.
The value of $`\gamma `$ depends on both velocities, $`v_s`$ and $`v_b`$. To estimate theoretically the value of $`\gamma `$, we consider the deviation, due to a collision, of the particle s moving at $`v_s`$ in the $`x`$ direction. The dissipative term $`\gamma \stackrel{}{v_i}`$ appearing in Eq. (6), in the $`x`$ direction, can therefore be formally written as
$$\gamma v_s=<\frac{\stackrel{}{v^{}}_s.\stackrel{}{x}v_s}{v_s}>\omega _cv_s$$
(10)
where $`\stackrel{}{v^{}}_s`$ is the velocity after the collision and $`\omega _c`$ is the rate of collision. The symbol $`<>`$ in Eq. (10) corresponds to the average over all collisions between the s particle and the b ones. To calculate the different terms, we proceed as follows. We consider the collision of s with a particle b moving at a velocity $`\stackrel{}{v}_b`$. The collision is characterized by two angles: $`\theta `$, the angle between $`(\stackrel{}{r_s}\stackrel{}{r_b})`$ and the $`x`$ axis, and $`\phi `$ the angle between $`\stackrel{}{v_b}`$ and the $`x`$ axis. Then, for such a collision, illustrated in Fig. 6, we can calculate theoretically $`\stackrel{}{v_s^{}}(\theta ,\phi )`$, the final velocity of the s particle.
For elastic collisions, the projection of $`\stackrel{}{v_s^{}}(\theta ,\phi )`$ on the $`\stackrel{}{x}`$ direction is given by,
$$\stackrel{}{v_s^{}}(\theta ,\phi ).\stackrel{}{x}=\frac{m_sm_b}{m_s+m_b}v_scos^2(\theta )+\frac{2m_b}{m_s+m_b}v_bcos(\theta )cos(\theta \phi )+v_ssin^2(\theta ),$$
(11)
with the collision taking place only if
$$v_scos(\theta )v_bcos(\theta \phi )>0.$$
(12)
Integrating over $`\phi `$ and taking into account Eq. (12) we can write
$$<\stackrel{}{v_s^{}}(\theta ).\stackrel{}{x}>_\phi =\frac{\underset{0}{\overset{2\pi }{}}_v(v_scos(\theta )v_bcos(\theta \phi ))\stackrel{}{v_s^{}}(\theta ,\phi )\stackrel{}{x}𝑑\phi }{\underset{0}{\overset{2\pi }{}}_v(v_scos(\theta )v_bcos(\theta \phi ))},$$
(13)
where $`_v`$ is the Heavyside function. We found for Eq. (13) two solutions depending on the velocities. If $`v_s<v_b`$, we have
$$<\stackrel{}{v_s^{}}(\theta ).\stackrel{}{x}>_\phi =\frac{m_sm_b}{m_s+m_b}v_scos^2(\theta )+v_ssin^2(\theta )\frac{2m_bv_bcos(\theta )sin(\theta _p)}{(\pi \theta _p)(m_s+m_b)}$$
(14)
for all $`\theta `$ ($`0\theta \pi `$) and with $`\theta _p=\mathrm{cos}^1(v_s\mathrm{cos}(\theta )/v_b)`$. For the second case, $`v_s>v_b`$, there is a critical angle, $`\theta _c=\mathrm{cos}^1(v_b/v_s)`$, such that for $`\pi \theta _c<\theta \pi +\theta _c`$ the collision does not take place. In this case the solution of Eq. (13) is
$$\begin{array}{c}<\stackrel{}{v_s^{}}(\theta ).\stackrel{}{x}>_\phi =\frac{m_sm_b}{m_s+m_b}v_scos^2(\theta )+v_ssin^2(\theta )\text{ for}\mathrm{\hspace{0.25em}0}\theta <\theta _c,\hfill \\ <\stackrel{}{v_s^{}}(\theta ).\stackrel{}{x}>_\phi =\frac{m_sm_b}{m_s+m_b}v_scos^2(\theta )+v_ssin^2(\theta )\frac{2m_bv_bcos(\theta )sin(\theta _p)}{(\pi \theta _p)(m_s+m_b)}\text{ for}\theta _c\theta \pi \theta _c.\hfill \end{array}$$
(15)
We call $`\nu (\theta )`$ the mean relative loss of velocity, $`\nu (\theta )=\frac{<\stackrel{}{v^{}}_s(\theta ).\stackrel{}{x}>_\phi v_s}{v_s}`$, averaging only over the angle $`\phi `$. In Fig. 8 we show $`\nu (\theta )`$ for the particular case $`R_s=0.25`$ and $`R_b=0.5`$, which means that $`v_s>v_b`$. Note that in our calculation, the terms $`v_s`$ and $`v_b`$ correspond to the averaged values with respect to the appropriate Maxwellian distribution. To obtain the mean value, $`\stackrel{~}{\nu }`$, of $`\nu `$, we average by integrating numerically over $`\theta `$.
To conclude the calculation of the dissipative term, $`\gamma v_s`$, we have to estimate, using Eq. (10), the collision frequency which also depends on the velocities of the two particles. A similar calculation of $`\nu `$ can be done . In the stationary state where $`v_s^2`$ and $`v_b^2`$ are constant and the distributions of the velocities are Maxwellian one can use
$$\omega _c=\chi \sqrt{\pi }(R_s+R_b)d\sqrt{v_s^2+v_b^2},$$
(16)
where $`d`$ is the density of b particles and $`\chi `$ is a correction factor which corresponds to the local radial distribution around the s particle.
In Fig. 9, we compare, for different values of $`R_s`$, the diffusion coefficient found from the simulation with the theoretical value, $`v^2(\mathrm{})/2\gamma `$, predicted by the Langevin equation combined with our analytical calculation of $`\gamma `$. The theoretical calculation of $`\gamma `$ agrees very well with the simulation results. We recall that $`1/\gamma `$ corresponds to the characteristic time for the diffusive behaviour. It is important to notice that $`\gamma `$ can be approximated by $`\omega _c`$ only when $`R_s<<R_b`$. Effectively, the calculation for $`m_s0`$ gives $`\stackrel{~}{\nu }=1`$. Larger s particles need to suffer more than one collision to lose memory of their previous condition. For $`m_s\mathrm{}`$, $`\stackrel{~}{\nu }`$ is found equal to zero. Using these methods, we find for the elastic monodisperse case ($`R_s=R_b`$) that relaxation (decorrelation) takes place after about three collisions.
The agreement between numerical results and theoretical predictions allows us to confirm our numerical algorithm.
## IV Forced system
In a real granular system dissipation occurs through collisions, a fact that must be taken into account. Experimental mechanical properties of grains (restitution and friction coefficients) and collision laws are used in our simulations. The collisions between grains and the walls are treated with the same inelastic properties. Due to dissipation, we need to feed energy into the system to maintain the particles agitated. To accomplish this, we choose random heating : At every time step $`\delta t`$ we give a random acceleration, $`\eta _i(t)`$, in both spatial directions to each particle. The equation of motion can now be written formally as:
$$\begin{array}{c}m\frac{d\mathrm{v}_i}{dt}=F_i^c+F_i^t,\\ <F_i^t(t)F_j^t(t^{})>=m^2\delta _{i,j}\delta (tt^{})\eta _0^2.\end{array}$$
(17)
$`F_i^c`$ is the collision force acting on a particle of mass $`m`$. We chose the random acceleration, $`F_i^t/m`$, to be independent of the mass of the particle. It is given by a Gaussian noise of variance $`\eta _0^2`$.
At long time, the loss of energy due to collisions and the gain due to $`F^t`$ balance each other such that the system reaches a steady state out of equilibrium. It can be shown that the velocity distribution in this steady state is well described by a Maxwellian.
### A Stationary state
In the stationary state energy loss and gain balance exactly. The energy loss per unit time, $`\mathrm{\Gamma }`$, for the s particle, can be expressed as:
$$\mathrm{\Gamma }=P(m_s,m_b)\omega _cm_sv^2,$$
(18)
where $`P(m_s,m_b)`$ is the relative energy loss of particle s due to collisions. Clearly as for $`\stackrel{~}{\nu }`$, $`\mathrm{\Gamma }`$ must depend on the mass of the particle and on the two velocities $`v_s`$ and $`v_b`$. On the other hand, the gain of energy due the stochastic force is
$$\frac{1}{2}m_s[v^2(t+\delta t)v^2(t)]=m_s\eta _0^2\delta t.$$
(19)
In the steady state of the monodisperse system ($`R=R_s=R_b`$, and $`v^2(\mathrm{})=`$constant), we find, using Eq. (16), the following scaling for $`v^2(\mathrm{}):`$
$$\begin{array}{c}v^2(\mathrm{})(\eta _0^2)^{2/3}\\ v^2(\mathrm{})\tau _c\end{array}$$
(20)
We checked these two scaling laws numerically (see Fig. 10) and obtained the correct exponent $`2/3`$ for the various coefficients of restitution used in the contact laws. We have also verified the predicted dependence on $`\tau _c`$ for different values of $`R`$. The good agreement between theory and simulation indicates that we can describe the system by macroscopic continuous equations if $`\delta t<<\tau _c`$. As we explain elsewhere , the term $`P(m,m)`$ (in the monodisperse case) is independent of mass and velocity, because all particles are identical. This value of $`P(m,m)`$ was found equal approximatively to 0.145 for the mechanical properties corresponding to acetate spheres . We can then, in the case of a mono-disperse system, predict the dependence of the mean square velocity on the various parameters and, consequently, characterize the stationary state. For the bi-disperse case, the calculation is more complicated. Indeed the loss of energy depends on the two types of colliding particles and also on the different coefficients of restitution and friction introduced in the collision laws. As we show the dependence of $`P(m_s,m_b)`$ on $`v_s/v_b`$ is not trivial.
In this paper we limit ourselves to the effect of the thermalization mode (or random force) on the diffusion coefficient. To this end, we will compare in the following section the simulation results for $`D`$ with $`v^2(\mathrm{})/2\gamma `$ from the Langevin equation.
### B Diffusion of one particle
To estimate $`D`$, we use reflecting boundaries and the same method explained in section II B. We consider here the bi-disperse case (a single particle of radius $`R_s`$ in a sea of particles of radius $`R_b`$). As we have not yet found a theoretical expression for $`P(m_s,m_b)`$ for this case, we use for the mean square velocities the values obtained from the simulations which are shown in Fig. 11a. Note that $`\eta _0^2`$ has been chosen such that the value of $`v_b^2`$ is the same as in the previous section. We see that $`v^2(\mathrm{})`$ first decreases with $`R_s`$ for $`R_s<R_b`$ but then increases when $`R_s>R_b`$. Because of dissipation and the random acceleration, the repartition of the energy with the mass is no longer proportional to $`1/m_s`$. In all cases it is possible to calculate the mean collision frequency for s with Eq. (16) and the associated $`\gamma `$ value with Eq. (10). We can then calculate the relaxation time $`\tau _r`$ for all couples ($`R_s`$,$`R_b`$) used. In Fig. 11b we show the diffusion coefficient $`D`$, and the relaxation time $`\tau _r`$ -in the inset- as functions of $`R_s`$. The behavior of $`v^2(\mathrm{})`$ strongly modifies the curve of $`\tau _r`$ and $`D`$ versus $`R_s`$ compared to the elastic case. Note that the relaxation time represented in Fig. 11 clearly increases as $`R_s`$ increases.
We have seen in the elastic case that $`D=\frac{v^2}{2\gamma }`$. In Fig. 11b, we show the numerical results for $`D`$ as a function of $`R_s`$ and the corresponding values given by $`\frac{v^2}{2\gamma }`$. One can see clearly that the external noise modifies the dynamics of the granular gas and in particular the diffusion coefficient, $`D`$. The numerical value of $`D`$ is found to be larger than that obtained by the corresponding random walk. Indeed, at short time, due to the random force, $`v^2(t)`$ is not constant. Between two collisions $`v^2(t)`$ increases linearly with $`t`$. Starting with the equation of motion of particle s (between two collisions),
$$\frac{d\mathrm{v}_\mathrm{i}(t)}{dt}=\eta _i(t),$$
(21)
and with the initial conditions $`x_i(0)`$ and $`v_i(0)`$, we can calculate mean square displacement
$$<(x_i(t)x_i(0))^2>=\underset{t_1=0}{\overset{t}{}}𝑑t_1\left(\mathrm{v}_i(0)+\underset{0}{\overset{t_1}{}}\eta _i(t_1^{})𝑑t_1^{}\right)\underset{t_2=0}{\overset{t}{}}𝑑t_2\left(\mathrm{v}_i(0)+\underset{0}{\overset{t_2}{}}\eta _i(t_2^{})𝑑t_2^{}\right).$$
(22)
In two dimensions, this yields for the interval between two collisions
$$<(r(t)r(0))^2)>=v^2(0)t^2+\frac{2\eta _0^2}{3}t^3.$$
(23)
On the other hand, in the case of a random walk (or elastic collisions) the mean square displacement at short time scales as $`t^2`$. This difference explains the disagreement between $`D`$ and $`\frac{v^2(\mathrm{})}{2\gamma }`$. As the velocity changes between two collisions the probability of collision is increasing with time too. The calculation of the coefficient of diffusion is not easy in this case, due to the correlation between the velocity and the probability of collision (see Eq. 16). For very small particles, if one approaches relaxation by the time of a new collision, i.e. $`\stackrel{~}{\nu }1`$, this calculation should be possible. Indeed we can assume that the velocities before and after a collision are not correlated and have the same distribution (molecular chaos). We can then compute the mean square displacement, knowing the dependence of the collision probability on the velocity . With this assumption we improve the estimate of $`D`$ for the smallest $`R_s`$. But for bigger particles we have seen that the velocities stay correlated over many collisions and we can no longer use molecular chaos.
## V Conclusions
We have presented here some general results about the diffusion process in an agitated granular gas. We first showed that the boundary conditions used in the simulations are of crucial importance. Indeed, periodic boundary conditions introduce artificially strong temporal correlations which alter the macroscopic properties of the gas. If we ensure that no correlations are induced by the algorithm, for example by using reflecting boundaries, the numerical results obtained for an elastic gas can be described very well by a Langevin equation. We have presented a theoretical calculation of the relaxation time which allows us to predict the diffusion coefficient in all cases studied. This was not a priori intuitive since the radius of the particles is of the order of the mean free path. Finally we have analyzed the influence of uniform heating (a random acceleration) on dissipative gases. We have shown that heating influences the dynamics at short time. This is evident through the value of the diffusion coefficient which is different from that expected from the Langevin description. We are now applying with success these results to the diffusion process in a granular mixture consisting of two type of grains (differing by mass or size) in equal proportion.
ACKNOWLEDGMENTS
This work was partially funded by the CNRS Programme International de Cooperation Scientifique PICS $`\mathrm{\#}753`$ and the Norwegian research council, NFR.
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# The deformation quantization of certain super-Poisson brackets and BRST cohomology
## Introduction
In the usual programme of deformation quantization (cf. ) the quantum mechanical multiplication is considered as a formal associative deformation (a so-called star product) of the pointwise multiplication of the classical observables, viz. the algebra of smooth complex-valued functions on a given symplectic manifold. The deformation is such that to first order in the deformation parameter $`\lambda `$ (corresponding to $`\mathrm{}`$) the commutator of the deformed product is proportional to the Poisson bracket. The difficult question of existence of these star products for every symplectic manifold was settled independently by DeWilde and Lecomte and Fedosov , , and even for general Poisson manifolds by M. Kontsevitch, .
Four years ago, adequate formulations for star-products in the theory of supermanifolds, however, did not seem to have appeared in the literature although the geometric quantization scheme had found its suitable generalization to the super case (see e.g. and references therein). In order to elaborate our understanding of supermanifolds (at Freiburg) I proposed to the diploma student Ralf Eckel to give a formulation thereof in terms of associated bundles of certain jet group bundles which he did rather nicely in his diploma thesis in April 1996. He also provided a star-product formula for the case where the ‘fermionic directions’ formed a trivial vector bundle (see further down for a precise statement). When preparing my habilitation thesis in May 1996 I suddenly realized that a simple Fedosov procedure could be set up for general vector bundles: however, I did not know in advance a possible super-Poisson-bracket, so I first constructed the deformed algebra à la Fedosov and a posteriori computed the super-Poisson bracket as its first order supercommutator in , a result which I included in my habilitation thesis. A month later I was made aware by Amine El-Gradechi that the super-Poisson bracket I had computed out of this quantization exactly coincided with Rothstein’s super-Poisson bracket, see , found in 1991.
In this report I should like to give a detailed description of this Fedosov construction (thereby including an improved version of the preprint without some rather embarassing misprints). I shall also include sketches of two more recent constructions related to the above and done in collaboration with H.-C. Herbig and S. Waldmann (see and ), namely the quantum BRST cohomology for covariant star-products, and the classical BRST cohomology for arbitrary coisotropic constraint surfaces (the ‘reducible first-class case’ in the physics literature) where a so-called ‘ghosts-for-ghosts’-scheme is no longer necessary.
The supermanifolds I shall deal with in this report will only be ‘split’, more precisely, the set-up will be as follows: Let $`(M,\omega )`$ be a $`2m`$-dimensional symplectic manifold and $`E`$ be an arbitrary $`n`$-dimensional vector bundle over $`M`$. Then the algebra $`𝒞_0`$ of ‘classical superobservables’ can be considered as the space of all smooth sections of the complexified dual Grassmann algebra bundle of $`E`$ (see e. g. ), i.e.
$$𝒞_0:=\mathrm{\Gamma }(\mathrm{\Lambda }E^{}),$$
(1)
where the multiplication is the pointwise wedge product. Clearly, $`𝒞_0`$ is a $`_2`$-graded supercommutative algebra, i.e. $`\varphi \psi =(1)^{d_1d_2}\psi \varphi `$ for $`\varphi ,\psi \mathrm{\Gamma }(\mathrm{\Lambda }E^{})`$, $`\varphi `$ of degree $`d_1`$ and $`\psi `$ of degree $`d_2`$. A super-Poisson bracket for $`𝒞_0`$ is by definition a $`_2`$-graded bilinear map $`M_1:𝒞_0\times 𝒞_0𝒞_0`$ which is superanticommutative, i.e. $`M_1(\psi ,\varphi )=(1)^{d_1d_2}M_1(\varphi ,\psi )`$, satisfies the superderivation rule $`M_1(\varphi ,\psi \chi )=M_1(\varphi ,\psi )\chi +(1)^{d_1d_2}\psi M_1(\varphi ,\chi )`$, and the super Jacobi identity, i.e. $`(1)^{d_1d_3}M_1(M_1(\varphi ,\psi ),\chi )+cycl.=0`$ where $`\chi 𝒞_0`$ is of degree $`d_3`$.
It is general not difficult to find super-Poisson brackets of purely algebraic type, i.e. which vanish when one of their arguments is restricted to a smooth complex-valued function, by means of a fibre metric $`q`$ in $`E`$ (see e. g. , p. 123, eqn 5-1):
$$M_1^{}(\varphi ,\psi )=q^{AB}(j(e_A)\varphi )(i(e_B)\psi )$$
(2)
where $`q^{AB}`$ are the components of the induced fibre metric in the dual bundle $`E^{}`$ in the dual base to a local base $`(e_A)`$, $`1AdimE`$, of sections of $`E`$, and $`i(e_B)`$ and $`j(e_A)`$ denote the usual interior product left antiderivation and right antiderivation, respectively, which are also often denoted by $`\underset{A}{\overset{}{}}`$ and $`\underset{A}{\overset{}{}}`$ in the literature on supermanifolds. The above definition does not depend on the choice of that local base.
In case $`M`$ is $`^{2m}`$ with the standard Poisson bracket one can combine the standard bracket with the above super-Poisson bracket to get
$$M_1(\varphi ,\psi )=\frac{\varphi }{q^i}\frac{\psi }{p_i}\frac{\varphi }{p_i}\frac{\psi }{q^i}+q^{AB}(j(e_A)\varphi )(i(e_B)\psi ).$$
(3)
However, for nontrivial bundles it does not seem to be so obvious to generalize this bracket in the sense that it is equal to -at least in degree zero- the Poisson bracket of the base space $`M`$ when restricted to the sections of degree zero. Some time ago M.Rothstein has given a formula for this more general situation, :
$`\{\varphi ,\psi \}_R`$ $`=`$ $`\mathrm{\Lambda }^{ij}((12\widehat{R}^E)^1)_i^k__k^E\varphi __j^E\psi `$ (4)
$`+q^{AB}(j(e_A)(\varphi ))(i(e_B)(\psi ))`$
where $`^E`$ is a covariant derivative in the bundle $`E`$ preserving the fibre metric $`q`$ and $`\widehat{R}^E`$ is a suitable section in the bundle $`\mathrm{\Gamma }(Hom(TM,TM)\mathrm{\Lambda }^{\mathrm{even}}E^{})`$ constructed out of the curvature of $`^E`$ (see Section 1 for details).
The paper is organized as follows: In the first part I transfer Fedosov’s Weyl algebra bundle to the above situation by simply tensoring with the dual Grassmann bundle $`\mathrm{\Lambda }E^{}`$. The fibrewise multiplication has also a component in $`\mathrm{\Lambda }E^{}`$ built by means of the fibre metric in $`E`$. Then Fedosov’s procedure can completely be imitated without further difficulties: we show the existence of a Fedosov connection $`D`$ of square zero whose kernel in the space of antisymmetric degree zero is in linear $`11`$ correspondence to the space of formal power series in $`\lambda `$ with coefficients in $`𝒞_0`$
$$𝒞:=𝒞_0[[\lambda ]],$$
(5)
which immediately gives rise to the desired quantum deformation (Theorem 1.3). Then I explicitly compute the super-Poisson bracket $`M_1`$ as the term proportional to $`(𝐢\lambda )/2`$ by means of Fedosov’s recursion formulae (Theorem 1.4) and show that it is equal to the Rothstein superbracket. We evaluate the formulas a little bit further in part 2 in the case where the connection $`^E`$ is flat: using a local basis of covariantly constant sections the quantum multiplication is a sort of tensor product of a star-product on (the smooth complex-valued functions on) $`M`$ and a formal Clifford multiplication. Part 3 is concerned with a sketch of a quantized BRST formulation (see ). In Part 4 I shall sketch joint work with Hans-Christian Herbig where we have found a classical BRST complex for general coisotropic (reducible first class) constraint manifolds using the Rothstein superbracket, see also .
Notation: In all of this paper the Einstein sum convention is used that two equal indices are automatically summed up over their range unless stated otherwise. Moreover, we widely make use of Fedosov’s notation in with the following exceptions: we use the symbol $``$ to denote the covariant derivative and not Fedosov’s $``$ and describe the occurring symmetric tensor fields with $``$ products (see e.g. , p. 209-226) and use the symmetric substitution operator $`i_s`$ instead of Fedosov’s functions of $`y`$ and derivatives with respect to $`y`$.
## 1 A star-product for sections of Grassmann algebra bundles
This Section is -up to some corrected typos and additional remarks- identical to .
### 1.1 The Fedosov construction
Let $`(M,\omega )`$ be a $`2m`$-dimensional symplectic manifold and $`E`$ an $`n`$-dimensional real vector bundle over $`M`$ with a fixed nondegenerate fibre metric $`q`$. For the computations that will follow we shall use co-ordinates $`(x^1,\mathrm{},x^{2m})`$ in a chart $`U`$ of $`M`$. The base fields $`\frac{}{x^i}`$ will be denoted by $`_i`$ ($`1i2m)`$ for short. For computations in $`E`$ we shall use a local base $`(e_A)`$, $`(1An)`$ of sections of $`E`$ over $`U`$. Denote the dual base in the dual bundle $`E^{}`$ of $`E`$ by $`(e^A)`$, $`(1An)`$. Let $`\mathrm{\Lambda }\mathrm{\Gamma }(\mathrm{\Lambda }^2TM)`$ denote the Poisson structure of $`(M,\omega )`$, i.e. the Poisson bracket of two smooth real valued functions $`f,g`$ is given by $`\{f,g\}:=\mathrm{\Lambda }(df,dg)`$. Denoting the components of $`\omega `$ and $`\mathrm{\Lambda }`$ in that chart by $`\omega _{ij}:=\omega (_i,_j)`$ and $`\mathrm{\Lambda }^{ij}:=\mathrm{\Lambda }(dx^i,dx^j)`$ we use the sign conventions of where $`\mathrm{\Lambda }^{ik}\omega _{jk}=\delta _j^i`$. Fix a torsion-free symplectic connection $`^M`$ in the tangent bundle of $`M`$. This is well-known to always exist which can be seen by Heß’s formula $`\omega (_X^MY,Z):=\omega (\stackrel{~}{}_XY,Z)+\frac{1}{3}(\stackrel{~}{}_X\omega )(Y,Z)+\frac{1}{3}(\stackrel{~}{}_Y\omega )(X,Z)`$ where $`X,Y,Z`$ are arbitrary vector fields on $`M`$ and $`\stackrel{~}{}`$ is an arbitrary torsion-free connection on $`M`$ (see ). Fix a connection $`^E`$ in $`E`$ which is compatible with $`q`$, i.e. $`X(q(e_1,e_2))=q(_X^Ee_1,e_2)+q(e_1,_X^Ee_2)`$ for an arbitrary vector field $`X`$ on $`M`$ and sections $`e_1,e_2`$ of $`E`$. This is also well-known to always exist which can be seen by the formula $`q(_X^Ee_1,e_2):=q(\stackrel{~}{}_X^Ee_1,e_2)+\frac{1}{2}(\stackrel{~}{}_X^Eq)(e_1,e_2)`$ for an arbitrary connection $`\stackrel{~}{}^E`$ in $`E`$.
We are now forming the Fedosov algebra $`𝒲\mathrm{\Lambda }`$:
$$𝒲\mathrm{\Lambda }:=\left(\times _{s=0}^{\mathrm{}}\mathrm{\Gamma }((^sT^{}M\mathrm{\Lambda }E^{}\mathrm{\Lambda }T^{}M))\right)[[\lambda ]]$$
(6)
This is to say that the elements of $`𝒲\mathrm{\Lambda }`$ are formal sums $`_{s,t=0}^{\mathrm{}}w_{st}\lambda ^t`$ where the $`w_{st}`$ are smooth sections in the complexification of the vector bundle $`^sT^{}M\mathrm{\Lambda }E^{}\mathrm{\Lambda }T^{}M`$. In what follows we shall sometimes use the following factorized sections $`F:=f\varphi \alpha \lambda ^{t_1}`$ and $`G:=g\psi \beta \lambda ^{t_2}`$ where $`f\mathrm{\Gamma }(^{s_1}T^{}M)`$, $`g\mathrm{\Gamma }(^{s_2}T^{}M)`$, $`\varphi \mathrm{\Gamma }(\mathrm{\Lambda }^{d_1}E^{})`$, $`\psi \mathrm{\Gamma }(\mathrm{\Lambda }^{d_2}E^{})`$, $`\alpha \mathrm{\Gamma }(\mathrm{\Lambda }^{a_1}T^{}M)`$, and $`\beta \mathrm{\Gamma }(\mathrm{\Lambda }^{a_2}T^{}M)`$. Let $`deg_s,deg_E,deg_a,deg_\lambda `$ be the obvious degree maps from $`𝒲\mathrm{\Lambda }`$ to itself, i.e. those $``$-linear maps for which the above factorized sections $`f\varphi \alpha \lambda ^{t_1}`$ are eigenvectors to the eigenvalues $`s_1,d_1,a_1,t_1`$ respectively and which we refer to as the symmetric degree, the $`E`$-degree, the antisymmetric degree, and the $`\lambda `$-degree, respectively. Moreover, let $`P_E`$ and $`P_\lambda `$ be the corresponding maps $`(1)^{deg_E}`$ and $`(1)^{deg_\lambda }`$ which we refer to as the $`E`$-parity and the $`\lambda `$-parity, respectively. We say that a $``$-linear endomorphism $`\mathrm{\Phi }`$ of $`𝒲\mathrm{\Lambda }`$ is of $`\zeta `$-degree $`k`$ ($`\zeta =s,a,E,\lambda `$) iff $`[deg_\zeta ,\mathrm{\Phi }]=k\mathrm{\Phi }`$. Analogously, $`\mathrm{\Phi }`$ is said to be of $`\zeta `$-parity $`(1)^k`$ ($`\zeta =E,\lambda `$) iff $`P_\zeta \mathrm{\Phi }P_\zeta =(1)^k\mathrm{\Phi }`$. Let $`C`$ denote the complex conjugation of sections in $`𝒲\mathrm{\Lambda }`$.
We shall sometimes write $`𝒲`$ for the space of elements of $`𝒲\mathrm{\Lambda }`$ having zero antisymmetric degree and $`𝒲\mathrm{\Lambda }^a`$ for the space of those elements having antisymmetric degree $`a`$. The space $`𝒲\mathrm{\Lambda }`$ is an associative algebra with respect to the usual pointwise product where we do not use the graded tensor product of the two Grassmann algebras involved. More precisely, for the above factorized sections the pointwise or undeformed multiplication is simply given by
$$(f\varphi \alpha \lambda ^{t_1})(g\psi \beta \lambda ^{t_2}):=(fg)(\varphi \psi )(\alpha \beta )\lambda ^{t_1+t_2}.$$
(7)
Note that the above four degree maps are derivations and the above two parity maps are automorphisms of this multiplication. Moreover, $`𝒲\mathrm{\Lambda }`$ is supercommutative in the sense that
$$GF=(1)^{d_1d_2+a_1a_2}FG$$
(8)
A linear endomorphism $`\mathrm{\Phi }`$ of $`𝒲\mathrm{\Lambda }`$ of $`E`$-parity $`(1)^d^{}`$ and antisymmetric degree $`a^{}`$ is said to be a superderivation of type $`((1)^d^{},a^{})`$ of the undeformed algebra $`𝒲\mathrm{\Lambda }`$ iff $`\mathrm{\Phi }(FG)=(\mathrm{\Phi }F)G+(1)^{d^{}d_1+a^{}a_1}F(\mathrm{\Phi }G)`$. Let $`\sigma `$ denote the linear map
$$\sigma :𝒲\mathrm{\Lambda }\mathrm{\Gamma }(\mathrm{\Lambda }E^{}\mathrm{\Lambda }T^{}M)[[\lambda ]]$$
(9)
which projects onto the component of symmetric degree zero and clearly is a homomorphism for the undeformed multiplication.
We now combine the two covariant derivatives $`_X^M`$ in $`TM`$ and $`_X^E`$ in $`E`$ into a covariant derivative $`_X`$ in $`TME`$ in the usual fashion and extend it canonically to a connection $``$ in $`𝒲\mathrm{\Lambda }`$. On the above factorized sections we get in a chart:
$$(f\varphi \alpha )=\left((__i^Mf)\varphi +f(__i^E\varphi )\right)(dx^i\alpha )+f\varphi d\alpha .$$
(10)
In order to define a deformed fibrewise associative multiplication consider the following insertion maps for a vector field $`X`$ on $`M`$ and a section $`e`$ of $`E`$: let $`i_a(X)`$ and $`i(e)`$ denote the usual inner product antiderivations in $`\mathrm{\Gamma }(\mathrm{\Lambda }T^{}M)`$ and $`\mathrm{\Gamma }(\mathrm{\Lambda }E^{})`$, respectively, and extend them in a canonical manner to superderivations of type $`(1,1)`$ and $`(1,0)`$ of the undeformed algebra $`𝒲\mathrm{\Lambda }`$, respectively. Let $`j(e)`$ be defined by $`P_Ei(e)`$. Moreover, let $`i_s(X)`$ denote the corresponding inner product derivation (or symmetric substitution, , p.209-226) in $`\times _{s=0}^{\mathrm{}}\mathrm{\Gamma }(^sT^{}M)`$, again extended to a derivation of the undeformed algebra $`𝒲\mathrm{\Lambda }`$ in the canonical way. Let $`q^{AB}`$ denote the components of the induced fibre metric $`q^1`$ in $`E^{}`$, i.e. $`q^{AB}:=q^1(e^A,e^B)`$. Note that $`q^{AB}`$ is the inverse matrix to $`q(e_A,e_B)`$. Then for two elements $`F,G`$ of $`𝒲\mathrm{\Lambda }`$ we can now define the fibrewise deformed multiplication $``$:
$`FG`$ $`:=`$ $`{\displaystyle \underset{k,l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(i\lambda /2)^{k+l}}{k!l!}}\mathrm{\Lambda }^{i_1j_1}\mathrm{}\mathrm{\Lambda }^{i_kj_k}q^{A_1B_1}\mathrm{}q^{A_lB_l}`$ (11)
$`\left(i_s(_{i_1})\mathrm{}i_s(_{i_k})j(e_{A_1})\mathrm{}j(e_{A_l})F\right)`$
$`\left(i_s(_{j_1})\mathrm{}i_s(_{j_k})i(e_{B_1})\mathrm{}i(e_{B_l})G\right)`$
Moreover, let $`\delta `$ and $`\delta ^{}`$ be the canonical superderivations of the undeformed algebra $`𝒲\mathrm{\Lambda }`$ of type $`(1,1)`$ and $`(1,1)`$, respectively, which are induced by the identity map of $`T^{}M`$ to $`T^{}M`$ where in the case of $`\delta `$ the preimage of the identity is regarded as being part of $`T^{}M`$ and the image as being part of $`\mathrm{\Lambda }T^{}M`$, and vice versa in the case of $`\delta ^{}`$. On the above factorized sections these maps read in co-ordinates
$`\delta (f\varphi \alpha )`$ $`=`$ $`(i_s(_i)f)\varphi (dx^i\alpha )`$ (12)
$`\delta ^{}(f\varphi \alpha )`$ $`=`$ $`(dx^if)\varphi (i_a(_i)\alpha ).`$ (13)
Define the total degree derivation $`Deg`$:
$$Deg:=2deg_\lambda +deg_s+deg_E$$
(14)
A $``$-superderivation of type $`((1)^d^{},a^{})`$ is defined in an analogous manner as for the undeformed multiplication.
Note that $`(𝒲\mathrm{\Lambda },)`$ is not a graded associative algebra in the strict sense since it is equal to the cartesian product of the eigenspaces of $`Deg`$ and not to the direct sum of these eigenspaces. It is, however, filtered by those complex subspaces of $`𝒲\mathrm{\Lambda }`$ (indexed by a nonnegative integer $`r`$) which are given by the images of the maps $`Deg(Deg1)\mathrm{}(Deg(r1))`$.
We collect some properties of the above structures in the following
###### Proposition 1.1
With the above definitions and notations we have the following:
1. $`\delta ^2=0=(\delta ^{})^2`$ and $`\delta \delta ^{}+\delta ^{}\delta =deg_s+deg_a`$.
2. $`\delta +\delta =0`$.
3. $`Ker(\delta )Ker(deg_a)=𝒞`$.
4. $`P_E`$ is a $``$-automorphism and $`deg_a`$ is a $``$-derivation which equips the Fedosov algebra $`(𝒲\mathrm{\Lambda },)`$ with the structure of a $`_2\times `$-graded associative algebra.
5. $`\delta `$, $``$, and $`Deg`$ are $``$-superderivations of type $`(1,1)`$, $`(1,1)`$, and $`(1,0)`$, respectively.
6. The parity map $`P_\lambda `$ and the complex conjugation $`C`$ are graded $``$-antiautomorphisms, i.e. $`\mathrm{\Phi }(FG)=(1)^{d_1d_2+a_1a_2}\mathrm{\Phi }(G)\mathrm{\Phi }(F)`$ for $`\mathrm{\Phi }=P_\lambda ,C`$.
1. Straight forward.
2. This follows from the vanishing torsion of $`^M`$.
3. Without the factor $`\mathrm{\Lambda }E^{}`$ the kernel of $`\delta `$ in the space of antisymmetric degree zero consists of the constants, which proves this statement.
4. The associativity of $``$ is known, see e.g. , p. 123, eqn 5-2, and can be done by a long straight forward computation. We shall sketch a shorter proof: $``$ is defined on each fibre (for $`mM`$) $`𝒲_m:=(\times _{i=0}^{\mathrm{}}(^iT_mM^{}\mathrm{\Lambda }E_m^{}\mathrm{\Lambda }T_mM^{}))[[\lambda ]]`$ on which we can rewrite the multiplication in the more compact form ($`F,G𝒲_m`$)
$$FG=\mu \left(e^{\frac{𝐢\lambda }{2}\left(R+S\right)}\left(FG\right)\right)$$
(15)
where the tensor product is over $`[[\lambda ]]`$, $`\mu `$ denotes the undeformed fibrewise multiplication, and $`R:=\mathrm{\Lambda }_m^{ij}i_s(_i)i_s(_j)`$, $`S:=q_m^{AB}j(e_A)i(e_B)`$. Due to the derivation properties of $`i_s(_i)`$, $`i(e_A)`$, and $`j(e_B)`$ we get formulas like
$`R\mu 1`$ $`=`$ $`\mu 1\left(R_{13}+R_{23}\right)`$
$`R1\mu `$ $`=`$ $`1\mu \left(R_{12}+R_{13}\right)`$
$`S\mu 1`$ $`=`$ $`\mu 1\left(S_{13}\left(P_E\right)_2+S_{23}\right)`$
$`S1\mu `$ $`=`$ $`1\mu \left(S_{12}+S_{13}\left(P_E\right)_2\right)`$
where the index notation is borrowed from Hopf algebras and indicates on which of the three tensor factors of $`𝒲_m`$ the maps $`R`$, $`S`$, and $`P_E`$ should act, e.g. $`R_{23}:=1R`$, $`(P_E)_2:=1P_E1`$. These “pull through formulas” can be used to pull through the corresponding formal exponentials. Since all the maps $`i_s(_i)`$ commute with $`j(e_A)`$ and $`i(e_B)`$ and since the $`j(e_A)`$ commute with all $`i(e_B)`$ whereas $`j(e_A)`$ and $`j(e_B)`$ anticommute as well as $`i(e_A)`$ and $`i(e_B)`$ we can conclude that all the six maps $`R_{12}`$, $`R_{13}`$, $`R_{23}`$, $`S_{12}`$, $`S_{13}(P_E)_2`$, and $`S_{23}`$ pairwise commute. This is the essential step for associativity. The gradation properties are immediate.
5. The derivation properties of $`\delta `$ and $`Deg`$ are clear, for the corresponding statement for $``$ the fact that $`^M`$ preserves the Poisson structure $`\mathrm{\Lambda }`$ and that $`^E`$ preserves the dual fibre metric $`q^1`$ is crucial.
6. Straight forward. Q.E.D.
Due to the first part of this proposition we can construct a $`[[\lambda ]]`$-linear endomorphism $`\delta ^1`$ of the Fedosov algebra in the following way: on the above factorized sections $`F`$ we put
$$\delta ^1F:=\{\begin{array}{cc}\frac{1}{s_1+a_1}\delta ^{}F& \hfill \mathrm{if}s_1+a_11\\ 0& \hfill \mathrm{if}s_1+a_1=0\end{array}$$
(16)
Since $`𝒲\mathrm{\Lambda }`$ is an $`_2\times `$-graded associative algebra we can form the $`_2\times `$-graded super Lie bracket which reads on the above factorized sections:
$$[F,G]:=ad(F)G:=FG(1)^{d_1d_2+a_1a_2}GF$$
(17)
It follows from the associativity of $``$ that $`ad(F)`$ is $``$-superderivation of the Fedosov algebra $`(𝒲\mathrm{\Lambda },)`$ of type $`((1)^{d_1},a_1)`$. Note that the map $`\frac{𝐢}{\lambda }ad(F)`$ which we shall often use in what follows is always well-defined because of the supercommutativity of the undeformed multiplication (8).
Consider now the curvature tensors $`R^M`$ of $`^M`$ and $`R^E`$ of $`^E`$, i.e. for three vector fields $`X,Y,Z`$ on $`M`$ and a section $`e`$ of $`E`$ we have $`R^M(X,Y)Z=_X^M_Y^MZ_Y^M_X^MZ_{[X,Y]}^MZ`$ and $`R^E(X,Y)e=_X^E_Y^Ee_Y^E_X^Ee_{[X,Y]}^Ee`$. Define elements $`R^{(M)}`$ and $`R^{(E)}`$ of the Fedosov algebra which are contained in $`\mathrm{\Gamma }(^2T^{}M\mathrm{\Lambda }^2T^{}M)`$ and $`\mathrm{\Gamma }(\mathrm{\Lambda }^2E^{}\mathrm{\Lambda }^2T^{}M)`$, respectively, as follows where $`V,W`$ are vector fields on $`M`$ and $`e_1,e_2`$ are sections of $`E`$:
$`R^{(M)}(V,W,X,Y)`$ $`:=`$ $`\omega (V,R^M(X,Y)W)`$ (18)
$`R^{(E)}(e_1,e_2,X,Y)`$ $`:=`$ $`q(e_1,R^E(X,Y)e_2).`$ (19)
Note that this is well-defined: since $`^M`$ preserves $`\omega `$ and $`^E`$ preserves $`q`$ it follows that $`R^{(M)}`$ is symmetric in $`V,W`$ and $`R^{(E)}`$ is antisymmetric in $`e_1,e_2`$. In co-ordinates these two elements of the Fedosov algebra can be written in the form $`R^{(M)}=(1/4)R_{klij}^{(M)}dx^kdx^l1dx^idx^j`$ and $`R^{(E)}=(1/4)R_{ABij}^{(E)}1e^Ae^Bdx^idx^j`$. Set
$$R:=R^{(M)}+R^{(E)}.$$
(20)
Then the following Proposition is immediate:
###### Proposition 1.2
With the above definitions and notations we have:
1. $`^2=\frac{𝐢}{\lambda }ad(R)`$.
2. $`P_E(R)=R`$, $`P_\lambda (R)=R`$ and $`C(R)=R`$.
3. $`\delta R=0`$.
4. $`R=0`$.
1. Straightforward computation.
2. Obvious.
3. This is a consequence of the vanishing torsion of $`^M`$ (first Bianchi identity).
4. This is a reformulation of the second Bianchi identity for linear connections in arbitrary vector bundles. Q.E.D.
We shall now make the ansatz for a Fedosov connection $`D`$, i.e. we are looking for an element $`r𝒲\mathrm{\Lambda }^1`$ of even E-parity, i.e. $`P_E(r)=r`$, such that the map
$$D:=\delta ++\frac{𝐢}{\lambda }ad(r)$$
(21)
has square zero, i.e. $`D^2=0`$. The following properties of $`D`$ for any $`r`$ are crucial:
###### Lemma 1.1
Let $`r`$ be an arbitrary element of $`𝒲\mathrm{\Lambda }^1`$ of even $`E`$-parity. Then
1. $`D^2=\frac{𝐢}{\lambda }ad(\delta r+r+R+\frac{𝐢}{\lambda }rr)`$.
2. $`D(\delta r+r+R+\frac{𝐢}{\lambda }rr)=0`$.
This is straight forward using Proposition 1.2 and the fact that $`rr=\frac{1}{2}[r,r]`$ for the above elements $`r`$ of even $`E`$-parity and odd antisymmetric degree. Q.E.D.
For an arbitrary element $`w𝒲\mathrm{\Lambda }`$ we shall make the following decomposition according to the total degree $`Deg`$:
$$w=\underset{k=0}{\overset{\mathrm{}}{}}w^{(k)}\mathrm{where}Deg(w^{(k)})=kw^{(k)}$$
(22)
Note that each $`w^{(k)}`$ is always a finite sum of sections in some $`\mathrm{\Gamma }(^sT^{}M\mathrm{\Lambda }E^{}\mathrm{\Lambda }T^{}M)`$. The subspaces of all elements of $`𝒲\mathrm{\Lambda }`$, $`𝒲`$, $`𝒲\mathrm{\Lambda }^a`$, and $`𝒞`$ of total degree $`k`$ will be denoted by $`𝒲^{(k)}\mathrm{\Lambda }`$, $`𝒲^{(k)}`$, $`𝒲^{(k)}\mathrm{\Lambda }^a`$, and $`𝒞^{(k)}`$, respectively.
As in Fedosov’s paper there is the following
###### Theorem 1.1
With the above definitions and notations: Let $`r𝒲\mathrm{\Lambda }^1`$ be defined by the following recursion:
$`r^{(3)}`$ $`:=`$ $`\delta ^1R`$
$`r^{(k+3)}`$ $`:=`$ $`\delta ^1\left(r^{(k+2)}+{\displaystyle \frac{𝐢}{\lambda }}{\displaystyle \underset{l=1}{\overset{k1}{}}}r^{(l+2)}r^{(kl+2)}\right)`$
Then $`r`$ has the following properties: it is real ($`C(r)=r`$), depends only on $`\lambda ^2`$ ($`P_\lambda (r)=r)`$, has even $`E`$-parity, and is in the kernel of $`\delta ^1`$.
Moreover, the corresponding Fedosov derivation $`D=\delta ++(𝐢/\lambda )ad(r)`$ has square zero.
The behaviour of $`r`$ under the parity transformations and complex conjugation immediately follows from the fact that they commute with $`\delta ^1`$ and from their (anti)homomorphism properties (Prop.1.1, 3., 5.; Prop.1.2, 2.).
Let $`A:=\delta r+r+R+\frac{𝐢}{\lambda }rr=:\delta r+R+B`$. Recall the equation $`\delta \delta ^1+\delta ^1\delta =1`$ on the subspace of the Fedosov algebra where $`deg_s+deg_a`$ have nonzero eigenvalues. Clearly, $`A^{(2)}=\delta r^{(3)}+R=0`$ because $`\delta R=0`$ (Prop.1.2,3.) hence $`R=\delta \delta ^1R`$. Suppose $`A^{(l)}=0`$ for all $`2lk+1`$. By Lemma 1.1, 2. we have $`0=(DA)^{(k+1)}=\delta A^{(k+2)}=\delta B^{(k+2)}`$. Hence $`B^{(k+2)}=\delta \delta ^1B^{(k+2)}=\delta r^{(k+3)}`$ proving $`A^{k+2}=0`$ which inductively implies $`D^2=0`$ since we had already shown that $`r`$ is of even $`E`$-parity. Q.E.D.
We shall now compute the kernel of the Fedosov derivation. More precisely, define
$$𝒲_D:=Ker(D)Ker(deg_a).$$
(23)
As in Fedosov’s paper we have the important characterization:
###### Theorem 1.2
With the above definitions and notations: $`𝒲_D`$ is a subalgebra of the Fedosov algebra $`(𝒲\mathrm{\Lambda },)`$. Moreover, the map $`\sigma `$ (9) restricted to $`𝒲_D`$ is a $`[[\lambda ]]`$-linear bijection onto $`𝒞`$.
The kernel of a superderivation is always a subalgebra. Since $`D`$ and $`\sigma `$ are $`[[\lambda ]]`$-linear the subalgebra $`𝒲_D`$ is a $`[[\lambda ]]`$-submodule of $`𝒲`$.
Let $`w𝒲`$. Decompose $`w=w_0+w_+`$ where $`w_0:=\sigma (w)`$ and $`w_+:=(1\sigma )(w)`$. We shall prove by induction over the total degree $`k`$ that $`w𝒲`$ is in $`𝒲_D`$ iff for all nonnegative integers $`k`$ $`w_0^{(k)}`$ is arbitrary in $`𝒞^{(k)}`$ and $`w_+^{(k)}`$ is uniquely given by the equation
$$w_+^{\left(k\right)}=\delta ^1(w^{\left(k1\right)}+\frac{𝐢}{\lambda }\underset{l=1}{\overset{k2}{}}[r^{\left(l+2\right)},w^{\left(k1l\right)}])=:\delta ^1\left(\left(Aw\right)^{\left(k1\right)}\right)$$
(24)
where of course an empty sum is defined to be zero and $`w_+^{(0)}=0`$. Note that $`Dw=\delta w+Aw`$ and that the $``$-linear map $`A`$ does not lower the total degree of $`w`$.
Now the equation $`(Dw)^{(k)}=0`$ is equivalent to the inhomogeneous equation $`\delta (w^{(k+1)})=(Aw)^{(k)}`$. A necessary condition for this equation to be solvable for $`w^{(k+1)}`$ clearly is $`\delta ((Aw)^{(k)})=0`$. But this is also sufficient since then $`(Aw)^{(k)}=\delta \delta ^1(Aw)^{(k)}`$ and we have the particular solution $`w_+^{(k+1)}=\delta ^1(Aw)^{(k)}`$ (since $`\sigma \delta ^1=0`$) which satisfies (24). To this particular solution any solution to the homogeneous equation $`\delta (w^{(k+1)})=0`$ can be added which precisely is the space $`𝒞^{(k+1)}`$.
It remains to show that conversely every initial piece $`w^{}:=w_0^{(0)}+w_0^{(1)}+w_+^{(1)}+\mathrm{}+w_0^{(k)}+w_+^{(k)}`$ where $`w_0^{(l)}`$ was arbitrarily chosen in $`𝒞^{(l)}`$, $`w_+^{(l)}`$ is determined by (24) for all $`0lk`$, and $`(Dw^{})^{(l)}=0`$ for all $`1lk1`$ can be continued to $`w^{\prime \prime }:=w^{}+w_0^{(k+1)}+w_+^{(k+1)}`$ with $`w_0^{(k+1)}`$ arbitrary in $`𝒞^{(k+1)}`$, $`w_+^{(k+1)}`$ determined by (24), and $`(Dw^{\prime \prime })^{(k)}=0`$. By induction, this will eventually lead to $`w𝒲_D`$ characterized by the above properties. Indeed, since $`D^2=0`$ we have $`0=(D^2w^{})^{(k1)}=\delta ((Dw^{})^{(k)})=\delta ((Aw^{})^{(k)})`$. Define $`w_+^{(k+1)}`$ by $`\delta ^1((Aw^{})^{(k)})`$ and choose any $`w_0^{(k+1)}𝒞^{(k+1)}`$. It follows at once that $`w_+^{(k+1)}`$ satisfies (24) and that we get $`(Dw^{\prime \prime })^{(k)}=0`$ which proves the induction and the Theorem. Q.E.D.
Let
$$\tau :𝒞𝒲_D𝒲$$
(25)
be the inverse of the restriction of $`\sigma `$ to $`𝒲_D`$. For $`\varphi \mathrm{\Gamma }(\mathrm{\Lambda }E^{})`$ we shall speak of $`\tau (\varphi )`$ as the Fedosov-Taylor series of $`\varphi `$ and refer to the components $`\tau (\varphi )^{(k)}`$ as the Fedosov-Taylor coefficients. We collect some of the properties of $`\tau `$ in the following
###### Proposition 1.3
With the above definitions and notations:
1. $`\tau `$ commutes with $`P_E`$, $`P_\lambda `$, and $`C`$.
2. Let $`\varphi =_{d=0}^n\varphi ^{(d)}\mathrm{\Gamma }(\mathrm{\Lambda }E^{})`$ where $`n:=dimE`$. Then $`Deg(\varphi ^{(d)})=d\varphi ^{(d)}=deg_E(\varphi ^{(d)})`$.
Moreover
$`\tau (\varphi )^{(0)}`$ $`=`$ $`\varphi ^{(0)}`$ (26)
$`\tau (\varphi )^{(1)}`$ $`=`$ $`\delta ^1(\varphi ^{(0)})+\varphi ^{(1)}`$ (27)
$`\mathrm{}`$ $`\mathrm{}`$
$`\tau (\varphi )^{(n)}`$ $`=`$ $`\delta ^1\left((\tau (\varphi )^{(n1)})+{\displaystyle \frac{𝐢}{\lambda }}{\displaystyle \underset{l=1}{\overset{n2}{}}}[r^{(l+2)},\tau (\varphi )^{(n1l)}]\right)+\varphi ^{(n)}`$
$`\tau (\varphi )^{(n+1)}`$ $`=`$ $`\delta ^1\left((\tau (\varphi )^{(n)})+{\displaystyle \frac{𝐢}{\lambda }}{\displaystyle \underset{l=1}{\overset{n1}{}}}[r^{(l+2)},\tau (\varphi )^{(nl)}]\right)`$ (29)
$`\mathrm{}`$ $`\mathrm{}`$
$`\tau (\varphi )^{(k+1)}`$ $`=`$ $`\delta ^1\left((\tau (\varphi )^{(k)})+{\displaystyle \frac{𝐢}{\lambda }}{\displaystyle \underset{l=1}{\overset{k1}{}}}[r^{(l+2)},\tau (\varphi )^{(kl)}]\right)`$ (30)
where $`kn`$. The Fedosov-Taylor series $`\tau (\varphi )`$ depends only on $`\lambda ^2`$.
3. For any nonnegative integer $`k`$ the map $`\varphi \tau (\varphi )^{(k)}`$ is a polynomial in $`\lambda `$ whose coefficients are differential operators from $`\mathrm{\Gamma }(\mathrm{\Lambda }E^{})`$ into some $`\mathrm{\Gamma }(^sT^{}M\mathrm{\Lambda }E^{})`$ of order $`k`$.
Since $`r`$ is invariant under the parity maps and complex conjugation, it follows that $`D`$ commutes with these three maps, hence $`𝒲_D`$ is stable under these maps. Since $`\sigma `$ obviously commute with them, so does the inverse of its restriction to $`𝒲_D`$, $`\tau `$. The rest is a consequence of the preceding Theorem and a straight forward induction. Q.E.D.
Define the following $`[[\lambda ]]`$-bilinear multiplication on $`𝒞`$: for $`\varphi ,\psi 𝒞`$
$$\varphi \psi :=\sigma (\tau (\varphi )\tau (\psi )).$$
(31)
We shall call $``$ the Fedosov star product associated to $`(M,\omega ,^M,E,q,^E)`$. For $`\varphi ,\psi \mathrm{\Gamma }(\mathrm{\Lambda }E^{})`$ the star product $`\varphi \psi `$ will be a formal power series in $`\lambda `$ which we shall write in the following form:
$$\varphi \psi =:\underset{t=0}{\overset{\mathrm{}}{}}\left(\frac{𝐢\lambda }{2}\right)^tM_t(\varphi ,\psi ).$$
(32)
We list some important properties of the Fedosov star product in the following
###### Theorem 1.3
With the above definitions and notations:
1. The Fedosov star product is associative and $`_2`$-graded, i.e. $`P_E`$ is an automorphism of $`(𝒞,)`$. The map $`P_\lambda `$ and the complex conjugation $`C`$ are graded antiautomorphisms of $`(𝒞,)`$.
2. The $``$-bilinear maps $`M_t`$ are all bidifferential, real, vanish on the constant functions in each argument for $`t1`$, and have the following symmetry property:
$$M_t(\psi ,\varphi )=(1)^t(1)^{d_1d_2}M_t(\varphi ,\psi ).$$
(33)
3. The term of order $`0`$ is equal to the pointwise Grassmann multiplication. Hence $`(𝒞,)`$ is a formal associative deformation of the supercommutative algebra $`(𝒞_0,)`$.
Basically, every stated property is easily derived from the definitions (31) and (32) and the corresponding behaviour of the fibrewise multiplication $``$ under $`P_E`$, $`P_\lambda `$, and $`C`$. The reality of the $`M_t`$ follows easily from the graded antihomomorphism property of $`C`$ once eqn (33) is proved by means of the graded antihomomorphism property of the $`\lambda `$-parity. Since $`\tau (1)`$ is easily seen to be equal to $`1`$ we have $`1\psi =\psi =\psi 1`$, and the $`M_t`$ must vanish on $`1`$ for $`t1`$. Finally, each $`M_t`$ obviously depends on only a finite number of Fedosov-Taylor coefficients whence it must be bidifferential. Q.E.D.
### 1.2 Computation of the super-Poisson bracket
In this section we are going to compute an explicit expression for the term $`M_1`$ of the Fedosov star product defined in the last section (compare (32) and Theorem 1.3). Only by means of the graded associativity of the deformed algebra $`(𝒞,)`$ we can derive the following
###### Lemma 1.2
Let $`\varphi ,\psi ,\chi `$ be sections in $`𝒞_0`$ of $`E`$-degree $`d_1,d_2,d_3`$, respectively. Then
$`M_1(\psi ,\varphi )`$ $`=`$ $`(1)^{d_1d_2}M_1(\varphi ,\psi )`$ (34)
$`M_1(\varphi ,\psi \chi )`$ $`=`$ $`M_1(\varphi ,\psi )\chi +(1)^{d_1d_2}\psi M_1(\varphi ,\chi )`$ (35)
$`0`$ $`=`$ $`(1)^{d_1d_3}M_1(M_1(\varphi ,\psi ),\chi )+\mathrm{cycl}.`$ (36)
Hence $`M_1`$ is a super-Poisson bracket on $`𝒞_0`$.
The first property is a particular case of (33). Consider now the graded commutator $`[\varphi ,\psi ]:=\varphi \psi (1)^{d_1d_2}\psi \varphi `$ on $`𝒞`$. Because of the graded associativity of $``$ we have the superderivation property $`[\varphi ,\psi \chi ]=[\varphi ,\psi ]\chi +(1)^{d_1d_2}\psi [\varphi ,\chi ]`$. Writing this out with the $`M_t`$ and taking the term of order $`\lambda `$ we get the second property. For the third, take the term of order $`\lambda ^2`$ in the super Jacobi identity for the graded commutator. Q.E.D.
Before we are going to compute $`M_1`$ directly it is useful to introduce the following notions:
For $`\varphi `$ in $`𝒞_0`$ let $`\varphi _1`$ and $`\rho `$ denote the component of symmetric degree one and $`\lambda `$-degree zero of the Fedosov-Taylor coefficient $`\tau (\varphi )`$ and the section $`r`$ (Theorem 1.1), respectively. Note that $`\varphi _1`$ is a smooth section in the bundle $`T^{}M\mathrm{\Lambda }E^{}`$. Denote by $`\mathrm{\Lambda }_0E^{}`$ the subbundle of the dual Grassmann bundle consisting of elements of even degree. Then $`\rho `$ is a smooth section in $`T^{}M\mathrm{\Lambda }_0E^{}T^{}M`$. Consider now the bundle $`TM\mathrm{\Lambda }_0E^{}T^{}M`$. There is an obvious fibrewise associative multiplication $``$ in that bundle which comes from the identification of $`TMT^{}M`$ with the bundle of linear endomorphism of $`TM`$: let $`X,Y`$ be vector fields on $`M`$, $`\varphi ,\psi \mathrm{\Lambda }_0E^{}`$, and $`\alpha ,\beta `$ one-forms on $`M`$. Then
$$(X\varphi \alpha )(Y\psi \beta ):=(\alpha (Y))X(\varphi \psi )\beta .$$
(37)
Let $`\widehat{R}^E`$ be the section in $`\mathrm{\Gamma }(TM\mathrm{\Lambda }^2E^{}T^{}M)`$ whose components in a bundle chart read
$$\widehat{R}^E:=\frac{1}{4}\mathrm{\Lambda }^{ik}R_{ABkj}^{(E)}_ie^Ae^Bdx^j=:_i(\widehat{R}^E)_j^idx^j,$$
(38)
and let $`\widehat{\rho }\mathrm{\Gamma }(TM\mathrm{\Lambda }_0E^{}T^{}M)`$ be defined by
$$\widehat{\rho }:=_i\mathrm{\Lambda }^{ik}i_s(_k)\rho =:_i\widehat{\rho }_j^idx^j.$$
(39)
Note that we can form arbitrary power series in $`\widehat{R}^E`$ by using the multiplication $``$ since $`\widehat{R}^E`$ is nilpotent.
We have the following
###### Lemma 1.3
With the above notations and definitions:
$`M_1(\varphi ,\psi )`$ $`=`$ $`\mathrm{\Lambda }^{ij}(i_s(_i)\varphi _1)(i_s(_j)\psi _1)+q^{AB}(j(e_A)(\varphi ))(i(e_B)(\psi ))`$
$`\varphi _1`$ $`=`$ $`dx^j((1\widehat{\rho })^1)_j^i__i^E\varphi `$ (41)
$`\widehat{\rho }`$ $`=`$ $`1(12\widehat{R}^E)^{1/2}.`$ (42)
where $`(1\widehat{\rho })^1`$ and $`(12\widehat{R}^E)^{1/2}`$ denote the corresponding power series with respect to the $``$ multiplication.
The first equation is a straight forward computation.
For the second, use the Fedosov recursion for $`\tau (\varphi )`$, (Proposition 1.3), note that $`\varphi _1^{(k)}`$ is zero for $`kn+2`$ and that only the component $`\rho `$ of $`r`$ matters since both $`\tau (\varphi )`$ and $`r`$ depend only on $`\lambda ^2`$, sum over the total degree which yields the equation
$$\varphi _1=\delta ^1^E\varphi +dx^j\left(\widehat{\rho }\right)_j^i\left(i_s\left(_i\right)\varphi _1\right)$$
which proves the second equation.
For the third, use the Fedosov recursion for $`r`$, (Theorem 1.1), take the component of symmetric degree 1 and $`\lambda `$-degree zero, sum over the total degree, and arrive at the quadratic equation
$$\widehat{\rho }\widehat{R}^E=\frac{1}{2}\widehat{\rho }\widehat{\rho }.$$
Since $`r`$ and hence $`\rho `$ does not contain components of symmetric degree zero, there is only one solution to this equation, namely the above third equation. Q.E.D.
This Lemma immediately implies the desired formula for the super-Poisson bracket:
###### Theorem 1.4
The super-Poisson bracket $`M_1`$ obtained by the Fedosov star product takes the following form:
$`M_1(\varphi ,\psi )`$ $`=`$ $`\mathrm{\Lambda }^{ij}((12\widehat{R}^E)^{1/2})_i^k((12\widehat{R}^E)^{1/2})_j^l__k^E\varphi __l^E\psi `$
$`+q^{AB}(j(e_A)(\varphi ))(i(e_B)(\psi ))`$
Clear from the Lemma ! Q.E.D.
###### Corollary 1.1
The above super Poisson bracket coincides with the Rothstein super Poisson bracket $`\{,\}_R`$, see (4) and .
Since by definition $`\mathrm{\Lambda }^{ij}(\widehat{R}^E)_i^k=\mathrm{\Lambda }^{ki}(\widehat{R}^E)_i^j`$ the same relation holds for any power series (with respect to $``$) $`(f(\widehat{R}^E))_i^k)`$ whence the result. Q.E.D.
Remarks:
1. In case $`(M,\omega )`$ is Kähler there exist star products of Wick type on $`M`$ (see , ): they are characterized by the property that for any two complex-valued smooth functions $`f,g`$ on $`M`$ the star product $`f^{}g`$ is made out of bidifferential operators which differentiate $`f`$ in holomorphic directions only and $`g`$ in antiholomorphic directions only. It seems to me very likely that super analogues of these star products can readily be formulated for any complex holomorphic hermitean vector bundle over $`M`$ as it has been done in geometric quantization, see .
2. If the dual Grassmann bundle $`\mathrm{\Lambda }E^{}`$ is replaced by the symmetric power $`E^{}`$ and the fibre metric $`q`$ by some antisymmetric bilinear form on the fibres covariantly constant by some connection in $`E`$ the whole construction can presumably carried through as well (see also Neumaier’s related construction for differential operators in , Section 3). As we shall explain further down this can be interpreted as a particular case of a symplectic fibration.
3. It may also be interesting to compute this construction in the particular case where $`M`$ is the cotangent bundle of an arbitrary semi-Riemannian manifold $`Q`$ and $`E`$ is the tangent bundle of $`Q`$ pulled back to $`T^{}Q`$ by the bundle projection. Star-products on $`T^{}Q`$ are strongly related to (pseudo) differential operator calculus on $`Q`$, see , , , and . In that situation one could study asymptotic representation theory incorporating Dirac operators. T. Voronov has studied the algebra $`𝒞`$ using symbol calculus and its representations on the space of differential forms on $`Q`$ (which is an intermediate step towards spinors), see .
Note added: The above Fedosov construction is not the full Fedosov construction one would expect in supermanifold theory as I have been made aware by the referee: there the super-Fedosov algebra should rather consist of a sort of completed tensor product of supersymmetric tensor fields (generalizing $`\mathrm{\Gamma }(T^{}M)`$) and superdifferential forms (generalizing $`\mathrm{\Gamma }(\mathrm{\Lambda }T^{}M)`$) which would include our $`𝒲\mathrm{\Lambda }`$, but also -roughly speaking- additional symmetric tensors and differential forms ‘in the purely fermionic directions’. Moreover the fibrewise multiplication would involve the full Rothstein superbracket. It is very probable that such a super Fedosov construction will go through without any big conceptual problem and, since to my best knowledge this has not yet been done in the literature, will be an interesting problem to attack.
I believe that the rôle of the above Fedosov construction can perhaps best be compared with the constructions which have been done in the meantime by B.Fedosov and O.Kravchenko for ordinary (i.e. non super) symplectic fibrations (see , ): they are using an intermediate Fedosov construction which starts with a ‘purely vertical’ star-product on the symplectic fibres satisfying some compatibility conditions which is supposed to already exist; in a second step the Fedosov construction proper is then only done for the base, but ‘tensored’ with the ‘vertical algebras’: the result is a star-product on the total space. The curvature of the fibre bundle underlying the symplectic fibration enters in the symplectic form of the total space when it is expressed in terms of the symplectic form on the base and on the fibres.
It seems to me that an even symplectic split supermanifold can be regarded as a ‘supersymplectic fibration’ with symplectic base and ‘purely fermionic’ fibres, and the simple nature of the Rothstein super symplectic form exactly corresponds to that picture. Moreover, in the Fedosov construction presented in this contribution the ‘fermionic vertical direction’, viz: the algebra $`\mathrm{\Gamma }(\mathrm{\Lambda }E^{})`$ already carries a simple explicit vertical star-product, namely a sort of formal Clifford multiplication (see also the next Section), and the construction is intermediate insofar that symmetric and antisymmetric tensor fields only come from the base. It is an interesting question under which circumstances the ‘full’ Fedosov construction for even symplectic supermanifolds (which will no doubt be much more complicated) reduces to the above ‘intermediate construction’.
## 2 Flat vector bundles
An important particular case is given by a vector bundle $`E`$ with fibre metric $`q`$ on which there exists a flat covariant derivative $`^E`$, for instance in the case of a trivial bundle $`M\times ^n`$ with $`q`$ being a nondegenerate bilinear form on $`^n`$ not depending on $`M`$.
Note first the standard fact for flat vector bundles that there is an open cover $`(U_\alpha )_{\alpha I}`$ of $`M`$ together with a basis of local sections $`e_A`$, $`1An`$ defined on each $`U_\alpha `$ which are covariantly constant and which are related by constant transition matrices on the overlaps of any two of the $`U_\alpha `$.
We have the following
###### Lemma 2.1
With the above additional assumptions the following holds:
1. The map $`r`$ as defined in Theorem 1.1 does not depend on $`\mathrm{\Lambda }E^{}`$, i.e. is contained in $`\times _{s=0}^{\mathrm{}}\mathrm{\Gamma }((^sT^{}M\mathrm{\Lambda }^1T^{}M))[[\lambda ]]`$.
2. The Fedosov-Taylor series of a function $`fC^{\mathrm{}}(M)`$ does not depend on $`\mathrm{\Lambda }E^{}`$, i.e. is contained in $`\times _{s=0}^{\mathrm{}}\mathrm{\Gamma }((^sT^{}M))[[\lambda ]]`$.
3. The Fedosov-Taylor series of a local covariantly constant section $`\varphi `$ in $`\mathrm{\Gamma }(\mathrm{\Lambda }E^{})`$ is equal to $`\varphi `$.
1. Since $`r^{(3)}=R^{(M)}`$, and since $`\delta ^1`$ and $``$ preserve $`\times _{s=0}^{\mathrm{}}\mathrm{\Gamma }((^sT^{}M\mathrm{\Lambda }T^{}M))[[\lambda ]]`$ which is a fibrewise subalgebra of $`𝒲\mathrm{\Lambda }`$ the statement follows by induction using Theorem 1.1.
2. The proof is completely analogous to part 1. upon using the formulas in Prop. 1.3.
3. Again by induction using Prop. 1.3 where the fact is used that $`r`$ supercommutes with $`\mathrm{\Gamma }(\mathrm{\Lambda }E^{})`$ according to 1. Q.E.D.
This immediately implies the following formula for the star-product:
###### Theorem 2.1
We make the above assumptions. Let $`\varphi ,\psi `$ two sections in $`𝒞`$ and express them locally as $`\varphi =_{d=0}^n\frac{1}{d!}\varphi _{A_1\mathrm{}A_d}e^{A_1}\mathrm{}e^{A_d}`$ and likewise for $`\psi `$ where $`e^A`$, $`1An`$ is a local base of covariantly constant sections of $`E^{}`$ and $`\varphi _{A_1\mathrm{}A_d}`$ are local $`C^{\mathrm{}}`$-functions. Then
$$\varphi \psi =\underset{d,d^{}=0}{\overset{n}{}}\frac{1}{d!}\frac{1}{d^{}!}(\varphi _{A_1\mathrm{}A_d}_\mathrm{F}\psi _{B_1\mathrm{}B_d^{}})(e^{A_1}\mathrm{}e^{A_d}_{\mathrm{Cl}}e^{B_1}\mathrm{}e^{B_d^{}})$$
(43)
where $`_\mathrm{F}`$ denotes the usual Fedosov star-product on $`M`$ defined by the map $`r`$ (restricted to $`\times _{s=0}^{\mathrm{}}\mathrm{\Gamma }((^sT^{}M\mathrm{\Lambda }^1T^{}M))[[\lambda ]]`$) and $`_{\mathrm{Cl}}`$ denotes the formal tensorial Clifford multiplication in $`𝒞`$ defined by
$`\varphi _{\mathrm{Cl}}\psi =`$ $`:=`$ $`{\displaystyle \underset{l=0}{\overset{n}{}}}{\displaystyle \frac{(i\lambda /2)^l}{l!}}q^{A_1B_1}\mathrm{}q^{A_lB_l}`$
$`\left(j(e_{A_1})\mathrm{}j(e_{A_l})\varphi \right)\left(i(e_{B_1})\mathrm{}i(e_{B_l})\psi \right).`$
The above formula (43) does not depend on the chosen covariantly constant local trivializaton.
The subalgebra $`\times _{s=0}^{\mathrm{}}\mathrm{\Gamma }((^sT^{}M\mathrm{\Lambda }T^{}M))[[\lambda ]]`$ of $`𝒲`$ is clearly preserved by the Fedosov derivative $`D`$ whence it follows at once that $`fg=f_\mathrm{F}g`$ for all $`f,gC^{\mathrm{}}(M)[[\lambda ]]`$. Moreover, for a covariantly constant section $`\chi `$ of $`𝒞`$ we clearly have $`f\chi =\sigma (\tau (f)\tau (\chi ))=\sigma (\tau (f)\tau (\chi ))`$ using the above Lemma and the properties of $``$ whence $`f\chi =f\chi =\chi f`$. Finally, note that $`\chi \chi ^{}=\chi \chi ^{}=\chi _{\mathrm{Cl}}\chi ^{}`$ for two covariantly constant sections, where the result is again covariantly constant, and therefore
$$\left(f\chi \right)\left(g\chi ^{}\right)=f\chi g\chi ^{}=fg\chi \chi ^{}=\left(f_\mathrm{F}g\right)\left(\chi _{\mathrm{Cl}}\chi ^{}\right)$$
which proves the above formula. Since the transition functions are constant it follows that (43) does not depend on the chosen local basis of covariantly constant sections. Q.E.D.
Conversely, it is easy to see that the above formula (43) always defines an associative $`_2`$-graded deformation of $`𝒞_0`$ where $`_\mathrm{F}`$ can be replaced by any given star-product on $`M`$: It is locally given by the tensor product over $`[[\lambda ]]`$ of the associative algebra $`(C^{\mathrm{}}(M)[[\lambda ]]`$ with the formal Clifford algebra $`(\mathrm{\Lambda }^n[[\lambda ]],_{\mathrm{Cl}})`$.
For a trivial flat bundle without holonomy the above formula had been given by R. Eckel in his thesis , p. 66.
## 3 A quantum BRST complex for quantum covariant star-products
The results of this Section have been obtained in collaboration with Hans-Christian Herbig and Stefan Waldmann in .
Let $`(M,\omega )`$ a symplectic manifold. Suppose that a Lie group $`G`$ (with Lie algebra $`𝔤`$) symplectically and properly acts on $`M`$ (e.g. when $`G`$ is compact) allowing for a classical momentum map $`J:M𝔤^{}`$: for each $`\xi 𝔤`$ let $`\xi _M`$ be the fundamental field $`md/dt(exp(t\xi )m)|_{t=0}`$, then $`\omega ^{\mathrm{}}(\xi _M)=dJ,\xi `$ and $`J(gm)=\mathrm{Ad}^{}(g)J(m)`$ for all $`gG,mM`$. This implies the Lie homomorphism property
$$\{J,\xi ,J,\eta \}=J,[\xi ,\eta ]$$
(44)
for all $`\xi ,\eta 𝔤`$. Recall the Marsden-Weinstein phase space reduction scheme, : suppose for the rest of this Section that $`0`$ is a regular value of $`J`$ and that the constraint surface $`C:=J^1(0)`$ is nonempty. Then $`G`$ acts locally freely on $`C`$, and supposing that $`G`$ acts freely and properly on $`C`$ the quotient manifold $`M_{\mathrm{red}}:=C/G`$ becomes a symplectic manifold, its symplectic form $`\omega _{\mathrm{red}}`$ being determined by the condition that its pull-back to $`C`$ by the canonical projection equals the restriction of $`\omega `$ to $`TC`$. Note that each $`G`$-invariant smooth function on $`C`$ naturally projects to $`M_{\mathrm{red}}`$.
Now let $``$ be a star-product on $`M`$. According to and a formal power series $`𝑱=_{r=0}^{\mathrm{}}\lambda ^r𝑱_rC^{\mathrm{}}(M,𝔤^{})[[\lambda ]]`$ will be called a quantum momentum map and the star-product $``$ (quantum) covariant iff $`𝑱_0=J`$ and analogously to (44):
$$𝑱,\xi 𝑱,\eta 𝑱,\eta 𝑱,\xi =i\lambda 𝑱,[\xi ,\eta ]$$
(45)
for all $`\xi ,\eta 𝔤`$. We call $`(M,,G,𝑱,C)`$ satisfying the previous conditions a Hamiltonian quantum $`G`$-space with regular constraint surface. According to a Theorem by Fedosov \[19, Sect. 5.8\] an even stronger condition can be achieved for all such group actions preserving a connection, e.g. proper actions (since they always preserve a Riemannian metric), namely strong invariance:
$$J,\xi ffJ,\xi =i\lambda \{J,\xi ,f\},$$
(46)
which obviously implies (45) setting $`𝑱=J`$. A simple example is provided by the standard Moyal-Weyl-star-product on $`^{2n}`$ together with the Lie algebra of all infinitesimal linear symplectic transformations represented by the space of all quadratic homogeneous polynomials. The problem whether a general classical momentum map can be deformed into a quantum momentum map for a suitable star-product is still an open problem as far as I know.
We are now constructing a BRST complex related to that problem (see for a general introduction the book and our article for more references): consider the trivial bundle $`E:=(𝔤𝔤^{})\times M`$ together with the fibre metric $`q`$ defined by the natural pairing between $`𝔤`$ and $`𝔤^{}`$. Then the superobservable algebra $`𝒞`$ (called $`𝒜`$) in ) of the first Section equals
$$𝒞=\mathrm{\Lambda }𝔤^{}\mathrm{\Lambda }𝔤C^{\mathrm{}}(M)[[\lambda ]].$$
(47)
As a $`[[\lambda ]]`$-module this space carries natural $``$-gradings, namely the ghost degree (form degree in $`\mathrm{\Lambda }𝔤^{}`$), the antighost degree (form degree in $`\mathrm{\Lambda }𝔤`$), and the ghost number $`\mathrm{𝖦𝗁}`$ which is defined as the difference of the ghost degree and the antighost degree and which we shall consider as a $`[[\lambda ]]`$-linear map $`𝒞𝒞`$ with the ghost number integers as eigenvalues. We shall write $`𝒞^{i,j}`$ for the submodule of all those elements having ghost degree $`i`$ and antighost degree $`j`$ and $`𝒞^{(i)}`$ for the submodule of all those elements having ghost number $`i`$. We equip $`𝒞`$ with a star-product as in Section 2, (43) where the initial star-product on $`M`$ does not have to be of Fedosov type. Consider now the following three elements of $`𝒞`$: $`𝑱𝒞^{1,0}`$, $`\mathrm{\Omega }:=1/2[,]𝒞^{2,1}`$, and $`\gamma :=`$ one half of the identity homomorphism of $`𝔤`$, contained in $`𝒞^{1,1}`$. Let $`\mathrm{\Theta }:=𝑱+\mathrm{\Omega }`$, the so-called BRST-charge which is contained in $`𝒞^{(1)}`$. Define the BRST operator $`Q`$ by
$$Q(\varphi ):=\frac{1}{i\lambda }\left(\mathrm{\Theta }\varphi (1)^{a+b}\varphi \mathrm{\Theta }\right)a,b\varphi 𝒞^{a,b}$$
(48)
Then we have the following
###### Theorem 3.1
Let $`(M,,G,𝐉,C)`$ be a Hamiltonian quantum $`G`$-space with regular constraint surface. Then
1. The Ghost number operator $`\mathrm{𝖦𝗁}`$ is equal to $`\varphi \frac{1}{i\lambda }(\gamma \varphi \varphi \gamma )`$ and therefore is a derivation of $`(𝒞,)`$ which thus becomes a $``$-graded associative algebra.
2. $`\mathrm{\Theta }\mathrm{\Theta }=0`$.
3. The BRST operator has square zero, $`Q^2=0`$, and is a superderivation of ghost number one of $`(𝒞,)`$.
The proof of this statement is a rather straight-forward consequence of equation (45). For more details see .
Define the quantum BRST cohomology by $`\mathrm{Ker}Q/\mathrm{Im}Q=:𝑯_{\text{BRST}}(𝒞[[\lambda ]])`$. Then we have the following
###### Theorem 3.2
Let $`(M,,G,𝐉,C)`$ be a Hamiltonian quantum $`G`$-space with regular constraint surface. Then
1. $`𝑯_{\text{BRST}}(𝒞[[\lambda ]])`$ becomes a $``$-graded associative algebra in a canonical way.
2. There is a representation $`\mathit{\varrho }_C`$ of the Lie algebra $`𝔤`$ on the $`[[\lambda ]]`$-module $`C^{\mathrm{}}(C)[[\lambda ]]`$ deforming the representation $`\varrho _C`$ induced by the restriction of the fundamental fields to $`C`$ such that the quantum BRST cohomology is isomorphic to the Chevalley-Eilenberg cohomology of $`𝔤`$ with values in $`C^{\mathrm{}}(C)[[\lambda ]]`$ with respect to $`\mathit{\varrho }_C`$.
3. In particular, the component of ghost number zero of the quantum BRST cohomology is isomorphic to the submodule of all those elements in $`C^{\mathrm{}}(C)[[\lambda ]]`$ which are invariant under $`\mathit{\varrho }_C`$.
See again for a detailed proof.
In case the Hamiltonian action of the connected Lie group $`G`$ on $`M`$ is proper and the reduced phase space exists we can choose a strongly invariant star-product on $`M`$ (see (46)). Under these circumstance we have the stronger
###### Theorem 3.3
With the assumption of the previous Theorem and the above additional assumptions we have:
1. The quantum BRST-cohomology is isomorphic to the Chevalley-Eilenberg cohomology of $`𝔤`$ with values in $`C^{\mathrm{}}(C)[[\lambda ]]`$ with respect to the undeformed representation $`\varrho _C`$.
2. In particular, the component of ghost number zero of the quantum BRST cohomology is isomorphic to the submodule of all those elements in $`C^{\mathrm{}}(C)[[\lambda ]]`$ which are invariant under $`\varrho _C`$. This space being isomorphic to $`C^{\mathrm{}}(M_{\mathrm{red}})[[\lambda ]]`$ the algebra structure on the cohomology induces a star-product on the reduced space $`M_{\mathrm{red}}`$.
For a proof see .
Remarks:
1. The proofs of the last two theorems are rather technical. They heavily rely on one side on purely geometric considerations, namely the existence of tubular neighbourhoods (which can be chosen $`G`$-invariant for proper $`G`$-actions) and the triviality of the normal bundle of $`C`$ in $`M`$ (since $`0`$ is a regular value of $`J`$), which leads to the construction of an acyclic Koszul complex (first on the submodule of $`𝒞`$ of ghost degree zero which is in a standard way extended to all of $`𝒞`$), and a rather explicit chain homotopy for that complex analogous to the one used in the proof of Poincaré’s Lemma. Secondly, we have used a purely tensorial, explicit equivalence transformation which modifies the Clifford part of the multiplication in $`𝒞`$ in such a way that $`Q`$ splits into a boundary operator lowering the antighost degree by $`1`$ and leaving invariant the ghost degree (which turns out to be a deformation of the aforementioned Koszul boundary operator) and a coboundary operator raising the ghost degree by one and leaving invariant the antighost degree (which turns out to be equal to a certain Chevalley-Eilenberg operator of $`𝔤`$). Hence $`𝒞`$ becomes a double complex where one differential is acyclic. This fact has been known in the classical situation, but miraculously remains true in this deformed situation. Thirdly, to relate the total cohomology to the data on the constraint surface $`C`$ we use an augmentation of this complex consisting in a deformation of the restriction map by a formal series of differential operators which can be constructed out of $`Q`$ and the classical chain homotopies. Finally, in the case of a proper group action the resulting star-product on the reduced space can be related to the one on $`M`$ essentially by means of the deformed restriction map.
2. For quantum covariant, but not strongly invariant star-products it can happen that the above mentioned ghost number zero part of the cohomology, the space of ‘quantum $`G`$-invariant functions on $`C`$’, can be too small in the sense that it is no longer a deformation of the whole space of classical $`G`$-invariant functions, but of a subspace of the latter, which is quite an anomaly. In a simple example (see , Section 7) we have seen that the reduced algebra can ultimately become commutative which does no longer seem to resemble a reasonable reduction of quantization, but which –with a little bad luck– in principle is possible as the example shows.
## 4 Classical reducible BRST without ghosts of ghosts
The results of this Section have been obtained in collaboration with Hans-Christian Herbig in .
Let $`C`$ be an arbitrary closed coisotropic submanifold of a symplectic manifold $`(M,\omega )`$ of codimension $`n`$, i.e. the $`\omega `$-orthogonal space to each tangent space of $`C`$ is contained in that tangent space. Physicists would speak of $`C`$ as a ‘first class constraint surface’. Let $`TC^\omega `$ be the $`\omega `$-orthogonal bundle to $`TC`$. This is known to be an integrable subbundle of $`TC`$ and gives rise to a local foliation thanks to Frobenius’ Theorem. If this foliation allows for a smooth quotient manifold $`M_{\mathrm{red}}`$ it becomes a symplectic manifold in a canonical way, see e.g. \[1, p. 417–418\]. Fix a subbundle $`N`$ of $`TM|_C`$ such that $`TM|_C=NTC`$ (e.g. as the normal bundle to $`TC`$ with respect to some Riemannian metric). The symplectic form provides an identification of $`N`$ with the dual of $`TC^\omega `$ via $`v(w\omega (v,w))`$ where $`cC`$, $`vN_c`$ (the fibre of $`N`$ over $`c`$) and $`wT_cC^\omega `$ and an identification of $`TC^\omega `$ with the conormal bundle of $`TC`$, i.e. the subbundle $`TC^{\mathrm{ann}}`$ of $`T^{}M|_C`$ of all those cotangent vectors annihilating $`TC`$ via $`v(w\omega (v,w))`$ where $`cC`$, $`vT_cC^\omega `$ and $`wT_cM`$, whence
$$N^{}TC^\omega TC^{\mathrm{ann}}.$$
(49)
The nontriviality of the bundle $`N`$ (and hence of the two others in the above equation) is related to the physicists’ ‘reducible case’: here the submanifold $`C`$ is given as the zero locus of a finite set of in general not functionally independent smooth real valued functions.
Next, choose a tubular neighbourhood around $`C`$, i.e. an open neighbourhood $`U`$ of the zero-section of $`N`$ together with a diffeomorphism $`\mathrm{\Phi }`$ of $`U`$ onto an open neighbourhood $`V`$ of $`C`$ in $`M`$ such that $`\mathrm{\Phi }(c)=c`$ for all $`cC`$ (where we identify $`C`$ with the zero-section in $`N`$). Hence $`U`$ becomes a symplectic manifold with the pulled-back form $`\mathrm{\Phi }^{}\omega `$. Denoting the bundle projection $`NC`$ by $`p`$ we consider the pulled-back bundle $`p^{}N`$ over $`U`$. We shall denote the dual bundle of $`p^{}N`$ by $`F`$, whence $`p^{}N`$ can be identified with $`F^{}`$. We have made this choice of notation to have an analogy $`M\times 𝔤F`$ and $`M\times 𝔤^{}F^{}`$ with the previous section.
The main idea which will make the construction work is the fact that the bundle $`F^{}(=p^{}N)`$ admits the tautological section $`J`$ which maps each point $`u`$ of $`U`$ to the same point in the fibre over $`p(u)`$. $`J`$ can be seen as a generalization of the momentum map of the previous section. I had been inspired by a similar construction in Connes’s book \[13, p.210\], used for the computation of the Hochschild cohomology of the algebra of all complex-valued $`C^{\mathrm{}}`$-functions on a given manifold $`M`$.
Choosing an arbitrary covariant derivative $`^F`$ in $`F`$ (inducing a covariant derivative $`^F^{}`$ in $`F^{}`$ in the standard way) we set
$$E:=FF^{};^E:=^F+^F^{}$$
(50)
and choose the natural pairing between $`F`$ and $`F^{}`$ as fibre metric $`q`$. It is clear that the above $`^E`$ preserves $`q`$.
Consider now $`𝒞_0:=\mathrm{\Gamma }(\mathrm{\Lambda }F^{}\mathrm{\Lambda }F)`$ together with the Rothstein super-Poisson bracket $`\{,\}_R`$ constructed out of the above data. Define the ghost degree, antighost degree, and ghost number maps in the same way as in the previous section. Then we have the following
###### Theorem 4.1
We use the above-made assumptions. Then
1. The ghost number map $`\mathrm{𝖦𝗁}`$ is a derivation of the super-Poisson algebra $`𝒞_0`$ which thus becomes $``$-graded.
2. There is an element $`\mathrm{\Theta }:=_{i=0}^n\mathrm{\Theta }_i𝒞_0`$, the so-called classical BRST charge, such that $`\mathrm{𝖦𝗁}(\mathrm{\Theta })=1`$, $`\mathrm{\Theta }_0=J`$, the antighost degree of $`\mathrm{\Theta }_i`$ is $`i`$, and, most importantly, $`\{\mathrm{\Theta },\mathrm{\Theta }\}_R=0`$.
3. The classical BRST operator $`Q:=\{\mathrm{\Theta },\}_R`$ has square zero, increases the ghost-number by one, and its classical BRST-cohomology $`\mathrm{Ker}Q/\mathrm{Im}Q`$ carries a canonical $``$-graded super-Poisson algebra structure induced by the one on $`𝒞_0`$.
In order to compute the above cohomology we consider the space of *vertical differential forms on* $`C`$, i.e. the space of sections $`\mathrm{\Omega }_v(C):=\mathrm{\Gamma }(\mathrm{\Lambda }(TC^\omega )^{})`$ together with the vertical exterior derivative $`d_v`$ which is defined by the same formula as the standard exterior derivative but restricted to vertical vector fields, i.e. sections of the integrable subbundle $`TC^\omega `$. Then we have the following result (which should be known by other methods):
###### Theorem 4.2
We use the above-made assumptions and notations. Then the classical BRST-cohomology is isomorphic to the vertical de Rham cohomology, i.e. the cohomology of the complex $`(\mathrm{\Omega }_v(C),d_v)`$. This latter space thus carries the structure of a $``$-graded super-Poisson bracket. Moreover, the sector of the classical BRST-cohomology having vanishing ghost number exactly corresponds to the space of all complex-valued $`C^{\mathrm{}}`$-functions on $`C`$ which are constant on the connected leaves of the foliation defined by $`TC^\omega `$. In case the reduced space $`M_{\mathrm{red}}`$ exists this last space is equal to the space of all complex-valued $`C^{\mathrm{}}`$-functions on $`M_{\mathrm{red}}`$.
For details of the proof, see . The main tool is the fact that $`J`$ defines a Koszul boundary operator on the space $`\mathrm{\Gamma }(\mathrm{\Lambda }F)`$ in the same way as has been remarked in the previous Section, that the resulting complex is acyclic allowing for an augmentation map consisting of the restriction to $`C`$, and that the component of $`\{J,J\}_R`$ having vanishing antighost degree vanishes when restricted to $`C`$ thanks to the fact that $`C`$ is coisotropic and to the chosen connection $`^E`$. In the irreducible case where $`C`$ is given as the zero locus of $`n:=codimC`$ functionally independent functions the method of deforming $`J`$ is well-known, see e.g.
The advantage of the above construction is that it is contained in a simple, geometrically defined BRST-complex $`𝒞_0`$ with only a finite number of nonzero $`𝒞^{i,j}`$ (although it may become difficult to explicitly compute the tubular neighbourhoods) in contrast to the more elaborate multistep ghosts-of-ghosts methods based on spectral sequence techniques, . It is tempting to try a quantization of this complex by means of the Fedosov-type star-product constructed in the first section, but this would require a more sophisticated analysis of the (affine) geometry of the vicinity of $`C`$ to solve the obvious problem whether the component of antighost degree zero of $`JJ`$ vanishes when restricted to $`C`$.
## Acknowledgments
I would like to thank R. Eckel, A. El Gradechi, H.-C. Herbig, N. Neumaier, C. Paufler, S. Waldmann, and in particular the referee for many useful discussions and propositions, for finding the most embarassing typos (as for instance eqn (10)) in my 1996 preprint , and for a critical reading of the manuscript. Moreover I’d like to thank J.-C. Cortet and D. Sternheimer for encouraging me to write this report.
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# Electromagnetic Response of a 𝑘_𝑥±𝑖𝑘_𝑦 Superconductor: Effect of Order Parameter Collective Modes
## Abstract
Effects of order parameter collective modes on electromagnetic response are studied for a clean spin-triplet superconductor with $`k_x\pm ik_y`$ orbital symmetry, which has been proposed as a candidate pairing symmetry for Sr<sub>2</sub>RuO<sub>4</sub>. It is shown that the $`k_x\pm ik_y`$ superconductor has characteristic massive collective modes analogous to the clapping mode in the A-phase of superfluid <sup>3</sup>He. We discuss the contribution from the collective modes to ultrasound attenuation and electromagnetic absorption. We show that in the electromagnetic absorption spectrum the clapping mode gives rise to a resonance peak well below the pair breaking frequency, while the ultrasound attenuation is hardly influenced by the collective excitations.
Superconducting and superfluid properties of unconventional Cooper pairing states continue to be of interest in the research of strongly correlated fermion systems such as liquid <sup>3</sup>He, heavy fermion compounds, high-$`T_c`$ superconductors and the more recently discovered superconductor Sr<sub>2</sub>RuO<sub>4</sub>. The spin-triplet $`k_x\pm ik_y`$-wave state is an example of the unconventional pairing states and is known to describe the superfluidity of <sup>3</sup>He-A. Two-dimensional (2D) version of this state has recently attracted much attention as a candidate pairing state for Sr<sub>2</sub>RuO<sub>4</sub>.
The possibility of the spin-triplet $`p`$-wave pairing in Sr<sub>2</sub>RuO<sub>4</sub> was first suggested by Rice and Sigrist from the analogy to <sup>3</sup>He. There now exist considerable experimental evidences supporting the unconventional spin-triplet superconductivity in Sr<sub>2</sub>RuO<sub>4</sub>. The absence of the Hebel-Slichter peak in $`1/T_1`$ and substantial reduction in $`T_c`$ by non-magnetic impurities indicate that Sr<sub>2</sub>RuO<sub>4</sub> is at least not a conventional $`s`$-wave superconductor. The $`\mu `$SR experiments suggest a pairing state with broken time reversal symmetry. More strong evidence of the triplet pairing has been found by <sup>17</sup>O NMR measurements in which the temperature-independent Knight shift is observed for magnetic fields parallel to the RuO<sub>2</sub> plane. Possible pairing states in Sr<sub>2</sub>RuO<sub>4</sub> have been classified from a group theoretical point of view. It was shown that the quasi-2D electronic structure of Sr<sub>2</sub>RuO<sub>4</sub> leads to five possible $`p`$-wave states stabilized by the absence of gap nodes. Among these states, the pairing state compatible with all the above experiments is the $`k_x\pm ik_y`$ state which is the 2D analog of <sup>3</sup>He-A.
Study of the electromagnetic (EM) response provides valuable information on the properties of unconventional superconductors. Dynamical properties are, in particular, intriguing because they depend not only on the equilibrium gap structure but also on the collective excitations of the order parameter. As is known in the study of superfluid <sup>3</sup>He, internal degrees of freedom of the order parameter give rise to the order parameter collective modes (OPCM’s) specific to a given pairing symmetry. In this paper, we discuss, within the collisionless regime, how the OPCM’s in the $`k_x\pm ik_y`$ superconductor contribute to the EM response.
The OPCM’s in the 2D $`k_x\pm ik_y`$ state has been discussed by Tewordt. As in the case of <sup>3</sup>He-A, there exists the clapping mode, which is a characteristic massive collective mode resulting from the orbital degrees of freedom of the $`k_x\pm ik_y`$-wave order parameter. Tewordt showed that the coupling of this mode to external fields vanishes in the long wavelength limit ($`q0`$) and so he did not attempt to discuss the observability of the clapping mode. The effect of a finite wave vector $`𝐪`$ is, however, important even in type II superconductors when we consider the role of OPCM’s in the EM response. We in fact show in this paper, by calculating a $`q`$-dependent dynamical conductivity, that the clapping mode leads to an EM absorption peak. We also discuss the ultrasound attenuation. We show that, although ultrasound experiments have played a key role in the study of the OPCM’s in <sup>3</sup>He, the ultrasound cannot be a good probe of the OPCM’s in metals such as Sr<sub>2</sub>RuO<sub>4</sub>. This is because the sound velocity in ordinary metals is much smaller than the Fermi velocity $`v_\mathrm{F}`$, in contrast to <sup>3</sup>He.
The 2D $`k_x\pm ik_y`$ state is defined on a cylindrical Fermi surface and the orbital structure of the order parameter is expressed in terms of two basis functions $`\widehat{k}_x=\mathrm{cos}\theta `$ and $`\widehat{k}_y=\mathrm{sin}\theta `$ specifying the direction of the Fermi momentum $`𝐤_\mathrm{F}`$. The matrix order parameter is given by
$$\widehat{\mathrm{\Delta }}_𝐤=\mathrm{\Delta }(\widehat{k}_x\pm i\widehat{k}_y)\sigma _zi\sigma _y=\mathrm{\Delta }e^{\pm i\theta }\sigma _zi\sigma _y,$$
(1)
where $`\sigma _i`$’s are Pauli matrices. The gap $`(\widehat{\mathrm{\Delta }}_𝐤\widehat{\mathrm{\Delta }}_𝐤^{})^{1/2}=\mathrm{\Delta }`$ of the two unitary states is independent of $`𝐤`$ and the Bogoliubov-quasiparticle energy $`E_𝐤=(\xi _𝐤^2+\widehat{\mathrm{\Delta }}_𝐤\widehat{\mathrm{\Delta }}_𝐤^{})^{1/2}`$ is the same as that of the $`s`$-wave state.
The collective excitations of the order parameter are described as oscillations of order parameter fluctuations from the equilibrium state. In the $`k_x\pm ik_y`$ state, one of the two degenerate states is realized in equilibrium. We choose the $`k_x+ik_y`$ state as the equilibrium state and consider an order parameter fluctuation of a plane wave form $`\delta \widehat{\mathrm{\Delta }}_𝐤(𝐪,\omega )e^{i(𝐪𝐫\omega t)}`$. Such fluctuation can be excited by space-time dependent EM fields, a scalar potential $`\varphi (𝐪,\omega )e^{i(𝐪𝐫\omega t)}`$ and a vector potential $`𝐀(𝐪,\omega )e^{i(𝐪𝐫\omega t)}`$. The order parameter fluctuation is expressed using the two bases $`e^{\pm i\theta }=\widehat{k}_x\pm i\widehat{k}_y`$ as
$$\delta \widehat{\mathrm{\Delta }}_𝐤(𝐪,\omega )=[D_+(𝐪,\omega )e^{i\theta }+D_{}(𝐪,\omega )e^{i\theta }]\sigma _zi\sigma _y,$$
(2)
where $`D_+`$ and $`D_{}`$ are variables to be determined by the self-consistency equation. Since $`D_+`$ and $`D_{}`$ are complex variables, there are four degrees of freedom of the order parameter fluctuations in the representation of (2).
A well-established way to study the dynamical properties including the fluctuations of the order parameter is to introduce matrix distribution functions in spin space, $`n_𝐤(𝐪,\omega )`$ and $`f_𝐤(𝐪,\omega )`$ with matrix elements $`(n_𝐤)_{\alpha \beta }=𝑑te^{i\omega t}a_{𝐤_{}\beta }^{}a_{𝐤_+\alpha }`$ and $`(f_𝐤)_{\alpha \beta }=𝑑te^{i\omega t}a_{𝐤_{}\beta }a_{𝐤_+\alpha }`$ ($`𝐤_\pm =𝐤\pm 𝐪/2`$). The deviation $`\delta f_𝐤`$ of $`f_𝐤`$ from its equilibrium value determines the order parameter fluctuation via the self-consistency equation:
$$\delta \widehat{\mathrm{\Delta }}_𝐤(𝐪,\omega )=\underset{𝐤^{}}{}v_{𝐤,𝐤^{}}\delta f_𝐤^{}(𝐪,\omega ),$$
(3)
where $`v_{𝐤,𝐤^{}}=2g_1\widehat{𝐤}\widehat{𝐤}^{}`$ is the pairing interaction. The deviation $`\delta n_𝐤`$ is related to the charge current density as
$$𝐉(𝐪,\omega )=ev_\mathrm{F}\underset{𝐤}{}\widehat{𝐤}\mathrm{Tr}\delta n_𝐤(𝐪,\omega )\frac{ne^2}{mc}𝐀(𝐪,\omega ),$$
(4)
where $`n`$ is the number density. (Note that for the spin-independent perturbations considered here, $`\delta n_𝐤`$ is proportional to the unit matrix and thus the trace in Eq. (4) gives only a factor 2.) The distribution functions have been intensively studied in the context of ultrasound attenuation in superfluid <sup>3</sup>He (see, for a review, Ref. ). Assuming the particle-hole symmetry, we can write the distribution functions integrated over the energy variable $`\xi _𝐤`$ as
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\xi _𝐤\delta n_𝐤`$ $`=`$ $`(\delta ϵ_0+\delta ϵ_1)\eta /(\omega \eta )`$ (5)
$``$ $`2\widehat{\mathrm{\Delta }}_𝐤\widehat{\mathrm{\Delta }}_𝐤^{}(\omega \delta ϵ_0+\eta \delta ϵ_1)/(\omega \eta )`$ (6)
$`+`$ $`{\displaystyle \frac{}{2}}(\omega +\eta )(\delta \widehat{\mathrm{\Delta }}_𝐤\widehat{\mathrm{\Delta }}_𝐤^{}\widehat{\mathrm{\Delta }}_𝐤\delta \widehat{\mathrm{\Delta }}_𝐤^{}),`$ (7)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\xi _𝐤\delta f_𝐤`$ $`=`$ $`{\displaystyle _{\epsilon _c}^{\epsilon _c}}𝑑\xi _𝐤{\displaystyle \frac{\mathrm{\Theta }_𝐤}{E_𝐤}}\delta \widehat{\mathrm{\Delta }}_𝐤`$ (8)
$``$ $`[{\displaystyle \frac{1}{2}}(\omega ^22\widehat{\mathrm{\Delta }}_𝐤\widehat{\mathrm{\Delta }}_𝐤^{}\eta ^2)\delta \widehat{\mathrm{\Delta }}_𝐤`$ (10)
$`\widehat{\mathrm{\Delta }}_𝐤\delta \widehat{\mathrm{\Delta }}_𝐤^{}\widehat{\mathrm{\Delta }}_𝐤(\omega \delta ϵ_0+\eta \delta ϵ_1)\widehat{\mathrm{\Delta }}_𝐤],`$
where $`\delta ϵ_0=e\varphi (𝐪,\omega )`$ and $`\delta ϵ_1=(e/mc)𝐤_\mathrm{F}𝐀(𝐪,\omega )`$ represent the perturbation energies, $`\eta =𝐯_\mathrm{F}𝐪`$, $`\mathrm{\Theta }_𝐤=(1/2)\mathrm{tanh}(E_𝐤/2T)`$, $`\epsilon _c`$ is the usual cutoff energy and
$`={\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\xi _𝐤}{2E_{𝐤_+}E_𝐤_{}}}`$ $`[{\displaystyle \frac{(E_{𝐤_+}+E_𝐤_{})(\mathrm{\Theta }_𝐤_{}+\mathrm{\Theta }_{𝐤_+})}{\omega ^2(E_{𝐤_+}+E_𝐤_{})^2}}`$ (12)
$`+{\displaystyle \frac{(E_{𝐤_+}E_𝐤_{})(\mathrm{\Theta }_𝐤_{}\mathrm{\Theta }_{𝐤_+})}{\omega ^2(E_{𝐤_+}E_𝐤_{})^2}}].`$
Here we have dropped the arguments $`𝐪`$ and $`\omega `$ for brevity ($`\delta \widehat{\mathrm{\Delta }}^{}=\delta \widehat{\mathrm{\Delta }}^{}(𝐪,\omega )`$). In Eqs. (5) and (8), small corrections of order $`q/k_\mathrm{F}`$ have been neglected and accordingly $`\widehat{\mathrm{\Delta }}_{𝐤_\pm }`$ has been put as $`\widehat{\mathrm{\Delta }}_𝐤`$.
Now we discuss the OPCM’s in the $`k_x+ik_y`$ state. It is convenient to introduce the following new variables:
$`{\displaystyle \frac{1}{2}}[D_+(𝐪,\omega )\pm D_+^{}(𝐪,\omega )]=\{\begin{array}{cc}D_+^{}\hfill & \\ D_+^{\prime \prime }\hfill & \end{array},`$ (13)
$`{\displaystyle \frac{1}{2}}[D_{}(𝐪,\omega )e^{2i\theta _q}\pm D_{}^{}(𝐪,\omega )e^{2i\theta _q}]=\{\begin{array}{cc}D_{}^{}\hfill & \\ D_{}^{\prime \prime }\hfill & \end{array},`$ (14)
where $`e^{i\theta _q}=\widehat{q}_x+\mathrm{i}\widehat{q}_y`$. Using Eq. (8), the self-consistency equation (3) can be written in terms of these variables as
$`M_1\left(\begin{array}{c}D_+^{}\\ D_{}^{}\end{array}\right)`$ $`=`$ $`\left(\begin{array}{c}0\\ 2i\mathrm{\Delta }(\omega \delta ϵ_0+\eta \delta ϵ_1)\mathrm{sin}2\vartheta _𝐤\end{array}\right),`$ (15)
$`M_2\left(\begin{array}{c}D_+^{\prime \prime }\\ D_{}^{\prime \prime }\end{array}\right)`$ $`=`$ $`\left(\begin{array}{c}2\mathrm{\Delta }(\omega \delta ϵ_0+\eta \delta ϵ_1)_𝐤\\ 2\mathrm{\Delta }(\omega \delta ϵ_0+\eta \delta ϵ_1)\mathrm{cos}2\vartheta _𝐤\end{array}\right),`$ (16)
where $`\vartheta `$ is the angle between $`𝐤`$ and $`𝐪`$, $`\mathrm{}_𝐤`$ denotes the angle average over the Fermi surface and $`M_{1,2}`$ are $`2\times 2`$ matrices with the following elements:
$`(M_1)_{11}`$ $`=`$ $`(\omega ^24\mathrm{\Delta }^2\eta ^2)_𝐤,`$ (17)
$`(M_1)_{22}`$ $`=`$ $`[\omega ^22\mathrm{\Delta }^2(1+\mathrm{cos}4\vartheta )\eta ^2]_𝐤,`$ (18)
$`(M_1)_{12}`$ $`=`$ $`(M_1)_{21}=(\omega ^24\mathrm{\Delta }^2\eta ^2)\mathrm{cos}2\vartheta _𝐤,`$ (19)
$`(M_2)_{11}`$ $`=`$ $`(\omega ^2\eta ^2)_𝐤,`$ (20)
$`(M_2)_{22}`$ $`=`$ $`[\omega ^22\mathrm{\Delta }^2(1\mathrm{cos}4\vartheta )\eta ^2]_𝐤,`$ (21)
$`(M_2)_{12}`$ $`=`$ $`(M_2)_{21}=(\omega ^2\eta ^2)\mathrm{cos}2\vartheta _𝐤.`$ (22)
Note that the variables $`D^{}`$ and $`D^{\prime \prime }`$ are decoupled from each other.
The OPCM’s correspond to the eigen modes in the absence of external fields. The eigen frequencies $`\omega `$ are therefore determined by $`detM_{1,2}=0`$. In the long wavelength limit, $`q0`$, it is easy to obtain the eigen frequencies. In this limit, since the off-diagonal elements of $`M_{1,2}`$ may be neglected, all the modes decouple. The results for $`q0`$ are summarized in Table I. The $`D_+^{}`$-mode has the eigen frequency just at the pair breaking edge, $`\omega =2\mathrm{\Delta }`$. The variable $`D_+^{\prime \prime }`$ represents a mode corresponding to the phase fluctuation of the order parameter, namely, the so-called Anderson-Bogoliubov mode. This is a gapless mode with sound-like dispersion $`\omega =(v_\mathrm{F}/\sqrt{2})q`$ but is replaced by a plasmon when the Coulomb interaction between electrons is taken into account. The remaining two modes, $`D_{}`$-modes, have the same character as the clapping mode in superfluid <sup>3</sup>He-A. The eigen frequency of the clapping modes in the 2D system is $`\omega =\sqrt{2}\mathrm{\Delta }`$ and is below the pair breaking energy $`2\mathrm{\Delta }`$. At finite $`q`$, the dispersion relation of the clapping modes up to the order $`q^2`$ is given by
$$\omega _{\mathrm{cl}}^2(q)=2\mathrm{\Delta }^2+\frac{1}{2}(v_\mathrm{F}q)^2.$$
(23)
The coupling of the OPCM’s to the external fields is determined by the right-hand side of Eqs. (15) and (16). We see that the $`D_\pm ^{}`$-modes and the $`D_\pm ^{\prime \prime }`$-modes couple to the transverse field and the longitudinal field, respectively.
The presence of the coupling between the $`D_\pm ^{\prime \prime }`$-modes and the longitudinal field means that these modes can be excited by a longitudinal phonon field. Since the collective excitations couple to the density fluctuation $`\delta \rho (𝐪,\omega )=_𝐤\mathrm{Tr}\delta n_𝐤`$ via the last term in Eq. (5), there is a possibility that the clapping mode ($`D_{}^{\prime \prime }`$-mode) is observed by ultrasound measurements. We have estimated the ultrasound attenuation coefficient by assuming typical conditions $`\omega v_\mathrm{F}q\mathrm{\Delta }`$ for ultrasound measurements in ordinary metals and also in Sr<sub>2</sub>RuO<sub>4</sub>. (In Sr<sub>2</sub>RuO<sub>4</sub>, the sound velocity is $`10^5`$ cm/s , the Fermi velocity is $`10^7`$ cm/s and $`\mathrm{\Delta }1`$ K.) The attenuation coefficient from quasiparticle excitations coincides with that for the $`s`$-wave superconductor because of the isotropic gap. The contribution from the collective excitations can be obtained using Eqs. (5) and (16). Taking into account the above conditions and considering the limit $`\omega 0`$, one can find, after straightforward algebra, that the collective excitations give only a small correction of order $`(\omega /v_\mathrm{F}q)^210^4`$ to the quasiparticle contribution. Consequently, the ratio of the superconducting to normal attenuation coefficient in the low-frequency limit $`\omega 0`$ is given by the well-known BCS result $`\alpha _s/\alpha _n=2/(e^{\mathrm{\Delta }/T}+1)`$. Thus sound waves are not suitable probe for the OPCM. It is to be noted that, in contrast to ordinary metals, ultrasound propagation in superfluid <sup>3</sup>He is in the high frequency regime such that $`\omega \mathrm{\Delta }v_\mathrm{F}q`$ where the OPCM plays an important role. This is a reason why sound waves can be a good probe of the OPCM’s in superfluid <sup>3</sup>He.
Now we consider the transverse response of a $`k_x+ik_y`$ superconductor. The transverse response is conveniently described by the complex conductivity $`\sigma _t(𝐪,\omega )`$ for the transverse electric field $`E_t(𝐪,\omega )=(i\omega /c)A_t(𝐪,\omega )`$. The real part $`\sigma _1`$ of the complex conductivity $`\sigma _t=\sigma _1+i\sigma _2`$ determines the absorption of EM waves. We can readily obtain $`\sigma _t`$ by calculating the transverse current density using the above formulation. The conductivity consists of three characteristic parts, i.e., $`\sigma _t=\sigma _t^{\mathrm{qp}}+\sigma _t^{\mathrm{cm}}+\sigma _t^{\mathrm{dia}}`$; the first and the second terms are contributions from quasiparticle excitations and from collective modes, respectively, and the last term $`\sigma _t^{\mathrm{dia}}=ine^2/m\omega `$ arises from the diamagnetic current $`(ne^2/mc)A_t`$. The quasiparticle contribution $`\sigma _t^{\mathrm{qp}}`$ is the same as the BCS result for the $`s`$-wave superconductor:
$$\sigma _t^{\mathrm{qp}}(𝐪,\omega )=\frac{2ine^2}{m\omega }\frac{\widehat{k}_t^2\eta ^2}{\omega ^2\eta ^2}(12\mathrm{\Delta }^2)_𝐤,$$
(24)
where $`\widehat{k}_t=\mathrm{sin}\vartheta `$ is the transverse component of the unit vector $`\widehat{𝐤}`$. The collective-mode contribution $`\sigma _t^{\mathrm{cm}}`$ comes only from the $`D_{}^{}`$-clapping mode and is obtained as
$`\sigma _t^{\mathrm{cm}}`$ $`(𝐪,\omega )=2iev_\mathrm{F}\mathrm{\Delta }{\displaystyle \underset{𝐤}{}}\widehat{k}_t\eta \mathrm{sin}2\vartheta D_{}^{}/E_t`$ (25)
$`=`$ $`{\displaystyle \frac{4ine^2}{m\omega }}{\displaystyle \frac{\mathrm{\Delta }^2(\omega ^24\mathrm{\Delta }^2\eta ^2)_𝐤\eta \widehat{k}_t\mathrm{sin}2\vartheta _𝐤^2}{detM_1}},`$ (26)
where the last line is derived using $`D_{}^{}`$ determined by Eq. (15). Note that Eq. (25) includes $`detM_1`$ in the denominator. This implies that a resonance EM absorption due to the clapping mode occurs at the eigen frequency. To see this explicitly, it is useful to estimate the real part $`\sigma _1`$ by assuming the long wavelength condition such that $`v_\mathrm{F}q\mathrm{\Delta }`$. This assumption is justified for type II superconductors, since the penetration depth $`\lambda `$ is large compared with the coherence length $`\xi _0`$ and the important values of $`q`$ in the Fourier integral of the fields penetrating into the superconductor are $`\lambda ^1`$. In what follows, we discuss the $`\omega `$-dependence of $`\sigma _1`$ in this case.
Since we are interested in high frequencies $`\omega \mathrm{\Delta }`$, we may expand the conductivity in terms of $`q`$. Expanding Eqs. (24) and (25) up to the order $`q^2`$ and adding the results, we obtain
$$\sigma _t(𝐪,\omega )=\frac{ine^2v_\mathrm{F}^2q^2}{4m\omega ^3}\left(1_0\frac{\omega ^24\mathrm{\Delta }^2}{\omega ^2\omega _{\mathrm{cl}}^2}\right)+\sigma _t^{\mathrm{dia}},$$
(27)
where $`_0=\mathrm{\Delta }^2(𝐪=0,\omega )`$. Equation (27) clearly shows that $`\sigma _t`$ has a pole at the eigen frequency of the clapping mode, yielding a delta function peak in the real part $`\sigma _1`$. The diamagnetic term $`\sigma _t^{\mathrm{dia}}`$ in Eq. (27) is pure imaginary except at $`\omega =0`$, so that only the first term contributes to $`\sigma _1`$ at the high frequencies of interest.
We briefly discuss $`\sigma _1`$ in the low frequency region $`\omega \mathrm{\Delta }`$ where the OPCM does not come into play. There are two contributions to $`\sigma _1`$ at the low frequencies. One is a delta function at $`\omega =0`$, associated with superfluid motion excited by EM fields. The other is caused by thermally excited quasiparticles. The energy conservation $`\omega =|E_{𝐤_+}E_𝐤_{}|`$ for the corresponding quasiparticle-scattering process requires that this contribution is restricted in the region $`\omega <v_\mathrm{F}q\mathrm{\Delta }`$. We note that in the limit $`v_\mathrm{F}q/\mathrm{\Delta }0`$, the conductivity sum rule $`_0^{\mathrm{}}𝑑\omega \sigma _1(𝐪,\omega )=\pi ne^2/2m`$ is satisfied only by the two contributions. As can be shown from Eq. (24), in the limit $`v_\mathrm{F}q/\mathrm{\Delta }0`$, the $`\omega `$-integration yields $`\pi n_se^2/2m`$ for the $`\delta (\omega )`$ term and $`\pi n_ne^2/2m`$ for the contribution from thermally excited quasiparticles, where $`n_s`$ is the superfluid density and $`n_n=nn_s`$ is the normal-fluid density. This explains why $`\sigma _1`$ at $`\omega \mathrm{\Delta }`$, given by Eq. (27), is proportional to $`q^2`$.
Equation (27) depends on temperature $`T`$ through $`_0`$ and $`\mathrm{\Delta }`$. At sufficiently low temperatures, the hyperbolic tangent function, $`\mathrm{tanh}(E_𝐤/2T)`$, in $`_0`$ may be replaced by unity; then Eq. (27) is reduced to the result for $`T=0`$. Figure 1 shows the $`\omega `$-dependence of the real part $`\sigma _1`$ at $`T=0`$ obtained from Eq. (27). The structure above $`\omega =2\mathrm{\Delta }`$ is caused by pair-breaking processes. Below the pair-breaking edge, $`\omega <2\mathrm{\Delta }`$, there is a delta function peak at $`\omega _{\mathrm{cl}}=\sqrt{2}\mathrm{\Delta }`$, associated with the absorption of the clapping mode. This collective-mode contribution is explicitly given by $`\sigma _1^{\mathrm{cm}}(𝐪,\omega )=(\pi ^2ne^2v_\mathrm{F}^2q^2/64m\mathrm{\Delta }^2)\delta (\omega \omega _{\mathrm{cl}})`$.
Hirschfeld et al. have studied the power absorption from the OPCM in the Balian-Werthamer state (in a 3D system). The power absorption $`P(\omega )`$ is given as the dissipation of the field energy which appears as Joule heat. To consider $`P(\omega )`$, we need to take into account the presence of the vacuum-metal interface. They estimated $`P(\omega )`$ by assuming the specular surface boundary condition and by neglecting the surface scattering effect causing the pair breaking near the surface. In the 2D $`k_x\pm ik_y`$ state considered here, the same calculation is repeated by using Eq. (27) and considering the EM wave injected along the 2D plane. The resulting collective-mode contribution to $`P(\omega )`$ has the following $`\omega `$-dependence: $`P^{\mathrm{cm}}(\omega )=0`$ below $`\omega =\omega _{\mathrm{cl}}(0)`$. Above the threshold, $`\omega >\omega _{\mathrm{cl}}(0)`$, $`P^{\mathrm{cm}}(\omega )`$ increases rapidly with $`\omega `$ and then has a peak structure with a finite width $`(\xi _0/\lambda _L)\mathrm{\Delta }`$ ($`\lambda _L=(mc^2/4\pi ne^2)^{1/2}`$, the London penetration depth); this $`\omega `$-dependence is descried by $`P^{\mathrm{cm}}(\omega )[\omega ^2\omega _{\mathrm{cl}}^2(0)]^{1/2}/[\omega ^2\omega _{\mathrm{cl}}^2(0)+c_1^2\lambda _L^2]^2`$, where $`c_1=v_\mathrm{F}/\sqrt{2}`$ is the velocity determining the dispersion of the clapping mode. In the power absorption spectrum, the OPCM gives rise to not a delta function peak but a broadened peak. This is a consequence of taking into account the dispersion of $`\omega _{\mathrm{cl}}`$ in the denominator of Eq. (27). In actual metals, impurity scattering also brings about the broadening. The power absorption spectrum, however, has a definite peak structure due to the collective excitation if $`\xi _0`$ is small enough compared with $`\lambda _L`$ and at least the superconductor is clean. This demonstrates that the clapping mode in the $`k_x\pm ik_y`$-wave type II superconductor can be observed by EM absorption measurements.
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# X-ray Spectra of a large sample of Quasars with ASCA
## 1 Introduction
Discovered in 1963 (Schmidt 1963), quasars are the most luminous continuously emitting objects in the Universe and represent the high luminosity end of the class of objects known as Active Galactic Nuclei (AGN). Like their lower luminosity cousins - Seyfert 1 galaxies - the bulk of the energy produced in quasars is thought to arise from accretion onto a compact object (the putative super-massive black hole). This central engine is also thought to be where the X-rays, that are observed from both quasars and Seyfert 1s, originate from.
In one model, the UV photons produced by viscous dissipation in an accretion disk are Comptonised to X-ray energies by a hot corona above the surface of this disk (Haardt & Maraschi 1993). These hard X-rays in turn illuminate the accretion disk, being either ‘reflected’ back towards the observer or thermalised into optical or UV photons. Evidence for these ‘reflection’ features (in the form of an iron K$`\alpha `$ emission line, Fe K absorption edge and Compton scattering hump) is commonly observed in the X-ray spectral band in Seyfert 1 galaxies (e.g. Pounds et al. 1990, Nandra & Pounds 1994). The detection of these reflection features however, in quasars, remains more ellusive (Reeves et al. 1997, Lawson & Turner 1997). In the radio-loud quasars the situation is somewhat further complicated by the presence of a powerful relativistic radio-jet. In the X-ray band, these radio-loud quasars have flatter X-ray spectral emission (e.g. Wilkes & Elvis 1987, Lawson et al. 1992), and are generally more luminous than the radio-quiet quasars. The radio-loud quasars also have little or no X-ray (iron) line emission; this is often interpreted in terms of the Doppler boosting of the X-ray continuum, by the relativistic jet (see Reeves et al. 1997 and references therein).
This paper presents the results of a detailed spectral analysis of 68 quasars obtained from the ASCA public archive. The aims are to extend the results that were presented in Reeves et al. (1997), which contained a smaller sample of 24 objects. The objects considered in this sample contain a roughly equal mix of both radio-loud and radio-quiet quasars. The bigger sample, for instance, allows us to perform an investigation of the properties of iron K lines and reflection in quasars, which will be limited in some of the objects by signal-to-noise. It also permits us to investigate the properties of quasars over a large range of luminosity and redshift. A further aim of the paper is to make the results of this analysis available to the general community; the paper presents results from a considerable number of quasars that are currently unpublished.
In the following section, the selection and properties of the sample are discussed. Section 3 then outlines the spectral fitting that was performed on the quasars in the sample. The following sections (4-7) then present and discuss the results in terms of the X-ray continuum emission, soft X-ray excesses from quasars, the properties of the iron line and reflection associated with the putative accretion disk and also the effects of absorbing material on the quasar spectra. Values of $`H_0=50`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and $`q_0=0.5`$ are assumed and all fit parameters are given in the quasar rest frame.
## 2 The ASCA sample of Quasars
Data have been selected mainly from the ASCA public archive, using all quasars that were available up until January 1998. The observations included in this paper are shown in table 1; in general when there has been a multiple observation of a given quasar we have taken the first observation. In total 68 quasars have been included in the sample, with the objects selected predominantly type I AGN. This covers a wide range of redshift (from z=0.06 to z=4.3), and also a wide range of luminosities (M$`{}_{V}{}^{}=23`$ to -30 and L$`{}_{210keV}{}^{}10^{43}`$ erg s<sup>-1</sup> to $`>`$10<sup>47</sup> erg s<sup>-1</sup>). Note that redshift and luminosity cut-offs of z $`>0.05`$ and M$`{}_{V}{}^{}<23.0`$ respectively have been used to define the sample. The distribution of quasar redshift for this sample is included as figure 1. Note that at z $`>3`$, the sample is predominantly made up of bright core-dominated radio-loud quasars.
Standard (to conservative) ASCA screening criteria and data selection have been used in analysing the data, this leaves 62 quasars that have sufficient signal-to-noise for further spectral analysis and interpretation. The selection criteria used in reducing these observations are outlined in Reeves et al. (1997). Of these 62 quasars, 35 are radio-loud and 27 are radio-quiet; according to the definition of radio-loudness (R<sub>L</sub>) used by Wilkes & Elvis (1987), where R<sub>L</sub> = log (F<sub>5GHz</sub>/F<sub>B</sub>). The cut-off between radio-quiet and radio-loud is defined arbitrarily at R=1. The roughly even numbers of objects in both classes allows us to investigate the properties within both.
Some objects, as described above, have been left out of the subsequent analysis. The radio-quiet object IRAS P09104+4109 may be a buried quasar, i.e. surrounded by a very high column of material; however most of the X-ray emission is thought to originate from a surrounding cluster and X-ray luminous cooling flow (Fabian & Crawford 1995). Similarly the RQQ E 0015+162 is strongly contaminated by the powerful X-ray emitting cluster CL 0016+16 (e.g. Neumann & Bohringer 1997) and the spatial resolution of ASCA is not sufficient to resolve these two objects. Therefore further analysis of these 2 sources have been excluded from our sample. Finally spectra have only been extracted where the source is detected at the 5$`\sigma `$ level in the ASCA detectors. This excludes the radio-quiet quasars QSO 0215-504, MS 12487+5706 and QSO 1725+503 from further study. The radio-quiet broad absorption line (BAL) quasar PHL 5200 has only been detected in the GIS3 and SIS0 detectors, but not GIS2 nor SIS1. As the data from this quasar contain few counts with which to constrain the spectral form, this object has also been excluded from further study in the sample.
## 3 X-ray Spectral Fitting
The ASCA data reduction process provides four seperate spectra, one for each of the four (GIS and SIS) instruments. Spectral fitting was performed by simultaneously fitting the data from each instrument, allowing the relative normalisations to vary as necessary. The processed spectra were fitted with a spectral model using the X-ray spectral fitting software XSPEC v10.0, over the energy range 0.8-10 keV for the GIS instruments and from 0.6-10 keV for the SIS.
The standard model used consists of a power-law with neutral absorption, together with a Gaussian-shaped narrow ($`\sigma =0.01`$ keV) Fe K line at 6.4 keV in the rest frame of the quasar. The absorption column that was fitted has two components; (i) the Galactic column density which was fixed in value in the z=0 frame, and (ii) the intrinsic absorption of the quasar in its rest frame, which was allowed to vary in value. The Galactic column densities have been obtained from the getnh program, run under the xanadu system, which uses data from Stark et al. (1992) and other surveys; absorption cross-sections have been taken from Morrison & McCammon (1983). Where possible more accurate values of $`N_\mathrm{H}`$ towards a particular quasar have been taken from individual observations using either Elvis et al. (1989) or Murphy et al. (1996), which use higher spatial resolution.
The results from the spectral fitting are presented in table 2. The flux is given over the range 0.5 keV to 10 keV (in the observed ASCA frame and extrapolated where necessary), with luminosity quoted from 2-10 keV (corrected for absorption) in the quasar rest-frame. The value quoted for $`N_\mathrm{H}`$ is the measured value of intrinsic absorption in the quasar rest frame, in excess of any Galactic absorption. The equivalent width (EQW) and flux of the iron fluorescence line are given, with the EQW being in the quasar rest frame. The line parameters are those corresponding to a narrow line ($`\sigma =0.01`$ keV); the line energy is included as a free parameter in the model fit unless otherwise stated in the table. Note also that the photon indices $`\mathrm{\Gamma }`$ are those of the underlying power-law continuum, with spectral features such as low energy absorption (cold or ionised), soft X-ray excess, iron line emission or Compton reflection, taken into account where these features are significant. When the fitted spectrum is particularly complex, the 2-10 keV band index has been used. However in most cases there is little difference between the underlying 0.6-10 keV index and the 2-10 keV index. All errors in the table are quoted at 68% confidence, for the appropriate number of interesting parameters.
For each quasar the best fitting value for reduced chi-squared ($`\chi ^2`$/degrees of freedom) is given. The values of $`\mathrm{\Delta }\chi ^2`$ quoted in table 2 are for adding an additional spectral feature to the fit, such as an intrinsic absorption column or Fe line. For adding an absorption column, 1 additional parameter is added to the fit and for an iron line, 1 or 2 parameters are added, depending on whether the line energy is fixed at 6.4 keV or not (i.e. the line normalization and line energy). As a rough guide, $`\mathrm{\Delta }\chi ^2=2.7`$ corresponds to 90% significance for the addition of 1 interesting parameter, with $`\mathrm{\Delta }\chi ^2=4.6`$ corresponding to 90% significance for 2 additional interesting parameters. Only measurements of $`N_\mathrm{H}`$ and the iron line which are at 90% significance or better have been quoted as best fit values in table 3.2; otherwise upper-limits only are quoted.
F-tests have also been performed (see Bevington & Robinson 1992), both on the intrinsic absorption and the iron line, to test the significance of these features. If the value obtained for F $`>3.0`$, for the addition of 1 extra interesting parameter, then the result has significance greater than $``$90%. This is used as a criterion for including line or absorption features in the spectral fitting. Actual spectral features such as an Fe line or an absorption column are only regarded as significant if these F-test criteria are met.
## 4 The X-ray Continuum Emission from Quasars
### 4.1 General Properties
The X-ray emission in the sample covers a wide range of photon index, from hard (e.g. PKS 2251+113, $`\mathrm{\Gamma }=0.95`$) to soft (e.g. IZWI, $`\mathrm{\Gamma }=2.37`$). The distribution of 2-10 photon index for the quasars in the sample is illustrated in figure 2 (plotted against radio-loudness). However the X-ray emission from the majority of quasars lies in the region from $`\mathrm{\Gamma }=1.52.1`$. The mean photon index for all 62 quasars is $`\mathrm{\Gamma }=1.76\pm 0.04`$, with a measured sample dispersion (1$`\sigma `$) of $`0.27\pm 0.03`$. Compared with typical measurement errors of the order $`\mathrm{\Delta }\mathrm{\Gamma }=\pm 0.05`$ to 0.10, this shows that the dispersion in the X-ray continuum emission of quasars is significant.
Spearman-Rank correlations have therefore been performed on the photon index ($`\mathrm{\Gamma }`$), for the 62 quasars. The strongest correlation present is a negative trend between $`\mathrm{\Gamma }`$ and radio-loudness R<sub>L</sub>, which is significant at $`>`$99.99% confidence (see figure 2), confirming the results found previously (e.g. Wilkes & Elvis 1987). The X-ray photon index $`\mathrm{\Gamma }`$ is also observed to decrease with X-ray luminosity L<sub>X</sub> (measured in the 2-10 keV band, QSO rest frame) and the quasar redshift z. However these correlations are of weaker significance (at 99.7% and 95.2% confidence respectively) and probably result from the strong correlation between $`\mathrm{\Gamma }`$ and R<sub>L</sub>. Indeed a partial S-R test performed on the data confirms that the trend between $`\mathrm{\Gamma }`$ and R<sub>L</sub> is the significant correlation present. In addition R<sub>L</sub> and L<sub>X</sub> are strongly correlated ($`>`$99.99%), in the sense that the most radio-loud objects tend to be more X-ray luminous. These correlations are therefore consistent with the increasing importance of the jet, in the more core-dominated radio-loud quasars, due to Doppler boosting, which may account for the flatter X-ray slopes and enhance the luminosity of the RLQs as a whole.
### 4.2 The Radio-Quiet sub-sample
We have calculated that the mean photon index for the radio-quiet quasars in this sample is $`\mathrm{\Gamma }=1.89\pm 0.05`$, with a significant dispersion of $`\sigma =0.27\pm 0.04`$. This is consistent with the mean 2-10 keV index found from an optically selected sample of PG quasars (George et al. 2000). Given this dispersion, correlations within just the radio-quiet sub-sample have been investigated, excluding all the radio-loud objects. The motivation for this is that it is possible to discount any effects from a relativistic jet, and hence just learn about the properties of the quasar central engine. No significant correlations were found between $`\mathrm{\Gamma }`$ and either L<sub>X</sub> or z for the 27 radio-quiet quasars in the sample. Note that the spectral indices were taken over the 2-10 keV band, in the quasar rest-frame.
Thus to extend the range of luminosity present, we next include the Seyfert 1s published in the Nandra et al. (1997) sample; the luminosity range covered then extends from L<sub>X</sub> $``$ 10<sup>41</sup> erg/s for the least luminous Seyfert 1 to L<sub>X</sub> $``$ 10<sup>47</sup> ergs/s for the most luminous quasars. However it was found that there was still no significant correlation between the X-ray photon index with either luminosity or redshift. This is in contrast to the results of the Reeves et al. (1997) paper, which found a positive correlation between $`\mathrm{\Gamma }`$ and L<sub>X</sub>, but which only considered a small sample of 9 radio-quiet quasars. Thus this suggests that there is little evolution in the underlying X-ray emission from radio-quiet quasars with either luminosity or redshift, over a wide range of luminosities and therefore presumably black hole masses. We have also carried out a direct comparison between our sample of radio-quiet quasars and the sample of Seyfert 1s analysed by Nandra et al. (1997). It is found that the mean slope for the radio-quiet quasars ($`\mathrm{\Gamma }=1.89\pm 0.05`$) is almost identical to the mean Seyfert 1 slope (i.e. $`\mathrm{\Gamma }=1.86\pm 0.05`$). This also suggests that there is little difference in the underlying X-ray emission between the Seyfert 1s and the radio-quiet quasars.
### 4.3 The Radio-Loud sub-sample
The radio-loud quasar sub-sample has a mean index of $`\mathrm{\Gamma }=1.66\pm 0.04`$ and dispersion $`\sigma =0.22\pm 0.03`$. Thus the difference in spectral slope between the RLQs and RQQs is significant at $`>`$99% confidence. We also checked for any correlation involving $`\mathrm{\Gamma }`$ within just the radio-loud sub-sample. No significant correlations were found for this sub-sample at the 99% confidence level.
### 4.4 On The Correlation between $`\mathrm{\Gamma }`$ and H$`\beta `$ FWHM
What is interesting is that significant dispersion in $`\mathrm{\Gamma }`$ is present in our sample within the radio-quiet quasar class alone, that cannot be attributed to the properties of the radio-jet and does not seem dependent on the object luminosity (see above). Therefore the radio-quiet sample was split up into broad and narrow line objects (i.e. from the widths of the permitted optical lines, such as H$`\beta `$), whereby the narrow line objects are defined to have H$`\beta `$ FWHM $`<2000`$ km/s (e.g. Osterbrock & Pogge 1985). This gives a sample of 8 ‘narrow-lined’ QSOs. It is found that the sample of 8 narrow-line QSOs have 2-10 keV spectra that are significantly steeper ($`\mathrm{\Gamma }=2.14\pm 0.07`$) than the other normal broad-lined radio-quiet QSOs ($`\mathrm{\Gamma }=1.81\pm 0.05`$). This is analogous to the difference between narrow and broad-line Seyfert 1s found previously in the 2-10 band (e.g. Brandt et al. 1997, Vaughan et al. 1999), and also in soft X-rays (Laor et al. 1997). Additionally most of the narrow line QSOs have strong soft excesses apparent in their spectra (see later).
So to investigate this trend further, we performed Spearman-rank correlations between $`\mathrm{\Gamma }`$ and H$`\beta `$ FWHM on all the radio-quiet quasars in this sample with known H$`\beta `$ widths (21 objects). The correlation was performed using the 2-10 keV photon indicies, fitted in the quasar rest-frame. It was found that the photon index for quasars increases with decreasing H$`\beta `$ width; this yields a correlation coefficient of r=-0.61, significant at the 99.4% confidence level for our 21 radio-quiet quasars. In order to extend the number of objects and also to increase the range of luminosity, we next include both the broad and narrow-lined Seyferts from the samples of Nandra et al. (1997) and Vaughan et al. (1999). The total sample of objects have subsequently been split into 2 bins consisting of (i) low luminosity AGN (24 objects, L$`{}_{210}{}^{}=10^{42}10^{44}`$ erg/s) and (ii) high luminosity quasars (25 objects, L$`{}_{210}{}^{}=10^{44}10^{46}`$ erg/s). The results of the Spearman-Rank analysis is summarised in table 3 and the correlation is plotted in figure 4.
Firstly for all of the 49 radio-quiet objects, we find a strong negative correlation between $`\mathrm{\Gamma }`$ and H$`\beta `$ at $`>99.9\%`$ significance, consistent with the results found in previous samples (e.g. Boller, Brandt & Fink 1996, Brandt et al. 1997, Leighly 1999). Interestingly when the 2 different luminosity sub-samples are used, a significant correlation is not found within the low luminosity bin, but a strong correlation is found for the high luminosity objects (at $`>`$99.9% significance). There does not seem an obvious reason to explain the lack of a correlation at low luminosites, this may just be a selection effect. Also note that the flattest radio-quiet quasar in this sample (PHL 909, with $`\mathrm{\Gamma }1.1\pm 0.1`$) also has the widest H$`\beta `$ line profile (FWHM = 11000 km/s) of all the objects considered. However the removal of PHL 909 makes no difference to any of the correlations. The important finding here is that the apparent correlation between $`\mathrm{\Gamma }`$ and H$`\beta `$ FWHM, found previously in low luminosity samples of Seyfert galaxies, appears to extend to higher luminosities and hence the quasars in this sample. Thus in this sample the ‘narrow-line’ quasars tend to have steeper underlying (2-10 keV) X-ray spectra than the more common broad-lined quasars.
So a tempting hypothesis exists that can explain the dispersion in quasar spectral slope, whereby the photon index is negatively correlated with H$`\beta `$ FWHM. One explanation already postulated is that these softer ‘narrow-line’ objects are accreting at a higher fraction of the Eddington-limit, which leads to greater Compton-cooling of the hard X-ray emitting corona (Pounds, Done & Osborne 1995, Laor et al. 1997). This can naturally account for the steep hard X-ray power-law and strong soft excesses (associated with the disk) generally observed in this class of object. The other interesting finding in the last section was the lack of a correlation between the photon index and the X-ray luminosity, for the radio-quiet quasar sub-sample. The implication here is that the underlying X-ray emission from radio-quiet AGN does not depend on the black hole mass (hence the null correlation between $`\mathrm{\Gamma }`$ and luminosity). Instead the important factor may be the fractional accretion rate ($`\dot{m}`$), i.e. the ratio of the black hole mass accretion rate to the Eddington-limited rate. Thus the underlying X-ray spectra could depend on the fractional accretion rate $`\dot{m}`$, whereby the objects with steeper X-ray spectra are accreting at a higher fraction of the Eddington limit.
## 5 The Soft X-ray Excess
We have searched systematically for a soft excess in the spectra of all the low redshift quasars in this sample. Firstly the spectra were fitted in the harder 2-10 keV band, and then the spectra were extrapolated back to 0.6 keV (or 0.8 keV for the GIS) to see if there was any spectral curvature at lower energies. An excess of counts below 2 keV, above that of the hard power-law continuum, indicates the presence of a ‘soft excess’. In this case we have added a blackbody component (in addition to the power-law and other spectral components - such as an iron line) to parameterize the soft excess emission and have subsequently re-fitted the spectrum. The results of this fitting are shown in table 4. The strength of the soft excess is given by the parameter R, which represents the ratio of the blackbody to power-law component in the 0.6-2 keV range; also given are the temperature of the blackbody kT (units eV) and the improvement ($`\mathrm{\Delta }\chi ^2`$) in the spectral fit from adding the blackbody component to the previous best-fit model (2 additional free parameters). Note that we have restricted this fitting procedure to objects of redshift z$`<`$0.3, as at higher redshifts the soft excess will be shifted out of the ASCA bandpass.
The table shows that some soft excess emission is significant in 9 quasars, approximately half of the low z (z$`<`$0.3) objects. This soft X-ray emission is illustrated in figure 5, which shows the data-model ratio residuals to power-law fits in 2 objects, the broad-line QSO HE 1029-1401 and the narrow-line object Markarian 478. In one object (HE 1029-1401) there is only a gradual curvature of the spectrum below 2 keV, whereas the soft excess in Mrk 478 rises sharply above the power-law continuum. Interestingly the majority of the soft excess emission occurs in the narrow-line quasars (6 objects: TON S180, NAB 0205+0204, PKS 0558-504, PG 1211+143, PG 1404+226, Mrk 478), with soft excesses only apparent in 3 broad-line quasars. This would indicate that soft excesses are more common in the narrow-line objects, which may be expected if these objects indeed accrete at a higher rate for a given black hole mass (e.g. Ross, Fabian & Mineshinge 1992). Only one of the nine objects is radio-loud (PKS 0558-504).
The interesting question to ask here is what is the origin of the soft X-ray excess emission? The standard explanation is that it results from thermal emission that originates directly from the hot inner accretion disk (Malkan & Sargent 1982) and hence is the high energy tail of the so-called ‘Big Blue Bump’. The temperature of this soft excess component for the objects in this sample varies between 100 and 300 eV. In some cases, such as in the narrow-line objects PG 1404+226 and Markarian 478, the soft excess is quite steep, with temperature $`100`$ eV, perhaps consistent with emission from the disk for such objects. In particular the soft excess of PG 1404 is very strong, energetically dominating the X-ray power-law component; this does seem to suggest that the soft excess in PG 1404+226 is a primary emission component and probably does originate directly as thermal emission from the accretion disk.
However other cases may not be so straightforward. For instance the soft excess in the quasar HE 1029-1401 has a blackbody temperature of $``$300 eV and it is likely that temperatures of this order are perhaps too hot to result directly from the quasar accretion disk. Note that in a typical quasar, with a central black hole mass of $`10^8`$M, temperatures of $`<50`$ eV would be expected, implying that the observed soft excesses in many objects are probably too hot to be the direct emission from the putative disk. Furthermore the temperature of the disk component varies as $`T_{BB}`$M$`{}_{}{}^{0.25}{}_{BH}{}^{}`$ and hence is expected to be cooler for the more luminous objects with larger black hole masses.
Nevertheless it is plausible that some degree of Comptonisation by electrons in a hot corona could upscatter cooler EUV photons from the disk to soft X-ray energies and account for the observed emission. Another possibility is that in some objects the soft X-ray emission results from reprocessing of the hard X-ray power-law, for instance through reflection and scattering of X-rays off the optically thick disk (e.g. George & Fabian 1991). Reflection will become increasingly important at soft X-ray energies as the accretion rate $`\dot{m}`$ increases towards the Eddington limit (e.g, Ross & Fabian 1993). As the accretion rate rises, progressively heavier elements become fully ionised and the disk becomes more reflective at soft X-ray energies, which can produce a steepening of the X-ray spectrum. Furthermore ionised emission lines can be produced from abundant elements such as O, Ne and Mg, which could also contribute towards the X-ray flux near to 1 keV. Tentative evidence for this is present in 2 narrow-line objects; in Ark 564 (Vaughan et al. 1999b) an ionised reflection component contributes significantly to the X-ray flux below 2 keV, whilst in the narrow-line QSO TON S180, Turner et al. (1998) report evidence for an emission-like feature at $``$ 1 keV. We also note that some of the quasars in this sample (for instance TON S180, HE 1029-1401, PKS 0558-504) may exhibit ionised iron K emission lines (see section 6), which could provide further support of this hypothesis.
## 6 The Iron K$`\alpha `$ Fluorescence Line
### 6.1 General Properties
Table 5 shows the list of quasars for which Fe fluorescence line emission was found to be significant (i.e. at the 90% level or better). The spectral fitting of the iron K lines was carried out where possible over the 2-10 keV ASCA energy range. The line parameters are first shown for a narrow line ($`\sigma =0.01`$ keV), with energy fixed at 6.4 keV, and then again for a narrow line but with the line energy as a free parameter in the fit. We have also considered broad line fine fits, where the intrinsic velocity width of the line is a extra free parameter. However the constraints on the line width is poor in most of the objects and so this is not considered further. The line parameters have all been fitted in the quasar rest frame, i.e. corrected for redshift effects.
The F-test and F-distribution (Bevington & Robinson 1992) were used to test the significance level of any line features, the results of which are also shown in the table. The change in $`\mathrm{\Delta }\chi ^2`$ for a fit with narrow line fixed at 6.4 keV, compared with a fit with no Fe line, is quoted along with the associated probability from performing an F-test with 1 additional parameter. Then the fit to a narrow line with the line energy free is compared to one with no lines, the values for both $`\mathrm{\Delta }\chi ^2`$ and the F-test probability (for 2 additional parameters; the line normalisation and the line energy) are both quoted. Also quoted is the additional probability (and $`\mathrm{\Delta }\chi ^2`$) for freeing the line energy in the fit; i.e. the fit to a line with free energy is compared with a line fixed at 6.4 keV. This then can be used as a measure of whether the line energy is significantly different to the neutral value of 6.4 keV. Line emission is deemed significant if the F-test probability is $`>`$ 90%; in some cases lines are not significant at 6.4 keV, but are when the line energy is free to vary, indicating that the emitted line energy is different to 6.4 keV.
Iron line emission has been detected in 21 out of the final 62 objects in the sample. There is a large spread in the best-fit value of line equivalent width (EQW), ranging from 32eV for 3C 273 up to $`>`$400eV for the radio-quiet objects IZwI and PG 1404+226, although the main peak of quasars with detected emission between 50 and 200eV is apparently consistent with Compton reflection origins (e.g. George & Fabian 1991). In the case of the radio-quiet quasars there is evidence of Fe K$`\alpha `$ line features, with 14 out of 27 objects showing evidence for line emission above the 90% confidence level. At the 99% (or better) confidence level, Fe K emission lines are detected in IZwI, TON S180, HE 1029-1401, PG 1114+445, PG 1116+215, PG 1211+143, Mrk 205 and E 1821+643. Interestingly lines are not detected in the most luminous (L$`{}_{X}{}^{}>10^{46}`$ ergs/s) of the radio-quiet quasars at higher redshifts; in particular PG 1634+706 ($`<`$35eV), HE 1104-285 ($`<`$53eV), PG 1407+265 (EQW$`<`$70eV) and PG 1718+481 ($`<`$88eV).
In contrast to the radio-quiet quasars, significant line emission has been detected in only 7 of the 35 radio-loud quasars. Generally, the line EQW found here is smaller than that for the radio-quiet objects, with EQW$`<`$100eV in many cases. Where no line emission has been detected, tight upper limits have been placed on the objects in several cases. For example, the blazars 3C 279 and CTA 102 have line EQWs $`<`$28eV and $`<`$32eV respectively, whilst the distant radio-loud quasars PKS 2126-158 and PKS 2149-306 also have tight limits of $`<`$32eV and $`<`$22eV respectively. Most of these quasars with low upper-limits on the line EQW are core-dominated quasars.
As a whole we find that the mean line equivalent width (detections only, narrow line fits) in the RQQs is $`163\pm 17`$ eV, whereas for RLQs the mean is $`85\pm 15`$ eV. This suggests that the amount of iron line emission is weaker in the radio-loud objects. To investigate this further we performed a Spearman-rank correlation (using survival statistics - see Isobe, Feigelson & Neilson 1986), which showed that the line EQW (for narrow lines) decreases with quasar radio-loudness at 99.99% confidence. This trend is also illustrated in figure 6, the mean iron line EQW (including upper limits) has been plotted against radio-loudness, with the quasars split up into different bins for R<sub>L</sub>. As a whole Fe lines are detected in the radio-quiet objects, but not in the radio-loud quasars. The simple interpretation of this trend is that the Fe emission diminishes as the jet angle approaches the line of sight. Thus the increased Doppler boosting of the X-ray continuum in relation to the line, weakens the relative strength of the disk reflection component in core-dominated RLQs.
We have also investigated any trends in the line emission for the radio-quiet objects, excluding the radio-loud quasars. This negates the effect that the powerful relativistic jet has on any of the correlations. In order to extend the range in luminosity, we include the 18 Seyfert 1s from the Nandra et al. (1997) sample, as well as our own 27 RQQs. A Spearman-rank analysis shows that the line EQW is negatively correlated with X-ray luminosity at 99.98% confidence, even when only the radio-quiet AGN are considered. This correlation confirms the ‘X-ray Baldwin effect’ that was found previously from Ginga (Iwasawa & Taniguchi 1993), and also ASCA data (Nandra et al. 1997b). Thus the Fe line emission appears to be absent in the most luminous of AGN, regardless of whether they are radio-loud or quiet.
### 6.2 Iron Line Energy
Overall it is found that in 11 out of 21 of the quasars that show Fe K$`\alpha `$ line emission, the line emission is at rest energies $`>`$6.4 keV at 90% confidence or better. Thus in at least half of the quasars in the sample, the line originates from matter that is partly ionised. For the RQQs, 7 objects have partially ionised Fe K$`\alpha `$ line emission; these are IZwI, TON S180, HE 1029-1401, PG 1116+215, IRAS 13349+2438, E 1821+643 and MR 2251-178. For the radio-loud quasars, 3C 109.0, PKS 0558-504, 3C 273 and PKS 1510-089 appear to have Fe lines that originate from partially ionised material.
The overall range of line energy, from 6.4 keV to 6.9 keV, represents emission from a variety of ionisation states from neutral iron (Fe i to Fe xvi at $``$6.4 keV) up to helium or even hydrogen-like iron (Fe xxv at 6.68 keV, Fe xxvi at 6.96 keV). This is in contrast to the line properties of lower luminosity AGN such as Seyfert1 galaxies, where the Fe K$`\alpha `$ emission normally originates from neutral iron, with the line energies closely distributed near to 6.3-6.4 keV (Nandra et al. 1997). For instance the mean line energy of the Nandra et al. (1997) sample of Seyfert 1 galaxies is $`6.37\pm 0.02`$ keV, whereas for our quasar sample the mean line energy is 6.62$`{}_{0.07}{}^{}{}_{}{}^{+0.05}`$ keV. A Spearman-rank correlation also shows that the line energy increases with object luminosity at $`>`$99.9% significance (also see figure 7), again providing strong confirmation of the Nandra et al. (1997b) result whereby the iron lines in high luminosity sources tend be weaker and shifted bluewards of 6.4 keV.
### 6.3 The Compton Reflection Hump
Having ascertained that there is substantial iron line emission from several of the quasars in this sample, we next determine the properties of the Compton reflection hump that would be expected to accompany the line, on the assumption that the line results from reflection off optically thick matter. Constraints have therefore been placed on the amount of Compton reflection that occurs in the spectra of quasars in this sample. This is particularly important because if the iron lines that are described above result from the accretion disk, then we would expect to see evidence for a reflection ‘hump’ in the spectra of quasars. To test this, a Compton reflection model has been used (the pexrav code in XSPEC, Magdziarz & Zdziarski 1995) with solar abundances, a disk inclination angle of 30 degrees has been assumed (unless otherwise stated) and an exponential cut-off at $`>100`$ keV. The strength of the reflection component (R) is measured in terms of the solid angle $`\mathrm{\Omega }`$ subtended by the primary X-ray source to the disk; then R is then given by R $`=\mathrm{\Omega }/2\pi `$. Although R is left as a free parameter in the spectral fitting, normally for a geometrically thin disk a value of R near to 1 would be expected. The model also used assumes that the ‘reflecting’ material (in this case presumably the accretion disk) is of a neutral or low ionisation state.
The Compton reflection component was systematically fitted to the brightest quasars in the sample (both at low and high z), i.e. where the signal-to-noise ratio allows a sufficient constraint of these additional model parameters. The results are now presented here and are shown in table 6 for those quasars in which the reflection component has been constrained; the table either shows the best fitting value of the relative strength, R, of the reflection component or the statistical upper limit to R.
It can be seen that in the low redshift quasars it has generally not been possible to constrain the amount of reflection present in the X-ray spectra. The only low redshift quasars in which the amount of reflection has been constrained are in the radio-loud quasars 3C 273 and 4C 74.26. The 68% upper limit on the strength of the reflection component in 3C 273 is R $`<0.25`$ (or R $`<0.33`$ at 90% confidence). This result is perhaps not surprising as 3C 273 is a core dominated radio-loud quasar, the reflection component should be considerably diluted by the beaming effects of a relativistic jet. An unusually strong reflection component is found in the lobe-dominated radio-quasar 4C 74.26 (also see Brinkmann et al. 1998); however the presence of a strong soft excess (with temperature KT$`300`$ eV for a blackbody) could confuse the detection of this reflection component, given the limited ASCA bandpass. Constraints are not possible on the amount of reflection in the other low z objects.
However, unlike for low z quasars, in the high redshift quasars there is a greater band-pass in hard X-rays, as this ($`>10`$ keV) part of the spectrum is redshifted into the ASCA rest-frame. Therefore in the high redshift quasars, which have sufficient signal to noise, it should be possible to constrain the higher energy Compton reflection component. Firstly for the high redshift radio-loud quasars, it can be seen from the table that the reflection component is very weak or consistent with no reflection, with the relative strength of the component constrained to R $`<<1`$. In fact the presence of a reflection component with R=1 (i.e. as would be expected from scattering of X-rays off an accretion disk) is excluded in all the distant radio-loud quasars considered here, at $`>`$99.99% confidence. Examples are S5 0014+813, NRAO 140, S5 0836+715, Q 1508+571, PKS 2126-158 and PKS 2149-306. However the lack of a reflection component in these bright radio-loud quasars is perhaps not surprising as many of these objects are believed to be jet-dominated AGN, where the Doppler boosted jet component can dominate over the reflection component from near to the central engine; this is also consistent with the weak iron lines that are observed in this class of object.
It has also been possible to constrain the amount of reflection in some of the high redshift (z $`>1`$) radio-quiet quasars. In the luminous quasars PG 1634+706, PG 1718+481 and HE 1104-1805, the reflection component is constrained to R $`<<1`$. The upper limits (90% confidence) on these 3 quasars are R $`<0.38`$, R $`<0.25`$ and R $`<0.15`$ respectively; in addition a reflection component in the high z QSO HS 1946+765 is also constrained to R $`<0.2`$. The presence of a reflection component with R=1 (as expected from an accretion disk), is excluded at $`>`$99.9% confidence in PG 1634+706, PG 1718+418 and HE 1104-1805. All 3 of these quasars are particularly luminous (with L$`{}_{210}{}^{}10^{46}`$ erg/s and M<sub>V</sub> = -28 to -29, using q<sub>0</sub>=0.5); both PG 1634 and HE 1104 are radio-quiet (R$`{}_{L}{}^{}=0.587`$ and R$`{}_{L}{}^{}=0.94`$ respectively), and PG 1718 is radio-intermediate (R$`{}_{L}{}^{}=1.58`$). Compton reflection has not been constrained in the other distant radio-quiet quasars (lack of signal-to-noise). However the evidence suggests that the amount of neutral reflection in highly luminous high redshift quasars is much weaker than what is observed in lower luminosity Seyfert 1 galaxies (e.g. Pounds et al. 1990 Nandra & Pounds 1994). As this trend is observed in the radio-quiet as well as the radio-loud quasars, the lack of reflection features in quasars cannot just be interpreted in terms of increased Doppler boosting from a relativistic jet. This, along with the above trends in the iron line emission, will be discussed in the next section.
### 6.4 The Nature of Iron lines and Disk Reflection in Quasars
The analysis performed in the above section show that properties of iron line emission and the associated reflection hump in quasars seem less straightforward than in the Seyfert 1s. As was explained earlier, the common picture in the Seyfert 1s is for the hard X-rays to scatter off the inner accretion disk, producing a broad iron K$`\alpha `$ line, iron K edge and a Compton reflection hump. The fact that the line is highly broadened and fits the relativistic motions expected in the inner disk around a black hole (e.g. Tanaka et al. 1995), strongly supports the interpretation that these ‘reflection’ features originate from the accretion disk. In the Seyfert 1s, this reprocessing is almost completely from material that is neutral or of a low ionisation state (i.e. Fe i to Fe xvi).
However this is apparently not the case in quasars. It was found that of the 21 quasars in this sample with significant iron K line detections, a large proportion of these (at least 11) have iron line energies ($`>6.6`$ keV) that are consistent with originating from highly ionised material (e.g. Fe xxiv to Fe xxvi). In the high luminosity end of the quasar distribution (both radio-loud and radio-quiet), the actual strength of the iron K line component decreases considerably (i.e. from EQW $`=100150`$ eV to EQW $`<50`$ eV), providing strong confirmation of the results in Nandra et al. (1997b). Additionally we also find that the strength of the neutral reflection hump is diminished considerably in the high luminosity quasars.
So how can one explain the observations in the quasars, which seem to contrast with the situation in the Seyfert 1 galaxies? Firstly the energy of the line emission in these quasars implies that the ionisation state of the reprocessing material (the surface layers of the putative accretion disk) is clearly much higher than in the Seyfert 1s. As has been pointed out from photoionsation modeling (e.g. see Ross & Fabian 1993; Matt, Fabian & Ross 1993; Ross, Fabian & Young 1999), the surface layers of such disks can become substantially photoionised as the fractional accretion rate ($`\dot{m}`$) of the central engine increases. Perhaps at the Eddington limit very little if any iron K emission is seen, as the disk becomes fully ionised down to several Thomson depths.
This picture is consistent with the most luminous quasars (either radio-loud or radio-quiet), with very weak iron lines, accreting near the Eddington limit. Additionally, in the radio-loud objects, the relativistic jet can further weaken the line and reflection component. Furthermore this hypothesis can also explain the apparent lack of a reflection hump from neutral material in the high z quasars, found in section 6.3. Indeed it has been postulated (Ross, Fabian & Young 1999), that the disk becomes more reflective in high luminosity quasars (as the photoelectric opacity below Fe K decreases at higher ionisation) which can reproduce the apparently featureless spectra of some high z quasars. The lack of contrast between the continuum and the reprocessed X-rays can then explain the apparent absence of a reflection ‘hump’. Further support for the presence of a highly ionised accretion disk in one object comes from an observation ( with ASCA and RXTE) of the luminous quasar PDS 456. Although the line emission is fairly weak, a deep, highly ionised edge (at 8.7 keV) is present in the spectrum of this quasar. This iron K edge could originate from a high ionisation reflector. The detailed spectral fitting of this quasar, from simultaneous ASCA and RXTE data, is discussed in a separate paper (Reeves et al. 2000).
## 7 X-ray Absorption in Quasars
### 7.1 General Properties
The intrinsic neutral column densities for the 62 quasars analyzed are shown in table 2, fitted in the rest frame of the quasar. An additional Galactic absorption component was also fitted in the observers rest frame, as described previously. Unless an intrinsic absorption column is detected at 90% confidence (in addition to the Galactic absorption), an upper limit only has been quoted. Intrinsic absorption (at $`>`$90% significance) is present in 20 radio-loud quasars (RLQs) and 10 radio-quiet quasars (RQQs), i.e. 30 out of the 62 quasars.
The strongest absorption is seen in the most distant radio-loud quasars. Out of the 20 RLQs with significant intrinsic absorption, 9 of the 13 RLQs at z $`>2`$ show evidence for absorption with $`N_\mathrm{H}`$ in the order of 10$`{}_{}{}^{22}10^{23}\mathrm{cm}^2`$. In particular, strong intrinsic absorption is seen in the distant radio-loud quasars S5 0014+813, PKS 0528+134, S5 0836+715, PKS 2126-158 and PKS 2149-306 all at $`>`$99.9% significance and columns of about $`10^{22}\mathrm{cm}^2`$ or greater. The largest absorption column is seen in S5 0014+813, at red-shift z = 3.4, with a column density of ($`6.5\pm 2.6`$)$`\times 10^{22}\mathrm{cm}^2`$ in the rest frame of the quasar (although note that the Galactic absorption towards S5 0014+813 is rather high).
Considering radio-loud quasars at intermediate redshifts ($`1<\mathrm{z}<2`$), 6 out of nine quasars also have evidence for absorption, but with lower columns than at high z, typically $`10^{21}10^{22}\mathrm{cm}^2`$. However low z radio-loud quasars at redshifts z$`<`$1 generally only show weak absorption ($`N_\mathrm{H}`$$`<10^{21}\mathrm{cm}^2`$). Significant absorption is only seen in 5 out of 13 low z radio-loud quasars, these are 3C 109.0, 3C 273, 3C 279, PG 1425+267 and 4C 74.26. 3C 109.0 and 4C 74.26 are more lobe dominated than some of the other RLQs, which may imply some dependence on orientation; for these two objects $`N_\mathrm{H}`$ $`10^{21}\mathrm{cm}^2`$. The columns towards 3C 273 and 3C 279 are both low with $`N_\mathrm{H}`$$`<10^{21}\mathrm{cm}^2`$.
In comparison, significant absorption is also seen in only 10 radio-quiet quasars out of 27. At high redshifts (z $`>2`$) there are only 2 radio-quiet quasars with statistically significant intrinsic absorption; HE 1104-285 (at z=2.319) has a column of $`1.81\times 10^{22}\mathrm{cm}^2`$ at $`>`$99% significance, whilst RDJ 13434+0001 (at z=2.350) has an intrinsic column of $`4.8\times 10^{22}\mathrm{cm}^2`$ although this is less well constrained. HE 1104-285 is unusual in that it is associated with a gravitational lens system (Wisotski et al. 1993), and RDJ13434+0001 may be a rare example of a type II QSO (Almaini et al. 1995). In the case of the low z RQQs, virtually all the objects are consistent with columns of $`N_\mathrm{H}`$$`<10^{21}\mathrm{cm}^2`$, i.e. similar to Galactic sized columns.
### 7.2 The Warm Absorber
However there are some low z quasars which show evidence for absorption from warm (or partially ionised) matter. An example of this is the radio-quiet quasar PG 1114+445, which has a warm absorber with a column of typically $`N_\mathrm{H}`$$`10^{21}`$ to $`10^{22}\mathrm{cm}^2`$, with the main contribution from the ionised Ovii and Oviii edges. Similar results for PG 1114 were reported by George et al. (1997). Another low z warm absorber quasar similar to PG 1114+445 is MR 2251-178 (Reeves et al. 1997, Pan, Stewart & Pounds 1990). We also detect a significant ionised absorption component in the quasar IRAS 13349+243, this has been reported previously by Brandt et al. (1997b) who suggest that a dusty warm absorber is likely.
There are only 2 other quasars, not reported in the literature to date, with possible evidence for a warm absorber in this sample. One of these is the radio-loud quasar PG 1425+267 which has a warm absorber of column a few $`\times 10^{21}\mathrm{cm}^2`$, although the detection is marginal statistically (95% confidence) and the ionisation ($`\xi 10\mathrm{erg}\mathrm{cm}\mathrm{s}^1`$) is rather poorly constrained. However perhaps the most interesting case is the luminous (M$`{}_{V}{}^{}27`$) nearby radio-quiet quasar PDS 456. The complex residuals to a power-law fit for this quasar are shown in figure 8. An iron K edge is present at 8.7$`\pm `$0.2 keV, and the optical depth of the edge is very large ($`\tau =1.6\pm 0.5`$). The edge is detected at $`>`$99.99% confidence (or $`\mathrm{\Delta }\chi ^2=56`$ for 2 interesting parameters). In addition we can also fit a low energy warm absorber component to PDS 456, with a high column of material ($`5\times 10^{22}\mathrm{cm}^2`$) and ionisation parameter ($`500\mathrm{erg}\mathrm{cm}\mathrm{s}^1`$). The spectrum of PDS 456 is discussed in detail in a separate paper (Reeves et al. 2000). The main interpretation of the fitting is that there is a warm absorber component at lower energies, but that a high ionisation disk reflector could be responsible for the deep iron K edge. There is no other evidence for the presence of such a strong iron K edge in any of the other quasars, although the constraints in many cases are quite poor.
On the whole in this sample there are only a few examples of warm absorbers in quasars. This is also supported by ROSAT PSPC observations of PG quasars (Laor et al. 1997), in which only about 5% have X-ray spectral features due to a warm absorber. George et al. (2000) also found relatively few warm absorbers in an ASCA sample of PG quasars. In contrast the occurrence of warm absorber features is common in the Seyfert 1s (e.g. Reynolds 1997, George et al. 1998). This may be due to the absorbing material being at a different level of ionisation in the quasars. For instance if such matter was very highly ionised, the material would be essentially transparent at Ovii and Oviii energies. Another possibility is that there is a lack of such material in quasars. However a more straightforward explanation may simply be due to quasar redshift, where the Ovii and Oviii features are red-shifted out of the ASCA bandpass.
### 7.3 On the correlation between X-ray absorption and redshift
The observations that have been discussed in section 7.1 suggest that the amount of neutral X-ray absorption seems to be higher in the high z objects than in the low z quasars. We have therefore performed Spearman-rank correlations in order to investigate any apparent correlation between neutral $`N_\mathrm{H}`$ and z. As in previous sections, the use of survival statistics has been employed in order to allow upper-limits to be used in the correlations. (Note that all quasars with flux levels $`<1\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup> were excluded from these correlations, as useful constraints cannot be placed on the amount of absorption present.) A simple correlation between $`N_\mathrm{H}`$ (fitted in the quasar frame - corrected from Galactic absorption) and z is significant at $`>`$99.99% confidence; thus the amount of absorption intrinsic to the quasar is seen to increase with redshift. This apparent trend has been reported previously, with smaller samples of objects (see Elvis et al. 1994, Cappi et al. 1997, Reeves et al. 1997, Fiore et al. 1998). Typically columns in high redshift (z$`>`$2) quasars are of the order $`10^{22}`$cm<sup>-2</sup> (QSO frame), whereas at low redshift N<sub>H</sub> is only of the order few $`\times `$ 10<sup>20</sup>cm<sup>-2</sup>. The positive correlation between $`N_\mathrm{H}`$ and z is illustrated in figure 9.
It as has been reported (e.g. Cappi et al. 1997, Mukai 1998) the average neutral column in the SIS0 and SIS1 detectors tends to be overestimated by $`23\times 10^{20}\mathrm{cm}^2`$, due to uncertainties in the low energy calibration of these instruments. However it is also worth noting that a simultaneous Beppo-Sax and ASCA observation of 3C 273 showed that the SAX LECS and MECS instruments were in agreement with the ASCA SIS (Orr et al. 1998). Nevertheless, in order to take into account any possible calibration effect, an additional column of $`N_\mathrm{H}`$$`=3\times 10^{20}\mathrm{cm}^2`$ has been added to the spectral fit, in the local z=0 ASCA rest-frame, and the intrinsic columns have subsequently been refitted. Note that all SIS data below 0.6 keV have also been ignored, where the calibration problems tend to be worse. Despite of this, the correlation between $`N_\mathrm{H}`$ and z was still found to be significant at $`>`$99.9% confidence. Excluding all those quasars that lie at low Galactic latitudes ($`\beta 20^{}`$), also has little effect on the significance level of the correlations.
By fitting the absorption in the quasar rest-frame it is possible that the amount of $`N_\mathrm{H}`$ could be overestimated, as the absorbing material could lie anywhere along the line-of-sight to the quasars. Indeed if there was just some local effect that could produce some spurious $`N_\mathrm{H}`$ measurement (such as SIS calibration or uncertainty in the Galactic column) then this would be magnified towards the quasars at higher z. Therefore the above correlations have also been performed by fitting the absorption column in the local z=0 frame (after correcting for the Galactic column). The correlation between $`N_\mathrm{H}`$ and z is still found to be significant at $`99.9`$% confidence even in this case. Again the calibration uncertainties of ASCA have been taken into account by adding an additional column of $`3\times 10^{20}\mathrm{cm}^2`$ and by removing all quasars at low $`\beta `$. However even after these checks were performed, the correlation remained significant. Thus the statistical evidence, even after accounting for calibration and uncertainties in the Galactic absorption, seems to suggest that the correlation between $`N_\mathrm{H}`$ and z is real.
Therefore the most striking observation here is the discovery of moderately large absorption columns in several of the high z quasars. As these high z quasars are predominantly type I radio-loud AGN, obscuration in terms of the molecular torus (as observed in Seyfert 2s) seems unlikely. The physical origins of this absorption can then either be local to the rest frame of the quasars at high z, or it may be associated with matter at intermediate redshifts (i.e. from intervening line of sight material) not physically connected with the quasars. It is possible that the if the absorbing material is intrinsic to these quasars, then it could be similar in origin to the high ionisation absorbers observed in more nearby AGN. The most substantial source of soft X-ray absorption from intervening matter would probably be from damped Lyman-$`\alpha `$ absorption systems, however the number density of these dense systems is reported to be quite low (O’Flaherty & Jakobsen 1997). A more detailed account of the possible origins of this X-ray absorption will not be discussed further in this paper, but can be found in the literature (e.g. Elvis et al. 1998, Cappi et al. 1997, Reeves et al. 1997, Elvis et al. 1994). However observations using the superior low energy throughput provided by XMM and Chandra are needed to confirm this trend, and to provide further clues as to the possible causes.
## 8 Summary & Conclusions
We now summarize the findings of this paper. In particular, comparison is drawn to our earlier paper (Reeves et al. 1997 or R97), which contained a smaller sample of quasars (24 objects compared with the current sample size of 62).
Firstly we confirm the following main results from the R97 paper:-
* A decrease in the X-ray photon index, with radio-loudness, for all quasars.
* A strong confirmation of the correlation between the neutral X-ray absorption column ($`N_\mathrm{H}`$) and quasar redshift (z), in the sense that intrinsic $`N_\mathrm{H}`$ increases with z. Furthermore in this sample, the correlation does not depend on any calibration effects nor the rest-frame of the absorbing column.
* A decrease of iron K line equivalent width with increasing radio-loudness. The interpretation is that the strength of the reflection disk component (and therefore any contribution from the iron K line) is diminished, due to Doppler boosting of the X-ray continuum by the relativistic jet, in the core-dominated radio-loud quasars.
However we have seen several new effects and correlations in this paper, that were not reported in our previous R97 sample:-
* In the R97 paper, a correlation was found whereby the photon index for the radio-quiet objects increased with X-ray luminosity. No such correlation was found in this paper, over a large range of both redshift and luminosity, for the radio-quiet sources. The difference may be due to the increase in sample size (from 9 radio-quiet quasars in R97 to 27 in the present paper).
* Two correlations were found involving the iron K emission line. Firstly the strength of the iron K emission was observed to decrease with luminosity (i.e. an ‘X-ray Baldwin’ effect), regardless of whether the objects are radio-loud or radio-quiet. In addition the energy (or ionisation) of the iron line was found to increase with luminosity. Both of these trends confirm the result found in the Nandra et al. (1997b) paper, whereby the composite line profiles for AGN tend to be weaker at higher luminosities as well as narrower and shifted bluewards of 6.4 keV. The ‘X-ray Baldwin effect’ was first proposed for AGN on the basis of Ginga data by Iwasawa & Taniguichi (1993).
* A new effect is found whereby the strength of the Compton reflection ‘hump’ is weaker in the most luminous quasars at high redshifts (z$`>`$1). The effect is observed not only in the jet-dominated radio-loud sources, but also in the radio-quiet quasars. This finding is consistent with the ‘X-ray Baldwin effect’ discussed above, and suggests that as a whole the neutral disk reflection component in the high luminosity quasars is generally weaker than in the lower luminosity sources such as the Seyfert 1s.
* A trend has also been found for the radio-quiet quasars in this sample, whereby the X-ray (2-10 keV) photon index increases with decreasing optical H$`\beta `$ width. Thus the quasars with the steepest X-ray spectra tend to have the narrowest H$`\beta `$ FWHM. This trend has previously been found in the lower luminosity Seyfert 1s (e.g. Brandt et al. 1997), but has not been reported before for the more luminous quasars.
* Soft X-ray excesses are also found a significant proportion (9) of the low z quasars in this sample. Interestingly the majority of the quasars with strong soft excesses are those with the narrowest optical H$`\beta `$ widths (where H$`\beta `$ FWHM $`<2000`$ km/s).
* The temperatures of the soft X-ray excesses in this paper vary in the range between kT = 100 - 300 eV, for simple blackbody fits. In a majority of the cases the temperatures are probably too hot to result by direct thermal emission from the putative quasar accretion disk. Instead, one possibility is that the soft excess originates via thermal Comptonisation of UV photons from the disk in a hot corona. Another possibility is that the ‘soft excess’ results from reprocessing. In particular emission and/or reflection from the surface of a highly ionised inner accretion disk could reproduce the observed excess in soft X-ray flux.
* A systematic search has been carried out for the presence of warm or ionised absorbers in this quasar sample. Only a smaller number were found. We confirm the presence of a warm absorber in three previously reported cases (PG 1114+445, IRAS 13349+243 and MR 2251-178). The only new warm absorbers reported here are in the radio-quiet quasar PDS 456 and a marginal detection in the radio-loud quasar PG 1425+267. Overall the apparent rarity of warm absorbers in quasars may be due to different (higher) ionisation, a smaller covering fraction, or is perhaps just due to the redshift effect.
So how can we place all these observations facts into a general scheme for quasars. Firstly the differences between radio-loud and radio-quiet quasars seem relatively straightforward. In the radio-loud quasars, a strong Doppler boosted emission component from the relativistic jet can account for the higher luminosities, the generally flatter X-ray spectra as well as the diminished iron K line and reflection component in these objects. One question of real interest for future study is whether the central engine is the same in the radio-loud quasars as it is in the radio-quiet quasars. For instance the structure of the accretion disk may be different in the radio-loud quasars; sensitive studies of the reflection component and iron line in the RLQs (i.e. with XMM) may help to determine this.
The properties of the radio-quiet quasars on the whole seem more complex. As has been seen in this paper, there is little or no dependence on the X-ray continua of quasars on luminosity and therefore presumably the black hole mass. However perhaps the one driving factor responsible for the individual properties of quasars may be the fractional accretion rate $`\dot{m}`$ of the central engine (i.e. the ratio of mass accretion to the Eddington rate - or the Eddington ratio). A high fractional accretion rate can result in the surface layers of the disk becoming highly photoionised, which subsequently can have several effects on the X-ray spectra. Depending on the degree of ionisation, ionised rather than neutral iron K emission lines can dominate the disk reflection spectrum, as observed. Furthermore at even higher ionisations, the neutral reflection component (and iron line emission) can appear to be weaker, particularly if the disk is fully ionised to several Thomson depths (e.g. Nayakshin et al. 1999, Ross et al. 1999), also in agreement with the apparent properties of the more luminous radio-quiet quasars. A further effect is that, for high Eddington ratios, stronger soft X-ray emission can be produced. This may partly arise as the intrinsic thermal emission from the disk can become stronger (Ross, Fabian & Mineshinge 1992). Another possibility is that as the disk is more highly ionised, it can become more reflective at soft X-ray energies, producing a steepening of the X-ray spectrum at low energies (i.e. as a result of the ionised disk reflection component). Thus a high accretion rate (relative to Eddington) may explain the strong soft excesses in some of the objects considered earlier.
It is also interesting to return to the question of the dichotomy between the broad and narrow line quasars that was considered earlier. It has been postulated in the literature that the narrow optical H$`\beta `$ lines may be an indicator of a high accretion rate (Pounds et al. 1995, Laor et al. 1997). If this is correct this could indeed account for the differences between the 2 types of objects. A high accretion rate may explain the strong soft excesses observed both in the narrow-line quasars (6 out of 8 objects in this sample) and also the lower luminosity narrow-line Seyfert 1s (NLS1s) in other samples (Vaughan et al. 1999, Leighly 1999). As explained this can be caused by increased intrinsic disk emission or an increase in the reflectiveness of the disk in the soft X-ray band. Also if the intrinsic disk emission is stronger, this can also account for the steeper 2-10 X-ray slopes, via increased Compton cooling (Pounds et al. 1995). Additionally some of the narrow-line quasars also show evidence for ionised iron K emission, also indicative of a high ionisation disk and thus a high accretion rate, although the evidence is tentative so far (also see Vaughan et al. 1999). So the narrow-line quasars (as well as the NLS1s) may radiate at a relatively high fraction of the Eddington rate, whereas in general the broad-lined quasars may be sub-Eddington, perhaps similar to the normal Seyfert 1s, but with more massive central black holes. An important question to ask is whether there are any narrow-line quasars at higher redshifts (z$`>1`$).
Finally the amount of soft X-ray absorption towards quasars was found to increase with redshift (also see Fiore et al. 1998 for a similar analysis of ROSAT quasars). This correlation is apparently robust, even when calibration effects and uncertainties in the amount of local absorption are taken into account. The main question that has arisen from this, is whether this absorption is intrinsic to the quasars or whether it originates from line-of-sight matter. Given the low-number density of high column systems (such as damped Ly-$`\alpha `$ systems) that could cause appreciable X-ray absorption (O’Flaherty & Jakobsen 1997), the most likely scenario is that the bulk of this absorbing material is local to the quasars or the host galaxy environment. As there seems to be a comparative lack of absorption in some radio-quiet quasars (see Fiore et al. 1998), this excess $`N_\mathrm{H}`$ may be associated with radio-loud quasars (also see Sambruna et al. 1999). However further data (with XMM and Chandra) is required to determine the exact location and cause of this absorption.
## Acknowledgments
We thank the ASCA support teams, at GSFC and ISAS, for their help. In particular, thanks to Ken Pounds and Simon Vaughan for proof reading the paper and providing useful discussions. We also thank the anonymous referee for providing suggestions to improve the paper. This research made use of data obtained from the Leicester Database and Archive Service (LEDAS) at the Department of Physics and Astronomy, Leicester University, UK, and the High Energy Astrophysics Science Archive Research Center (HEASARC), provided by NASA’s Goddard Space Flight Center.
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# Creation of Skyrmions in a Spinor Bose-Einstein Condensate
## Abstract
We propose a scheme for the creation of skyrmions (coreless vortices) in a Bose-Einstein condensate with hyperfine spin $`F=1`$. In this scheme, four traveling–wave laser beams, with Gaussian or Laguerre-Gaussian transverse profiles, induce Raman transitions with an anomalous dependence on the laser polarization, thereby generating the optical potential required for producing skyrmions.
03.75.Fi, 32.80.-t
The recent experimental success in all–optical trapping of an atomic Bose-Einstein condensate (BEC) opens the prospect of studies into the internal structure of spinor BECs, including the possibility of creating vortex states without core, or skyrmions, in the BEC. Skyrmions, which do not have an ordinary vortex core due to the spin degree of freedom, offer a myriad of new physical phenomena beyond those presented by other vortex states. One feature is the reduction of kinetic energy associated with the rotation by transferring this energy to the spin. Another interesting property of skyrmions is that they do not represent a topological excitation, although their flux is vortex-like far away from the line of symmetry . In this paper, we propose an optical method to create skyrmions in a condensate of alkali atoms. This method, which employs laser beams and Raman transition to generate skyrmions in the BEC, is related to proposals for the creation of vortices in a single–component BEC.
For a BEC with hyperfine spin $`F=1`$, a skyrmion wave function, which is axially symmetric around the $`z`$-axis, has the form
$$\left(\begin{array}{c}\psi _1\\ \psi _0\\ \psi _1\end{array}\right)=\sqrt{\rho }\left(\begin{array}{c}\mathrm{cos}^2\left(\frac{\beta }{2}\right)\\ \sqrt{2}e^{i\varphi }\mathrm{cos}\left(\frac{\beta }{2}\right)\mathrm{sin}\left(\frac{\beta }{2}\right)\\ e^{2i\varphi }\mathrm{sin}^2\left(\frac{\beta }{2}\right)\end{array}\right),$$
(1)
with $`\rho (z,r_{})`$ the total density of atoms, $`𝐫_{}=(x,y)`$ the transverse coordinate vector and $`r_{}=|𝐫_{}|`$ the radial distance from the $`z`$-axis. The angle $`\varphi `$ corresponds to the orientation of $`𝐫_{}`$ in the $`x`$-$`y`$–plane. The function $`\beta (r_{})`$, which characterizes the spin state of the condensed atoms, is related to the superfluid velocity $`𝐯_s`$ of the system by
$$𝐯_s=\frac{\mathrm{}}{Mr_{}}[1\mathrm{cos}(\beta (r_{}))]𝐞_\varphi ,$$
(2)
where $`𝐞_\varphi (\mathrm{cos}\varphi ,\mathrm{sin}\varphi ,0)`$ is the unit vector in the $`\varphi `$ direction. Because a skyrmion has no vortex core, $`\beta (0)=0`$ must hold in order to avoid a singularity. For $`\beta (r_{})=\pi /2`$, the superfluid velocity reduces to that of an ordinary vortex state.
To create a skyrmion we consider a BEC which first is magnetically trapped in the $`m=1`$ hyperfine state. We assume that the trap is then switched off and the optical potential is applied to the BEC. Thus, our initial state is given by $`\psi _1=\sqrt{\rho (r_{},z)}`$ and $`\psi _0=\psi _1=0`$. Our objective is to design an optical potential which transfers this state into the state (1).
Generally a coherent optical potential for atoms is created by applying several highly detuned laser beams which induce Raman transitions between different angular momentum states or internal states of the atomic BEC (cf Fig. 1). In order to preserve the cylindrical symmetry of the skyrmions along the $`z`$-axis, the laser beams must also propagate along this axis. This restriction raises the following challenge: the polarization orientation of laser beams propagating along the $`z`$-axis is generally presumed to be in a superposition of $`𝐞_x`$ and $`𝐞_y`$. However, as a Raman transition of such beams can only effect transfers of the type $`\psi _1`$ to $`\psi _1`$, the possibility of creating skyrmions of the form (1) seems to be excluded.
Fortunately the flexibility in the transverse structure of laser beams permits a portion of the laser mode to be polarized along the axis of propagation. Mathematically, this contribution to longitudinal polarization has its origin in the vanishing divergence of the electric field for any laser beam; that is,
$$\text{div}𝐄(𝐫)=0.$$
(3)
For a laser beam propagating along the $`z`$-axis, with transverse mode function $`u(𝐫_{})`$ and polarization vector $`𝐞_x`$, the divergence condition (3) makes it necessary that the complete positive-frequency part of the electric field is given by
$$𝐄^{(+)}(𝐫)=e^{\pm ikz}e^{i\omega t}\left[i𝐞_xu(𝐫_{})𝐞_z\frac{1}{k}u^{}(𝐫_{})\right],$$
(4)
where $`u^{}(𝐫_{})`$ denotes the derivative of $`u`$ with respect to $`x`$. The second term on the right-hand side of Eq. (4) is essential for inducing transitions to $`\psi _0`$ component.
Our proposal includes four traveling–wave laser beams and a homogenous magnetic field $`𝐁=B𝐞_z`$. The purpose of the magnetic field is to lift the degeneracy of the ground states $`|m_g`$ with $`m_g=0,\pm 1`$ being the magnetic quantum number of the states. We write the corresponding Hamiltonian in the form
$$H_B=E_B\underset{m_g=1}{\overset{1}{}}m_g|m_gm_g|.$$
(5)
The laser beams induce transitions between the $`F=1`$ manifold of the atomic ground state and an $`F=1`$ manifold of excited atomic states, such as the $`D_1`$ line between <sup>2</sup>S$`{}_{1/2}{}^{}(F=1)`$ and <sup>2</sup>P$`{}_{1/2}{}^{}(F=1)`$ in Rubidium (<sup>87</sup>Rb). All four beams are linearly polarized in the $`x`$-direction and have a large detuning to suppress spontaneous emission. Each beam has a different mode structure: we include a Gaussian beam with frequency $`\omega _0`$ and a Laguerre-Gaussian beam of second order with frequency $`\omega _2`$. Both beams are copropagating along the positive $`z`$-axis. Our scheme further includes two copropagating Laguerre-Gaussian beams (along the negative $`z`$-axis) of first order with opposite orbital angular momentum and frequencies $`\omega _1`$ and $`\omega _1`$, respectively. All beams are assumed to have the same width $`w`$ (where $`w`$ is the width parameter for the Gaussian envelope of each beam). The electric field of each beam then is of the form (4), whereby the transverse mode functions are approximately given by
$$u_l=A_l\frac{r_{}^l}{w^l}e^{il\varphi }e^{r_{}^2/w^2},l=1,0,1,2,$$
(6)
The interaction of the atoms with the laser beams in the electric dipole approximation is given by
$$H_{\text{int}}=\mathrm{}\underset{l=1}{\overset{2}{}}\widehat{\mathrm{\Omega }}_le^{i\omega _lt}+\text{H.c.}.$$
(7)
where $`\widehat{\mathrm{\Omega }}_l`$ is the Rabi-frequency operator for the $`l^{\mathrm{th}}`$ laser beam, defined by
$`\mathrm{}\widehat{\mathrm{\Omega }}_le^{i\omega _lt}`$ $``$ $`𝐝^{}𝐄_l^{(+)}`$ (8)
$`=`$ $`{\displaystyle \underset{m_e,m_g}{}}m_e|𝐝^{}𝐄_l^{(+)}|m_g|m_em_g|`$ (9)
where $`m_e`$ denotes the magnetic quantum numbers for the excited states. Using standard techniques for Clebsch-Gordon coefficients (eg, see Ref. ), we find for each laser beam the expression
$`\mathrm{}\widehat{\mathrm{\Omega }}_l`$ $`=`$ $`𝒟e^{\pm ikz}[{\displaystyle \frac{iu_l}{2}}(|0_e1_g|+|1_e0_g|`$ (12)
$`+|1_e0_g|+|0_e1_g|)`$
$`{\displaystyle \frac{u_l^{}}{k\sqrt{2}}}(|1_e1_g||1_e1_g|)]`$
with $`𝒟`$ is a reduced matrix element. For simplicity we assume $`𝒟`$ to be real. The first contribution describes the usual transitions which are induced by $`x`$-polarized laser beams. The term proportional to $`u_l^{}`$ has its origin from the small $`z`$-polarized term which results from the vanishing divergence of the electric field.
The optical potential which can be deduced from such a configuration for largely detuned laser beams has the general form (see, e.g., Ref. )
$$V_{\text{opt}}=\mathrm{}\underset{l,l^{}=1}{\overset{2}{}}e^{i(\omega _l^{}\omega _l)t}\frac{\widehat{\mathrm{\Omega }}_l^{}^{}\widehat{\mathrm{\Omega }}_l}{\omega _l\omega _{\text{res}}},$$
(13)
where $`\omega _{\text{res}}`$ is the resonance frequency for the electronic transition. The total Hamiltonian for this system then takes the form
$$H=d^3x\underset{m_g,m_g^{}}{}\mathrm{\Psi }_{m_g}^{}m_g\left|H_{\text{1-p}}\right|m_g^{}\mathrm{\Psi }_{m_g^{}}+H_{\text{nl}}$$
(14)
with the one-particle contribution
$$H_{\text{1-p}}=\frac{𝐩^2}{2M}+H_B+V_{\text{opt}}$$
(15)
and the nonlinear interaction energy
$`H_{\text{nl}}`$ $`=`$ $`{\displaystyle \frac{\lambda _s}{2}}{\displaystyle \underset{m_g,m_g^{}}{}}{\displaystyle d^3x\mathrm{\Psi }_{m_g}^{}\mathrm{\Psi }_{m_g^{}}^{}\mathrm{\Psi }_{m_g^{}}\mathrm{\Psi }_{m_g}}`$ (20)
$`+{\displaystyle \frac{\lambda _a}{2}}{\displaystyle }[(\mathrm{\Psi }_1^{})^2(\mathrm{\Psi }_1)^2+(\mathrm{\Psi }_1^{})^2(\mathrm{\Psi }_1)^2`$
$`+2\mathrm{\Psi }_1^{}\mathrm{\Psi }_0^{}\mathrm{\Psi }_1\mathrm{\Psi }_0+2\mathrm{\Psi }_1^{}\mathrm{\Psi }_0^{}\mathrm{\Psi }_1\mathrm{\Psi }_0`$
$`2\mathrm{\Psi }_1^{}\mathrm{\Psi }_1^{}\mathrm{\Psi }_1\mathrm{\Psi }_1+2(\mathrm{\Psi }_0^{})^2\mathrm{\Psi }_1\mathrm{\Psi }_1`$
$`+2\mathrm{\Psi }_1^{}\mathrm{\Psi }_1^{}(\mathrm{\Psi }_0)^2]d^3x`$
where $`\lambda _i4\pi \mathrm{}^2a_i/M`$ is proportional to the corresponding scattering length $`a_i`$. Transforming to the interaction picture, with respect to $`H_B`$ of Eq. (5), the optical potential assumes the form
$`V_{\text{opt}}`$ $`=`$ $`{\displaystyle \frac{𝒟^2}{\mathrm{}}}{\displaystyle \underset{l,l^{}=1}{\overset{2}{}}}{\displaystyle \frac{e^{i(\omega _l^{}\omega _l)t}}{\omega _l\omega _{\text{res}}}}\{(\mathrm{𝟏}|00|){\displaystyle \frac{u_l^{}^{}u_l^{}}{2k^2}}`$ (25)
$`+(\mathrm{𝟏}+|00|+e^{2iE_bt}|11|`$
$`+e^{2iE_bt}|11|){\displaystyle \frac{u_l^{}^{}u_l}{4}}`$
$`\pm i(e^{iE_bt}|01|e^{iE_bt}|01|){\displaystyle \frac{u_l^{}^{}u_l^{}}{2\sqrt{2}k}}`$
$`i(e^{iE_bt}|10|e^{iE_bt}|10|){\displaystyle \frac{u_l^{}^{}u_l}{2\sqrt{2}k}}\}`$
It is now our task to derive from this general expression a special potential that is suitable for the creation of skyrmions. As our initial state is the ground-state for a BEC trapped in $`m_g=1`$, it becomes obvious from Eq. (1) that we require matrix elements of the form $`|11|e^{2i\varphi }`$ , $`|01|e^{i\varphi }`$ and so on. As the transverse mode function $`u_l`$ is proportional to $`e^{il\varphi }`$, its derivative $`u_l^{}`$ is approximately proportional to $`e^{i(l1)\varphi }`$. This relation permits the following construction.
The detunings $`\mathrm{\Delta }_l\omega _l\omega _{\text{res}}`$ of the laser beams are the most significant frequencies in the problem. To construct the desired matrix elements we assume that $`\mathrm{\Delta }_2\mathrm{\Delta }_0`$ and $`\mathrm{\Delta }_1\mathrm{\Delta }_1`$. However, $`\mathrm{\Delta }_0`$ and $`\mathrm{\Delta }_1`$ should be sufficiently different so that a Raman transition involving the beams 0 and 1 is suppressed. Assuming that the second largest frequency is given by the magnetic interaction energy $`E_B`$ allows the rotating-wave approximation to be done with respect to this frequency in Eq. (25). By assuming that
$`\delta \omega _{02}`$ $``$ $`\omega _0\omega _2+2E_BE_B,\omega _0\omega _2,`$ (26)
$`\delta \omega _{11}`$ $``$ $`\omega _1\omega _1E_BE_B,\omega _1\omega _1,`$ (27)
many terms can be neglected, and we arrive at
$`V_{\text{opt}}`$ $``$ $`{\displaystyle \frac{𝒟^2}{2\mathrm{}k^2}}(\mathrm{𝟏}|00|)\left\{{\displaystyle \frac{|u_0^{}|^2+|u_2^{}|^2}{\mathrm{\Delta }_0}}+{\displaystyle \frac{|u_1^{}|^2+|u_1^{}|^2}{\mathrm{\Delta }_1}}\right\}`$ (32)
$`+{\displaystyle \frac{𝒟^2}{4\mathrm{}}}(\mathrm{𝟏}+|00|)\left\{{\displaystyle \frac{|u_0|^2+|u_2|^2}{\mathrm{\Delta }_0}}+{\displaystyle \frac{|u_1|^2+|u_1|^2}{\mathrm{\Delta }_1}}\right\}`$
$`+{\displaystyle \frac{𝒟^2}{4\mathrm{}\mathrm{\Delta }_0}}\{|11|u_0^{}u_2e^{i\delta \omega _{02}t}+\text{ H.c}\}`$
$`{\displaystyle \frac{i𝒟^2}{2\sqrt{2}\mathrm{}k\mathrm{\Delta }_1}}\{|01|u_1^{}u_1^{}e^{i\delta \omega _{11}t}\text{ H.c}\}`$
$`+{\displaystyle \frac{i𝒟^2}{2\sqrt{2}\mathrm{}k\mathrm{\Delta }_1}}\{|01|u_1^{}u_1^{}e^{i\delta \omega _{11}t}\text{ H.c}\}.`$
This term can be further simplified because the derivative of $`u_l`$ typically scales as $`1/w`$ compared to $`u_l`$ itself. Including the prefactor $`1/k`$ that appears together with each derivative, we see that every term with a derivative is suppressed by a factor of
$$\epsilon \frac{1}{kw}=\frac{\lambda }{2\pi w}.$$
(33)
As the width of a laser beam always exceeds its wavelength $`\lambda `$, the factor (33) is always much smaller than one. For instance, for an optical transition ($`\lambda =795`$ nm) a quite strongly focused beam with a width of 5 $`\mu `$m has $`\epsilon 0.025`$. This allows us to develop $`V_{\text{opt}}`$ in a power series in $`\epsilon `$. Before doing so we observe that in Eq. (32) the transitions between $`|0`$ and $`|\pm 1`$ are suppressed by a factor of $`\epsilon `$ as compared to the transition between $`|1`$ and $`|1`$. Since this is undesirable for the production of skyrmions we assume that the intensity $`A_0`$ (and hence $`u_0`$) of the Gaussian laser beam defined by Eq. (6) will be of the order of $`\epsilon `$ as compared to the other prefactors $`A_l`$. Inserting the mode functions (6) into Eq. (32) and developing $`V_{\text{opt}}`$ to first order in $`\epsilon `$, we find
$`V_{\text{opt}}`$ $``$ $`{\displaystyle \frac{𝒟^2e^{2\overline{r}^2}}{4\mathrm{}}}\{(\mathrm{𝟏}+|00|)({\displaystyle \frac{A_2^2}{\mathrm{\Delta }_0}}\overline{r}^4+{\displaystyle \frac{2A_1^2}{\mathrm{\Delta }_1}}\overline{r}^2)`$ (37)
$`+\left(|11|e^{2i\varphi }\overline{r}^2{\displaystyle \frac{A_0A_2}{\mathrm{\Delta }_0}}e^{i\delta \omega _{02}t}+\text{ H.c}\right)`$
$`+{\displaystyle \frac{i\sqrt{2}A_1^2}{\mathrm{\Delta }_1}}\epsilon \overline{r}[|01|e^{i\varphi }(1\overline{r}^2(1+e^{2i\varphi }))e^{i\delta \omega _{11}t}`$
$`+|01|e^{i\varphi }(1\overline{r}^2(1+e^{2i\varphi }))e^{i\delta \omega _{11}t}\text{ H.c}]\}`$
Eq. (37) is the main result of this paper. It consists of several parts which will now be analyzed.
The first row in Eq. (37) represents an optical trapping potential if the laser detuning $`\mathrm{\Delta }_l`$ is positive. For $`\overline{r}r_{}/w1`$ it is harmonic. Its physical origin is a Raman transition which returns to the same internal hyperfine level. As the state $`|m_g=0`$ has two possibilities ($`|m_e=\pm 1`$) to make such a transition the trapping potential is twice as strong as compared two the states $`|m_g=\pm 1`$. Note also that this term is of order $`\epsilon ^0`$ and therefore represents the dominant contribution.
The second row describes ordinary Raman transitions between the hyperfine states $`|m_g=1`$ and $`|m_g=1`$. As laser 2 corresponds to a Laguerre-Gaussian beam of second order, an orbital angular momentum of $`2\mathrm{}`$ will be transferred to the atoms, as it is required for the formation of the skyrmion (1). We remark that the same result could have been achieved by replacing the Gaussian beam 0 by a weak Laguerre-Gaussian beam of order 1, and beam 2 by a Laguerre-Gaussian beam of order $`1`$. The optical potential would only be changed in that the trapping potential in the first row of Eq. (37) would become harmonic (apart from the overall exponential factor). From an experimental point of view this alternative may be advantageous because it only requires two different types of laser beams (Laguerre-Gaussian of order 1 and -1) instead of four types.
The third and the fourth row in Eq. (37) describe anomalous Raman transitions in which the atoms absorb or emit a photon with polarization $`𝐞_z`$ in one of the two processes involved. In Fig. 1 this corresponds to a vertical transition from $`|m_g=1`$ (or $`|m_g=1`$) to $`|m_e=1`$ (or $`|m_e=1`$). The possibility of these transitions is a direct consequence of the full mode structure (4). Since these vertical transitions are proportional to the derivative $`u_l^{}`$ of the transverse mode functions, no orbital angular momentum is transferred even though the corresponding beams are of Laguerre-Gaussian type $`\pm 1`$ (more precisely, for $`\overline{r}1`$ the terms proportional to $`\mathrm{exp}[\pm 2i\varphi ]`$ in Eq. (37) are negligible). The second process which completes the Raman transition is an ordinary transition from $`|m_g`$ to $`|m_e=m_g\pm 1`$. Since the corresponding laser beam is of Laguerre-Gaussian type $`\pm 1`$ the total orbital angular momentum transferred to the atoms is $`\pm \mathrm{}`$. As a result, the anomalous Raman transitions create simple vortex states in $`|m_g=0`$ and doubly excited vortices in $`|m_g=1`$.
It might be counterintuitive that the $`z`$-polarized contribution to the electric field (4), which usually can safely be neglected, does now play an important role for the creation of skyrmions. In fact, its contribution remains as small as usual compared to the dominant parts (diagonal elements of $`V_{\text{opt}}`$). The important point in our proposal is that we can make all matrix elements inducing internal transitions equally large by using a weak laser beam 0 so that the usual Raman transition $`m_g=1m_g=1`$ is not stronger than the anomalous Raman transitions.
As the diagonal elements of $`V_{\text{opt}}`$ are much larger than the off-diagonal elements they provide a strong trapping potential. According to Eq. (1) a skyrmion is a superposition of the ground-state in $`m_g=1`$ and a singly (doubly) excited vortex state in $`m_g=0(m_g=1)`$, respectively. Because of the strong trapping potential these states all have different energies. Thus, to induce transitions using the off-diagonal elements of $`V_{\text{opt}}`$ these matrix elements need to rotate at some frequency to match a resonance condition. This can be done by adjusting the frequencies $`\delta \omega _{02}`$ and $`\delta \omega _{11}`$ appropriately. The efficiency of the skyrmion production strongly depends on $`\delta \omega _{02}`$ and $`\delta \omega _{11}`$.
We have numerically examined the time evolution of a spin-1 BEC in the potential (37). Our simulations are based on the split-step method in two spatial dimensions for a three-component nonlinear Schrödinger equation (It is assumed that the BEC is homogeneously extended along the $`z`$-axis with a length of 50 $`\mu `$m). We consider a BEC of 15000 Rubidium atoms ($`M=1.45\times 10^{25}`$ kg, $`a_s=5.4`$ nm, $`a_a=0.05`$ nm ). For a much larger number of atoms our method seems to be unsuitable because the potential (37) has a maximum height due to the exponential envelope. A strong repulsive interaction between the atoms then pushes the atoms out of the trapping area.
For the optical potential (37) we consider the situation that $`A_2^2/\mathrm{\Delta }_0=2A_1^2/\mathrm{\Delta }_1`$ holds and the amplitude of the Gaussian beam is given by $`A_0=3\epsilon A_2/\sqrt{2}`$. The effective Rabi frequency of the trapping potential is then given by $`𝒟^2A_2^2/(4\mathrm{}^2\mathrm{\Delta }_0)`$ which we take to be 66 kHz (For a detuning $`\mathrm{\Delta }_0`$ of 1 GHz this corresponds to a laser intensity of about 25 mW). The width of the laser beams is 5 $`\mu `$m. For the results shown in Fig. 2 we have used $`\delta \omega _{02}=2.4`$ kHz and $`\delta \omega _{11}=2.95`$ kHz. At time $`t=6.6`$ ms this results in a superposition of internal states with about 27% of the atoms in state $`m_g=1`$, about 34% in the state $`m_g=0`$, and 39% in state $`m_g=1`$. For various other choices of $`\delta \omega _{02}`$ and $`\delta \omega _{11}`$ we have found as little as 4% population in the initial state $`m_g=1`$ and as much as 70% in one of the other two states. The phase of the three components agrees very well with the situation described by the state (1). Animations of the time evolution are available on the internet .
Acknowledgement: K.-P. M. would like to thank the Optik Zentrum Konstanz for financial support. This project has been supported by an Australian Research Council Small Grant.
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# Ultrahigh sensitivity of slow–light gyroscope
## Acknowledgements
We are grateful to Malcolm Dunn and Stig Stenholm for valuable discussions. U.L. gratefully acknowledges the support of the Alexander von Humboldt Foundation and of the Göran Gustafsson Stiftelse.
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# THE HAWKING-UNRUH TEMPERATURE AND DAMPING IN A LINEAR FOCUSING CHANNEL
## I Introduction
Many of the effects of quantum fluctuations on the behavior of charged particles can be summarized concisely by an effective temperature first introduced in gravitational fields by Hawking , and applied to accelerated particles (with the neglect of gravity) by Unruh .
Hawking argued that the effect of the strong gravitational field of a black hole on the quantum fluctuations of the surrounding space is to cause the black hole to radiate with a temperature
$$T=\frac{\mathrm{}g}{2\pi ck},$$
(1)
where $`g`$ is the acceleration due to gravity at the surface of the black hole, $`c`$ is the speed of light, and $`k`$ is Boltzmann’s constant. Shortly thereafter, Unruh argued that an accelerated observer should become excited by quantum fluctuations to a temperature
$$T=\frac{\mathrm{}a^{}}{2\pi ck},$$
(2)
where $`a^{}`$ is the acceleration of the observer in its instantaneous rest frame.
In a series of papers, Bell and co-workers , have noted that electron storage rings provide a demonstration of the utility of the Hawking-Unruh temperature (2), with emphasis on the question of the incomplete polarization of the electrons due to quantum fluctuations of synchrotron radiation. The author has commented on how the Hawking-Unruh temperature can be used to characterize quickly the limits on damping of the phase volume of beams in electron storage rings , leading to well-known results of Sands .
## II Quantum Analysis of a Linear Focusing Channel
Recently, Chen, Huang and Ruth have discussed radiation damping in a linear focusing channel , finding that in such devices the beam can be damped to the quantum mechanical limit set by the uncertainty principle. I show here how this result follows very quickly from an application of the Hawking-Unruh temperature.
A linear focusing channel is a beam-transport system that confines the motion of a charged particle along a straight central ray via a potential that is quadratic in the transverse spatial coordinates. This potential can be characterized by a spring constant $`k`$, and hence the frequency $`\omega `$ of transverse oscillations (as observed in the laboratory frame) of a particle of mass $`m`$ and Lorentz factor $`\gamma `$ is
$$\omega =\sqrt{\frac{k}{\gamma m}}.$$
(3)
If the amplitude of the oscillation in transverse coordinate $`x`$ is called $`x_0`$, then the amplitude $`a_0`$ of the corresponding transverse acceleration is
$$a_0=x_0\omega ^2=\frac{kx_0}{\gamma m}.$$
(4)
To apply the Hawking-Unruh temperature, we consider the motion in the instantaneous rest frame of the particle. Supposing the transverse oscillations are small, the instantaneous rest frame is very nearly the frame in which the particle has no longitudinal motion. Quantities measured in the instantaneous rest frame will by denoted with the superscript $``$. Thus, in the instantaneous rest frame the amplitude of the transverse acceleration as measured is
$$a_0^{}=\gamma ^2a_0=\frac{\gamma kx_0}{m},$$
(5)
the frequency of the oscillation is
$$\omega ^{}=\gamma \omega ,$$
(6)
and hence the transverse spring constant of the focusing channel appears as
$$k^{}=m\omega ^2=\gamma k.$$
(7)
In the instantaneous rest frame, the charge particle finds itself in a bath of radiation of characteristic temperature given by eq. (2) with acceleration $`a^{}`$ given by eq. (5). This bath can be regarded as the effect of quantum fluctuations, which excite transverse oscillations (having two degrees of freedom) to characteristic energy $`U^{}`$ (as measured in the instantaneous rest frame) given by
$$U^{}=kT=\frac{\mathrm{}a_0^{}}{2\pi c}=\frac{\mathrm{}\gamma kx_0}{2\pi mc}.$$
(8)
The energy of transverse oscillation can also be written in terms of the (invariant) transverse amplitude $`x_0`$ as
$$U^{}=\frac{k^{}x_0^2}{2}=\frac{\gamma kx_0^2}{2}.$$
(9)
Hence, the amplitude of excitation of the transverse oscillations is
$$x_0=\frac{\mathrm{}}{\pi mc}=\frac{\mathrm{\lambda ̄}_C}{\pi },$$
(10)
where $`\mathrm{\lambda ̄}_C`$ is the (reduced) Compton wavelength of the particle.
The amplitude (10) must, however, be compared to the amplitude of the zero-point oscillations of the system, considered as a quantum oscillator:
$$x_{0,\mathrm{zero}\mathrm{point}}=\sqrt{\frac{\mathrm{}}{\gamma m\omega }}=\sqrt{\frac{\mathrm{\lambda ̄}_C\mathrm{\lambda ̄}}{\gamma }},$$
(11)
where $`\mathrm{\lambda ̄}=c/\omega `$ is the laboratory (reduced) wavelength of the transverse oscillation as measured along the beam axis. In practical laboratory devices, we will have $`\mathrm{\lambda ̄}\gamma \mathrm{\lambda ̄}_C`$. Hence, the excitation of the transverse oscillations by fluctuations in the radiation of the oscillating charge, as are described by the Hawking-Unruh temperature, is negligible compared to the zero-point fluctuations of the transverse oscillations. In this sense, we can say along with Huang, Chen and Ruth that the radiation does not excite the transverse oscillations, and those oscillations will damp to the quantum-mechanical limit.
In futuristic devices, for which $`\gamma >\mathrm{\lambda ̄}/\mathrm{\lambda ̄}_C`$, i.e., when
$$\gamma >\frac{mc^2}{k\mathrm{\lambda ̄}_C},$$
(12)
quantum excitations of oscillations in a linear focusing channel would become important. When (12) holds, the transverse oscillations would be relativistic even when their amplitude is only a Compton wavelength. The strength of the transverse fields in the channel would then exceed the QED critical field strength (in the average rest frame),
$$E_{\mathrm{crit}}=\frac{m^2c^3}{e\mathrm{}}=1.6\times 10^{16}\text{V/cm}=3.3\times 10^{13}\text{Gauss},$$
(13)
and the beam energy would be rapidly dissipated by pair creation.
Another way of viewing a practical linear focusing channel is that its Hawking-Unruh excitation energy, (8), is small compared to the zero-point energy, $`\mathrm{}\omega ^{}/2=\gamma \mathrm{}\omega /2`$ of transverse oscillations.
The quantum-mechanical limit for transverse motion can, of course, also be deduced from the uncertainty principle:
σxσpx
>
,
>
subscript𝜎𝑥subscript𝜎subscript𝑝𝑥Planck-constant-over-2-pi\sigma_{x}\sigma_{p_{x}}\mathrel{\vbox{\kern 0.0pt\hbox{$>$} \kern 0.0pt\hbox{$\sim$} }}\hbar, (14)
which leads to a minimum normalized emittance of
$$ϵ_N=\frac{\sigma _x\sigma _{p_x}}{mc}\mathrm{\lambda ̄}_C,$$
(15)
corresponding to geometric emittance of
$$ϵ_x=\frac{ϵ_N}{\gamma \beta _z}\frac{\mathrm{\lambda ̄}_C}{\gamma }.$$
(16)
## III Semiclassical Analysis
In a quantum analysis of a linear focusing channel, we found that the transverse oscillations can damp to the limit set by the uncertainty principle. Hence, in a classical analysis we would expect the damping to be able to proceed until the transverse amplitude was zero.
Indeed, a simple analysis confirms this. Transform to the longitudinal rest frame, in which the particle’s motion is purely transverse. The particle has nonzero kinetic energy in this frame, but its average momentum is zero. The radiation due to the transverse oscillation is reflection symmetric about the transverse plane in this frame, so the radiation carries away energy but not momentum. With time, all of the energy would be radiated away, and the particle would come to rest. The transverse oscillations will have damped to zero without affecting the longitudinal motion.
If we add the concept of photons to the preceding analysis, we can say that the radiated photons carry away momentum along the direction of emission, but the radiation pattern is symmetric, so the averaged radiated momentum is zero. Again, the radiation carries away energy, now in the form of photons.
Back in the lab frame, we view the photons as carrying away a small amount of longitudinal momentum on average, as a result of the Lorentz transformation of the energy radiated in the longitudinal rest frame. This momentum, however, is only that part of the particle’s longitudinal momentum associated with its transverse oscillation; the longitudinal velocity of the particle is unaffected.
On average, the photons carry away no transverse momentum in the lab frame, and the average momentum of the radiated photons is therefore parallel to the beam axis in lab frame. However, there is no need to argue that the momentum of individual radiated photons is parallel to the beam axis, nor to imply that the matter of the focusing channel absorbs transverse momentum in a manner than affects the kinematics of the radiation process .
## IV Comparison of a Linear Focusing Channel to a Wiggler
A comparison with the behavior of particle beams in a wiggler is instructive. Here the transverse confinement of the beam motion is provided by a series of alternating transverse magnetic fields. This has the notable effect that even if a particle enters the wiggle parallel to the beam axis, transverse oscillations will result whose amplitude is independent of the initial transverse coordinate.
In contrast, a particle that enters a linear focusing channel parallel to and along the axis undergoes no oscillation, no matter what is the particle’s longitudinal momentum.
We thereby see that radiation damping cannot reduce the oscillations in a wiggler to zero unless the longitudinal momentum falls to zero also, since the wiggler continually re-excites transverse oscillations for any particle with nonzero kinetic energy.
Another difference between a wiggler and a linear focusing channel can be seen by going to the longitudinal rest frame. In the case of the wiggler, the alternating magnetic fields in the laboratory transform to fields that are very much like a plane wave propagating against the direction of the laboratory motion of the beam. The radiation induced by this effective plane wave is not symmetric with respect to the transverse plane, but results in a net kick of the particle into the backward direction.
Viewed in the lab frame, we find that along with the damping of their transverse oscillations, the particles’ longitudinal momenta are significantly reduced. To maintain the initial longitudinal momentum, the beam must be reaccelerated. The momentum (and energy) added back into the beam then increases the amplitude of the transverse oscillations, and the damping cannot continue beyond some limit.
In contrast, in a linear focusing channel, the transverse damping proceeds without significant reduction in the longitudinal momentum of the particle, and the transverse oscillations can damp to the quantum limit without the need of adding energy back into the beam.
## Acknowledgments
I wish to thank Pisin Chen, Ron Ruth and Max Zolotorev for conversations on radiation damping. This work was supported in part by DoE grant DE-FG02-91ER40671.
## References
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# Optical studies of the X-ray transient XTE J2123-058 -I. Photometry
## 1 Introduction
Soft X-ray transients (SXTs) are a subclass of low-mass X-ray binaries (LMXBs) that are characterized by episodic X-ray outbursts (usually lasting for several months), when the X-ray luminosities can increase by as much as a factor of 10<sup>7</sup> (van Paradijs & McClintock, 1995). The observed optical flux is generated by X-ray re-processing in the accretion disc and the companion star. These outbursts recur on a time scale of decades, but in the interim the SXTs are in a state of quiescence and the optical emission is dominated by the radiation of the faint companion star. This offers the best opportunity to analyze the properties of this star and obtain dynamical information which eventually enables us to constrain the nature of the compact object. There are currently 12 SXTs with identified optical counterparts, with 8 dynamical black-holes and 3 confirmed neutron stars: CenX-4, Aql X-1 and 1608-522 (van Paradijs & McClintock, 1995). In addition, there are a few neutron star binaries exhibiting X–ray on and off states ( EXO0748–676, 4U 2129+47 and SAX J1748.9–2021 ), although they are not classified as SXTs because they do not show the classic fast rise and slow exponential/linear decay.
The X-ray transient J2123–058 was discovered on 29 June 1998 by the Rossi X-Ray Time Explorer (RXTE) (Levine et al. 1989) reaching a peak X-ray flux of 100 mCrabs (2-12 keV). Its high Galactic latitude (b=-36.2) is unusual among transients, an indication that J2123–058 might be a member of the galactic halo population. The optical counterpart was promptly identified with a variable star of R=17.2 (Tomsick et al. 1998a), which was only marginally visible on a digitized U.K. Schmidt plate, suggesting a preoutburst magnitude R$``$20 (Zurita et al. 1998). Spectra obtained early in the outburst showed strong high-excitation lines of He ii $`\lambda `$4686, C iii/N iii $`\lambda `$4640 and weak Balmer emission embedded in broad absorptions (Tomsick et al. 1998a, Hynes et al. 1999). These features are frequently observed in SXTs during outburst (e.g. Callanan et al 1995) and persistent LMXBs (e.g. Augusteijn et al. 1998). Type-I (thermonuclear) bursts have been detected both in X-rays (Takeshima and Strohmayer 1998) and optical (Tomsick et al. 1998b), a signature of a neutron star in the binary. The outburst light curve exhibited regular 0.7 mag deep triangular-shaped minima repeating every 6hrs (Casares et al. 1998, Tomsick et al. 1998b), a strong indication of high inclination. This provided the first evidence for the system orbital period (P= 5.957$`\pm `$0.003) which was later confirmed by a radial velocity study of the Heii $`\lambda `$4686 emission line (Hynes et al. 1998). In addition, Ilovaisky and Chevalier (1998) reported the presence of a 0.3 mag modulation with a period of 7.2 days, probably caused by the precessing disc. Since 26 Aug the system had settled down to its quiescent state at R$``$21.7 (Zurita & Casares 1998).
This paper presents the results of a comprehensive set of observations that have led to detailed optical light curves from outburst through the decay into quiescence. Our spectroscopy will be the subject of a second paper (Hynes et al. 2000).
| Date | HJD<sup>(∗)</sup> | Exp/Filter | Telescope |
| --- | --- | --- | --- |
| Jul02 98 | -3 | 1xR | IAC80 0.8m<sup>(1)</sup> |
| Jul04 ” | -1 | 1xB,1xV,1xR | IAC80 0.8m |
| Jul06 ” | 01 | 80xR | OGS 1m<sup>(2)</sup> |
| Jul07 ” | 02 | 79xR | OGS 1m |
| Jul08 ” | 03 | 20xR | OGS 1m |
| Jul09 ” | 04 | 33xR | OGS 1m |
| Jul10 ” | 05 | 20xR | OGS 1m |
| Jul12 ” | 07 | 42xR | OGS 1m |
| Jul13 ” | 08 | 24xR | OGS 1m |
| Jul14 ” | 09 | 2xV | M. Canopus 1m<sup>(3)</sup> |
| Jul16 ” | 11 | 2xV | M. Canopus 1m |
| Jul18 ” | 13 | 1xB,1xV,1xR | M. Canopus 1m |
| Jul19 ” | 14 | 40xBV | Crimean 0.5m<sup>(4)</sup> |
| | | 36xR | Crimean 1.25 m<sup>(5)</sup> |
| | | 2xB,2xV,2xR | M. Canopus 1m |
| Jul20 ” | 15 | 42xBV | Crimean 0.5m |
| | | 74xR | Crimean 1.25m |
| | | 1xB,5xV,1xR | M. Canopus 1m |
| Jul21 ” | 16 | 36xBV | Crimean 0.5m |
| | | 94xR | Crimean 1.25m |
| Jul22 ” | 17 | 25xBV | Crimean 0.5m |
| | | 43xR | Crimean 1.25m |
| Jul23 ” | 18 | 48xBV | Crimean 0.5m |
| | | 45xR | Crimean 1.25m |
| Jul24 ” | 19 | 61xR | Kryonerion 1.2m<sup>(6)</sup> |
| | | 55xBV | Crimean 0.5m |
| Jul26 ” | 21 | 70xR | OGS 1m |
| | | 46xBV | Crimean 0.5m |
| Jul27 ” | 22 | 77xR | OGS 1m |
| | | 33xBV | Crimean 0.5m |
| Jul28 ” | 23 | 77xR | OGS 1m |
| | | 40xBV | Crimean 0.5m |
| Jul29 ” | 24 | 74xR | OGS 1m |
| | | 49xBV | Crimean 0.5m |
| Jul30 ” | 25 | 81xR | OGS 1m |
| | | 33xBV | Crimean 0.5m |
| Jul31 ” | 26 | 40xBV | Crimean 0.5m |
| Date | HJD<sup>(∗)</sup> | Exp/Filter | Telescope |
| --- | --- | --- | --- |
| Aug01 98 | 27 | 23xBV | Crimean 0.5m |
| | | 4xR | Crimean 0.38m |
| Aug02 ” | 28 | 1xR | IAC80 0.8m |
| Aug03 ” | 29 | 1xR | IAC80 0.8m |
| Aug04 ” | 30 | 2xR | Crimean 0.38m |
| | | 1xR | IAC80 0.8m |
| Aug05 ” | 31 | 2xR | Crimean 0.38m |
| Aug12 ” | 38 | 1xR | IAC80 0.8m |
| Aug15 ” | 41 | 12xR | IAC80 0.8m |
| Aug16 ” | 42 | 26xR | IAC80 0.8m |
| Aug17 ” | 43 | 8xR | IAC80 0.8m |
| Aug19 ” | 45 | 3xV,1xR | M. Canopus 1m |
| Aug26 ” | 52 | 18xR | OGS 1m |
| | | 1xV | IAC80 0.8m |
| Aug27 ” | 53 | 14xR | OGS 1m |
| Sep02 ” | 59 | 4xR | IAC80 0.8m |
| Sep23 ” | 80 | 11xR | OGS 1m |
| Sep24 ” | 81 | 8xR | OGS 1m |
| Sep25 ” | 82 | 5xR | OGS 1m |
| Sep26 ” | 83 | 5xR | OGS 1m |
| Jun21 99 | 351 | 1xR,1xV,1xI | JKT 1m<sup>(7)</sup> |
HJD–2451000
<sup>1</sup>IAC80–80cm Telescope in the Observatorio del Teide (Tenerife).
<sup>2</sup> 1 m Optical Ground Station in the Observatorio del Teide (Tenerife).
<sup>3</sup> Mount Canopus 1m telescope at Tasmania.
<sup>4</sup> 0.5m telescope at Crimea
<sup>5</sup> 1.25m reflector of SAI Crimean Station.
<sup>6</sup> 1.2m telescope at the Kryonerion Astronomical station of the National Observatory of Athens at Kryonerio of Korinthia.
<sup>7</sup> 1m Jacobus Kapteyn Telescope in the Roque de los Muchachos Observatory (La Palma).
## 2 Observations and Data Reduction
We observed J2123-058 during the period 1998 July–September with the IAC80 and 1m Optical Ground Station (OGS) in the Observatorio del Teide (Tenerife), the 1.25m SAI Station and 0.5m and 0.38m telescopes at Crimea, the Mount Canopus 1m telescope at Tasmania and the 1.2m telescope at the Kryonerion Astronomical station of the National Observatory of Athens at Kryonerio of Korinthia. On the night of 21 of June of 1999, we obtained VRI images, using the 1m Jacobus Kapteyn Telescope in the Roque de los Muchachos Observatory (La Palma). The integration times ranged from 1 to 40 min, depending on the telescope, atmospheric conditions and the brightness of the target. The observing log is presented in Table 1. All images were corrected for bias and flat-fielded in the standard way.
We applied aperture photometry to our object and several nearby comparison stars within the field of view, using iraf. We selected four comparison stars which were checked for variability during each night, and over the entire data set. We calibrated the data was using 9 standard stars from several fields (Landolt 1992), from which we constructed a colour dependent calibration. The Crimean data in R were calibrated relative to the nearby comparison star from the USNO catalogue which was also in to the previous calibrated data. White light photometry was also obtained in Crimea using a 0.5m telescope equipped with a TV detector. The TV camera operates in the range $`\lambda \lambda `$3500–8000 with maximum sensitivity at $`\lambda `$4000 and effective wavelength at $`\lambda _{eff}`$4851. Magnitudes were measured relative to a nearby star from the USNO catalogue which was calibrated at B and V wavelengths. ’Equivalent BV magnitudes’, were calculated by interpolating the fluxes in B and V to find the flux at the effective wavelength of the TV detector. For further details see Pavlenko, Prokofieva & Dolgushin (1989).
## 3 Light Curves
### 3.1 Long term light curve
Figure 1 compares the overall light curve of J2123–058 in optical (VR bands) and X-rays (2–12 keV). The long-term behaviour in X-rays shows a classical FRED (fast-rise exponential-decay) morphology, with characteristic e-folding times of 2.4 d (rise) and 19 d (decay). These time scales coincide with the mean values of the distributions of SXTs (Chen, Shrader & Livio 1997). Twenty days after the peak of the outburst, the X-ray intensity reaches a secondary maximum (of less than half of the peak intensity). The secondary maximum is also suggested by our optical data.
We identify 3 different stages in the optical light curve: the outburst plateau (until $``$ 10 Aug), the decay phase ($``$ 10–26 Aug) and quiescence (from $``$ 26 Aug). In the plateau phase the object brightness decays at a moderately slow rate of $``$ 0.03 mag$``$day<sup>-1</sup> although a modulation is clearly visible in the nightly mean magnitudes. The time scale of this variability is consistent with the 7-d modulation attributed to disc precession by Ilovaisky & Chevalier (1998). From $``$ 10 Aug the optical light curve began an abrupt fall at a rate of $``$ 0.2 mag$``$day<sup>-1</sup> before reaching quiescence on 26 Aug at R=21.7 (see Figure 1).
Taking V(peak)$``$17.2 and V(quiescent)$``$22.9, we estimate a total outburst amplitude of 5.7 mags. Using the empirical relation $`\mathrm{\Delta }`$V=14.36-7.6$`\mathrm{lg}`$P<sub>orb</sub>(hr) (Shahbaz & Kuulkers 1998), we would expect a total amplitude $`\mathrm{\Delta }`$V$``$8.4, 2.7 mags larger that what is observed. This difference can be explained by 3 effects: (1) the outburst brightness of the disc being reduced by a factor cos(i), since the disc would be foreshorted by the high binary inclination angle, (2) we are also assuming that the quiescent flux is completely dominated by the companion star with no veiling from the accretion disc and (3) if the secondary star is sufficiently evolved for it to be degenerate it would appear intrinsically fainter. We believe the discrepancy is due to a combination of these effects and can only be resolved by obtaining optical spectroscopy in quiescence.
We see evidence for optical bursts both during outburst (29 July with an amplitude of 0.3 mag) and at the onset of quiescence (27 Aug and 2 Sept with amplitudes larger than 1 mag). The optical bursts observed during the ourburst decay is most probably caused by re-processing of X-ray bursts. However, the origin of the optical bursts observed during the onset of quiesence is more puzzling, given that the source had almost reached its quiescent (low luminosity) X-ray state. Note that no simultaneous X-ray observations exist, which might shed light on the origin of these optical flashes observed during quiescence.
Our V and ’BV’ (Figure 1), show that the BV magnitudes drop faster than R. Moreover the amplitudes in BV and R increase during the decay. Figure 2 presents colour information of J2123-058 as a function of time. Although we cannot directly compare the ’BV’–R colour with V–R it is clear from the plot that the system redden as the outburst decays and the secondary’s contribution increases. In quiescence, we obtained an upper limit to the quiescent V magnitude through 2x2400s images of J2123–058 with the IAC80 telescope. We also have a marginal detection in V on the night of 21 Jun 1999 using the JKT telescope. Our colours are consistent with the spectral type of a late–K main sequence star.
### 3.2 The orbital light curve
Representative light curves of J2123-058 in different stages of the outburst cycle are presented in Figure 3. The data have been folded on the updated ephemeris given by Zurita & Casares (1998): HJD 2451042.639(5) + 0.24821(3) E. Note the dramatic changes in amplitude and morphology of the light curves as the outburst decays.
The July light curves are flat topped with broad triangular minima. They have a full amplitude of 0.7 mag and are reminiscent of the 5.1 hr period eclipsing transient EXO 0748–676, although its long term behaviour is characterized by high and low X-ray/optical states rather than transient outbursts. We also find similarities with the eclipsing LMXB source 2A1822–371 whose complex optical light curve has been successfully modeled by X-ray heating and partial occultation of a thick non-axisymmetric accretion disc and a faint companion (Mason et al. 1980). Asymmetries in the eclipse minima, are also seen in J2123-058 on individual nights and also in the curves by Tomsick et al. 1998.
On the other hand, the light curve of Aug 16 is almost sinusoidal and shows a peak-to-peak amplitude of $``$1.4 mag. It resembles the outburst light curve of the 5.2 hr neutron star system 4U 2129+47, where the large amplitude has been attributed to X-ray reprocessed radiation from the heated face of the optical star (Thorstensen et al.1979). Note, however, that at the onset of the fast optical decay ($``$10 Aug) the X-ray flux had already dropped by a factor $``$ 10 with respect to the outburst peak. Note also the presence of two narrow 0.2 mag dips at phases 0 and 0.5 which suggest the presence of grazing eclipses.
Finally, the quiescent light curve displays a characteristic ellipsoidal modulation from the secondary star: a double-humped variation with a full amplitude of $``$0.4 mag. This light curve was produced by phase binning the entire quiescent data (from Aug 26 on) after detrending the night-to-night variability using a linear fit.
## 4 Distance estimate
In the context of King & Ritter’s (1998), the X-ray exponential decay seen in J2123-058 indicates that irradiation is strong enough to ionize the entire accretion disc. Also, a secondary maximum is expected one irradiated-state viscous time after the onset of the outburst and it can be used to calibrate the peak X-ray luminosity and hence the distance to the source $`D_{\mathrm{kpc}}`$ through
$$D_{\mathrm{kpc}}=4.3\times 10^5t_s^{3/2}\eta ^{1/2}f^{1/2}F_p^{1/2}\tau _d^{1/2}$$
where $`F_p`$ is the peak X-ray flux, $`t_s`$ the time of the secondary maximum after the peak of the outburst in days, $`\tau _d`$ the e-folding time of the decay in days, $`\eta `$ the radiation efficiency parameter and $`f`$ the ratio of the disc mass at the start of the outburst to the maximum possible mass (Shahbaz, Charles & King 1998). In our case, $`\tau _d`$=19 d, $`t_s`$20 d and $`F_p`$ can be estimated from the XTE count rate (6 counts s<sup>-1</sup> in the energy range 2-12 keV) which corresponds to 1.84 $`\times `$ 10<sup>-9</sup> erg cm<sup>-2</sup> s<sup>-1</sup>. Assuming $`\eta `$=0.15 and $`f`$=0.5 we find $`D_{\mathrm{kpc}}`$ = 5.7.
Alternatively, we can estimate the distance to the source by comparing the quiescent magnitude with the absolute magnitude of a main sequence star which fits within the Roche lobe of a 6hr period orbit. Combining Paczynski’s (1971) expression for the averaged radius of a Roche lobe with Kepler’s Third Law we get the well-known relationship between the secondary’s mean density and the orbital period: $`\rho =110/P_{\mathrm{hr}}^2`$ (g cm<sup>-3</sup>). Substituting the orbital period of J2123–058 we obtain $`\rho `$=3.1 g cm<sup>-3</sup> which corresponds to a K7V secondary star with mass $``$0.6 M and absolute magnitude $`M_R`$ 7. The dereddened quiescent magnitude is R=21.4 (using A<sub>V</sub>=0.37$`\pm `$0.15 as derived from the NaD1 line; see Hynes et al. 1998) which yields $`D_{\mathrm{kpc}}`$ = 7.7. Strictly speaking, this is a lower limit to the distance as we are neglecting any contribution by the accretion disc to the quiescent optical flux. However, note that the true distance is probably not too far off 8 kpc since our quiescence light curve does not show strong evidence for disc contamination (see section 5 and Figure 6). A conservative limit can therefore be provided by assuming a 50 percent contribution by the accretion disc to the continuum light. Allowing for 50 percent disc contamination (as observed in J0422+32, a black-hole transient with comparable orbital period; see Casares et al. 1995) we obtain R=22.2 for the companion star and hence $`D_{\mathrm{kpc}}`$ = 10.8. Hereafter we will adopt $`D_{\mathrm{kpc}}=8\pm 3`$ which is consistent, at the lower end, with other estimates based on photospheric expansion models of the X-ray bursts (Homan et al 1999, Tomsick et al 1999). A spectral type determination of the companion star is essential to refine this distance estimate.
## 5 Modelling
In an attempt to interpret the different light curves and derive the system parameters, we have used a model based on the work by de Jong et al. (1996, also see references included). The model assumes a flared accretion disc of the form $`hr^{9/7}`$ (with $`h`$ and $`r`$ the disc height and radius respectively) and a Roche lobe filling secondary and accounts for X-ray heating, shadowing effects and mutual eclipses of the disc and the secondary. The disc is assumed to radiate as a blackbody with a radial temperature distribution calculated according to Vrtilek et al. (1990). The intensity distribution on the secondary star is computed using Kurucz model atmospheres. The albedo of the accretion disc and the companion star are fixed to 0.95 and 0.40 respectively, following the results of de Jong et al (1996). The model parameters are the binary inclination (i), mass ratio (q=M<sub>1</sub>/M<sub>2</sub>), the accretion disc radius (R<sub>disc</sub>) defined as a fraction of the distance to the inner Lagrangian point (R$`_{\mathrm{L}_1}`$), the flaring angle of the accretion disc ($`\alpha `$) and the X-ray luminosity (L<sub>x</sub>).
In order to model the outburst light curve we have averaged our best quality light curves of the plateau phase (26-30 July) in 29 phase bins. Using L<sub>x</sub>=1.3$`\times `$ 10<sup>37</sup> erg s<sup>-1</sup> (for $`D_{\mathrm{kpc}}`$ = 8) and M<sub>2</sub>=0.6 M we performed a least-squares fit to the data. Our best fit solution gave a reduced chi squared $`\chi _\nu ^2`$ of 1.36 for $`i`$=76.0$`\pm `$1.0 degrees, R<sub>disc</sub>=0.75$`{}_{0.03}{}^{}{}_{}{}^{+0.06}`$ R$`_{\mathrm{L}_1}`$, $`\alpha `$=7.6$`{}_{0.2}{}^{}{}_{}{}^{+1.0}`$ degrees and $`q`$=4.6$`{}_{0.2}{}^{}{}_{}{}^{+0.5}`$. The uncertainties quoted are at the 99 percent confidence level and were obtained by grid searching the parameter of interest whilst optimizing the other model parameters. We have also rescaled the $`\chi _\nu ^2`$ values so that the minimum $`\chi _\nu ^2`$ is 1. The best model fit to the outburst data is shown in Figure 4 (solid line). In order to examine the effects of changing M<sub>2</sub>, we fitted the outburst data with M<sub>2</sub>=0.1M (Figure 4 – dotted line). We find that the derived parameters are the same, within the errors: $`\chi _\nu ^2=1.5`$, $`i`$=75.3$`{}_{1.2}{}^{}{}_{}{}^{+0.8}`$ degrees, R<sub>disc</sub>=0.83$`\pm `$0.05 R$`_{\mathrm{L}_1}`$, $`q`$=3.4$`\pm `$0.4, $`\alpha `$=8.8$`{}_{1.2}{}^{}{}_{}{}^{+0.8}`$ degrees. Regarding the decay stage, we have fitted the light curve of 16 Aug binned into 20 phase bins, when the X-ray luminosity was an order of magnitude lower, L<sub>x</sub>=1.3$`\times `$10<sup>36</sup> erg s<sup>-1</sup>. We fitted the decay stage data with this X-ray luminosity and M<sub>2</sub>=0.6M. Our best fit solution gave a reduced $`\chi _\nu `$ of 1.51 for $`i`$=72.0$`\pm `$3.0 degrees, R<sub>disc</sub>=0.56$`\pm `$0.06 R$`_{\mathrm{L}_1}`$, $`\alpha `$=5.7$`\pm `$0.5 degrees and $`q`$=4.2 ($`q`$ could not be constrained with this data set). The best model fit to the decay stage data is shown in Figure 5. Again 99 per cent confidence levels are quoted.
In Figure 6, we show the quiescent light curve which exhibits the characteristic ellipsoidal modulation of the secondary star: i.e. two equal maxima and two unequal minima. The minimun at phase 0.5 is expected to be deeper than the minimum at phase 0.0 because gravity darkening is more pronounced near the inner Lagrangian point $`L_1`$. This effect is important for high inclination systems (see eg. Avni and Bahcall, 1974). In the figure we also show model plots using i=73, q=4.6 (values which are consistent with those derived from the decay and outburst data), and no X–ray heating. The solid and dotted lines show plots with zero and 50 percent disc contamination. The former model probably best describes the data.
## 6 Discussion
XTE J2123–058 is a remarkable neutron star binary. Our optical light curve shows marked orbital modulations with dramatic variations as the outburst declines. These are very similar to the modulations observed in accretion disc corona (ADC) sources of comparable orbital periods (e.g. 4U 2129+47 and 2A 1822–371; see eg. McClintock et al. 1980 and Mason et al. 1982) although none of these are transient systems.
The X-ray light curve displays the classic SXT properties, namely a FRED morphology with typical e-folding rise and decay times and secondary maximum. Furthermore, the ratio of X-ray to optical luminosity \[$`\xi =B_0+2.5\mathrm{log}F_\mathrm{x}(\mu \mathrm{Jy})`$\] is also in excellent agreement with the observed distribution of LMXBs. Taking B=17.28 (Tomsick et al. 1998a) and F$`{}_{\mathrm{x}}{}^{}(212keV)100`$ mCrab (Levine, Swank & Smith 1998) at the outburst peak and assuming $`A_\mathrm{v}=0.37`$ (Hynes et al 1998) we obtain $`\xi `$=21.9, whereas the distribution peak of LMXBs gives $`\xi `$=21.8 $`\pm `$ 1.0 (see van Paradijs & McClintock 1995). This result implies that, despite its high inclination, the X-ray source in J2123-058 is not hidden by the accretion disc (i.e. $`\alpha `$90–i). This is consistent with the values our model favors for the binary inclination (i=73) and disc flaring angle ($`\alpha `$=5.7 – 7.6).
The longterm evolution of the optical light curve can be compared to those of other SXTs (e.g. GRO J0422+32, A0620-00, N. Muscae 91). They show a slow linear decay followed by a steeper fall. We find that J2123–058 also reproduces this behaviour, although the total amplitude and time scales are a factor of $``$ 2 shorter (see e.g. Callanan et al 1995).
We have modeled our R-band light curves of J2123–058 at different stages of the outburst including the obscuration effects and X-ray heating of the secondary star accretion disc. This led us to constrain the system inclination to i=73$`{}_{}{}^{}\pm 4`$. We find encouraging the excellent agreement between the inclination values obtained for the two independent light curves (July 26-30 and 16 Aug) during which $`L_\mathrm{x}`$ has dropped by one order of magnitude. The light curves at the plateau phase (July) are very similar to those of EXO 0748–676 and 2A 1822–371 with an extended depression of the luminosity from phase $``$0.7 until the eclipse and a steeper rise to maximum (see Mason et al. 1980 and Schmidtke et al 1987). Our model fit indicates that the accretion disc is the dominant source of light and the triangular shaped minima can be interpreted as eclipses of the accretion disc by the secondary star together with the changing aspect of the heated polar caps of the companion star. The dramatic changes observed in the light curves during decline, are triggered by large changes in the disc size and geometry. Our fit to the decay data (16 Aug) demands a thinner and smaller accretion disc which implies a smaller fraction of the disc is X-ray heated. Conversely, the secondary star is more exposed to the X-ray radiation and therefore the total amplitude of the modulation increases by a factor of 2 (to $`1.4`$ mag). The resulting light curve has a sinusoidal-like shape and is reminiscent of the LMXBs 4U 2129+47(=V1727 Cygni) and HZ Her. Our model fits implies a change of $``$30 percent in the disc size, as the system fades by 1.7 mags in the optical. The change in the radius of the disc size is what one expects. If angular momentum is transported outwards in the disc through viscous processes, then at outburst, since matter diffuses inwards, the angular momentum of that matter has to be transfered to the outer parts of the disc, and the radius of the disc is expected to expand. When the system is decaying, after the end of the mass transfer enhancement, the disc shrinks to its original radius (Livio & Verbunt 1988; Ichikawa & Osaki 1992). Observations of U Gem, Z Cha, OY Car, HT Car show that the accretion discs are indeed larger in outburst than in quiescence (Smak 1984b; O’Donohuge 1986; Harrop-Allin & Warner 1996). Comparing our results with disc radius variation in U Gem (Smak 1984a), we find approximately the same rate ofdecrease.
## 7 Acknowledgments
Part of this work is based on observations made with the European Space Agency OGS telescope operated on the island of Tenerife by the Instituto de Astrofisica de Canarias in the Spanish Observatorio del Teide of the Instituto de Astrofisica de Canarias.
We are grateful to M. Serra–Ricart, D.Alcalde, A.Gomez and P. Rodriguez–Gil for performing some of the observations. We are thank A. Dapergolas, E.T. Harlaftis and D. Galloway for making their data available to us. JC acknowledges support by the Spanish Ministry of Science grant 1995-1132-02-01.We thank J. van Paradijs and the Amsterdam group for the use of their X–ray binary model.
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# Supernovae versus Neutron Star Mergers as the Major 𝑟-Process Sources
## 1. Introduction
All of the actinides and approximately half of the stable nuclei with mass numbers $`A100`$ in the solar system were produced by rapid neutron capture, the $`r`$-process. While the astrophysical site of this process remains a mystery, the main candidate environments are neutrino-heated ejecta from core-collapse supernovae (SNe; Woosley & Hoffman 1992; Meyer et al. 1992; Takahashi, Witti, & Janka 1994; Woosley et al. 1994) and decompressed ejecta from neutron star mergers (NSMs; Lattimer & Schramm 1974, 1976; Symbalisty & Schramm 1982; Freiburghaus, Rosswog, & Thielemann 1999b). To be a major source for the $`r`$-process, an environment must satisfy two criterions: one on reproducing the solar $`r`$-process abundance pattern and the other on supplying the total amount of $`r`$-process material in the present Galaxy. Using more physical parametrizations than the previous approach based on constant neutron number density and temperature (e.g., Kratz et al. 1993), several recent studies (Hoffman, Woosley, & Qian 1997; Meyer & Brown 1997; Freiburghaus et al. 1999a) derived the $`r`$-process conditions that can produce relative abundance patterns with peaks at $`A130`$ and 195 as observed in the solar system. At present, it seems that conditions in models of neutrino-heated ejecta from SNe are deficient (e.g., Qian & Woosley 1996) while those in models of decompressed ejecta from NSMs are promising (e.g., Freiburghaus et al. 1999b) for an $`r`$-process. However, due to the uncertainties in the theoretical models (e.g., Qian & Woosley 1996; Freiburghaus et al. 1999b), a reliable comparison of the actual conditions in these two environments with the $`r`$-process conditions cannot be made yet to help establish or discriminate either environment as the $`r`$-process site. Furthermore, for both environments the models showed that if the ejecta were composed of $`r`$-process material, then the amount provided by a single event combined with the estimated number of SNe or NSMs over Galactic history would be adequate to account for the present Galactic $`r`$-process inventory (e.g., Mathews & Cowan 1990; Qian & Woosley 1996; Rosswog et al. 1999). Therefore, the above two criterions cannot readily identify SNe or NSMs as the major $`r`$-process sources.
In this Letter I discuss a phenomenological approach to test SNe and NSMs as the major $`r`$-process sources. By considering mixing of the ejecta from an individual event with the interstellar medium (ISM), I show that observations of metal-poor stars are difficult to explain if NSMs are the major $`r`$-process sources (§2). I further show that a self-consistent picture of $`r`$-process enrichment by SNe based on the same consideration is supported by meteoritic data on <sup>182</sup>Hf and <sup>129</sup>I in the early solar system and by observations of metal-poor stars (§3). To emphasize the importance of observations in establishing the major $`r`$-process sources, I conclude with a discussion of two possible direct tests of the SN $`r`$-process model: gamma rays from decay of $`r`$-process nuclei in SN remnants and surface contamination of the companion by SN $`r`$-process ejecta in binaries (§4).
## 2. $`r`$-Process Abundances and Neutron Star Mergers
I assume that as far as the $`r`$-process is concerned, the mergering of a neutron star (NS) with a black hole (BH) is similar to that of two neutron stars. These two kinds of events are collectively referred to as NSMs. The rate of such events in the Galaxy is quite uncertain. Here I take an average rate $`f_\mathrm{G}^{\mathrm{NSM}}(10^4\mathrm{yr})^1`$ over Galactic history, which is at the very high end of various estimates (e.g., Phinney 1991; Bethe & Brown 1998; Arzoumanian, Cordes, & Wasserman 1999). Numerical simulations of a NS-NS merger event by Rosswog et al. (1999) showed that a total mass $`M_{\mathrm{ej}}^{\mathrm{NSM}}4\times (10^3`$$`10^2)M_{}`$ of decompressed NS material may be ejected. Then the grand total from all the past NSMs over Galactic history of $`t_\mathrm{G}10^{10}`$ yr is $`M_{\mathrm{ej}}^{\mathrm{NSM}}f_\mathrm{G}^{\mathrm{NSM}}t_\mathrm{G}4\times (10^3`$$`10^4)M_{}`$. This is roughly equal to the present Galactic $`r`$-process inventory $`X_{,r}^{\mathrm{tot}}M_\mathrm{G}10^4M_{}`$, where $`X_{,r}^{\mathrm{tot}}10^7`$ is the total $`r`$-process mass fraction of nuclei with $`A100`$ in the solar system (Käppeler, Beer, & Wisshak 1989) and $`M_\mathrm{G}10^{11}M_{}`$ is the total Galactic mass in stars and gas. So it seems that NSMs could be the major $`r`$-process sources.
However, the amount of ejecta from a single NSM event discussed above is too much to explain the observed $`r`$-process abundances in metal-poor stars. This can be seen by considering mixing of the ejecta from each event with the ISM. Rosswog et al. (1999) showed that the total energy of the NSM ejecta is at most comparable to the SN explosion energy ($`10^{51}`$ erg). Therefore, when its energy/momentum is dispersed in the ISM, this ejecta can mix with at most the same amount of material as swept up by a SN remnant, $`M_{\mathrm{mix}}3\times 10^4M_{}`$ (e.g., Thornton et al. 1998). Consequently, if all of this ejecta were $`r`$-process material as required to account for the present Galactic $`r`$-process inventory, an ISM enriched by a single NSM event would have a total $`r`$-process mass fraction
$$X_{r,\mathrm{NSM}}^{\mathrm{tot}}=1.3\times 10^7\left(\frac{M_{\mathrm{ej}}^{\mathrm{NSM}}}{4\times 10^3M_{}}\right)\left(\frac{3\times 10^4M_{}}{M_{\mathrm{mix}}}\right).$$
(1)
The solar $`r`$-process mass fractions of elements with $`100A130`$ and $`A>130`$ are $`X_{,r}^{100A130}X_{,r}^{A>130}4\times 10^8`$ (Käppeler et al. 1989). According to equation (1), $`X_{r,\mathrm{NSM}}^{\mathrm{tot}}`$ would be approximately equal to $`X_{,r}^{\mathrm{tot}}X_{,r}^{100A130}+X_{,r}^{A>130}`$ even for the lowest $`M_{\mathrm{ej}}^{\mathrm{NSM}}`$ of interest. Therefore, whether the NSM ejecta were composed of $`r`$-process elements with $`100A130`$ or $`A>130`$, or both groups in solar proportion, equation (1) predicts that abundance ratios of e.g., Ag ($`A107`$) and/or Eu ($`A153`$) with respect to hydrogen in an ISM enriched by a single event would be at least comparable to the corresponding solar $`r`$-process values (Ag/H)<sub>⊙,r</sub> and (Eu/H)<sub>⊙,r</sub>. This predicted level of $`r`$-process enrichment is in disagreement with recent observations of $`r`$-process abundances in metal-poor stars, as Ag/H in stars with $`[\mathrm{Fe}/\mathrm{H}]1.3`$ to $`2.2`$ are $`10`$$`10^2`$ times lower than (Ag/H)<sub>⊙,r</sub> (Crawford et al. 1998) while Eu/H in stars with $`[\mathrm{Fe}/\mathrm{H}]3`$ are $`30`$$`10^3`$ lower than (Eu/H)<sub>⊙,r</sub> (McWilliam et al. 1995; Sneden et al. 1996, 1998).
The following interpretation of equation (1) gives some insights into how NSMs would explain the present Galactic $`r`$-process inventory and may help appreciate why this explanation is disfavored by observations. The Galaxy can be divided into $`3\times 10^6`$ units each having a mass $`M_{\mathrm{mix}}3\times 10^4M_{}`$. This division has a physical meaning as $`M_{\mathrm{mix}}`$ is the maximum mass within which the ejecta from an individual NSM event can be distributed. Due to the rather low rate of NSMs in the Galaxy, on average at most one NSM event occurred in a unit over Galactic history. Therefore, in order to account for the present Galactic $`r`$-process inventory, each unit would have to be enriched with an approximately solar $`r`$-process mass fraction by a single NSM event. As shown above, this picture of $`r`$-process enrichment is inconsistent with the observed $`r`$-process abundances in metal-poor stars. Furthermore, due to the high rate of SNe in the Galaxy, many SNe occurred in a unit where a single NSM event also occurred at sometime in Galactic history (see §4). As Fe enrichment of this unit was provided by these SNe, stars formed at different times in this unit would have varying \[Fe/H\] but either zero or approximately solar $`r`$-process mass fraction if NSMs were the major $`r`$-process sources. This is in disagreement with the observed correlation between abundances of $`r`$-process elements and Fe at $`[\mathrm{Fe}/\mathrm{H}]2.5`$ (Gratton & Sneden 1994; Crawford et al. 1998; see also McWilliam et al. 1995).
Of course, the above discussion of mixing of the NSM ejecta and the ISM is oversimplified. For example, one could imagine that a smaller than average fraction of $`r`$-process ejecta was mixed into the ISM near the edge of a NSM remnant. In this case stars formed near the edge of the remnant would have $`r`$-process mass fractions smaller than that given by equation (1). However, one would also expect that less than average enrichment was not unduly pervasive and a significant fraction of the metal-poor stars would have the $`r`$-process enrichment indicated by equation (1). The fact that no such stars have been observed suggests a difficulty of the NSM $`r`$-process model in explaining the observations at low metallicities. Furthermore, even if the observed $`r`$-process abundances in metal-poor stars could be attributed to pervasive less than average enrichment by NSMs, one would still face the other difficulty in explaining the observed correlation between abundances of $`r`$-process elements and Fe at $`[\mathrm{Fe}/\mathrm{H}]2.5`$. As Fe enrichment was controlled by SNe occurring at a much higher rate, widely-varying degrees of mixing of the $`r`$-process ejecta in an already existing NSM remnant with Fe produced by fresh SNe would result in large scatter in $`r`$-process abundances over a broad range of \[Fe/H\]. This is in disagreement with the rather early onset of the correlation between abundances of $`r`$-process elements and Fe at $`[\mathrm{Fe}/\mathrm{H}]2.5`$. In summary, observations of metal-poor stars would be difficult to explain if NSMs were the major $`r`$-process sources.
## 3. $`r`$-Process Enrichment by Supernovae
In contrast to the case of NSMs, observations of metal-poor stars as well as meteoritic data on <sup>182</sup>Hf and <sup>129</sup>I in the early solar system support a self-consistent picture of $`r`$-process enrichment by SNe (Qian & Wasserburg 2000, hereafter QW; Wasserburg & Qian 2000; hereafter WQ). In this picture the ejecta from each SN event is mixed with an average mass $`M_{\mathrm{mix}}3\times 10^4M_{}`$ of ISM swept up by the SN remnant (e.g., Thornton et al. 1998). It is assumed that the SN rate per unit mass of gas $`f_\mathrm{G}^{\mathrm{SN}}/M_{\mathrm{gas}}`$ is approximately constant over Galactic history (this seems reasonable as the star formation rate is proportional to the gas mass). Consequently, an average ISM in the Galaxy is enriched by newly-synthesized material from SNe at a frequency
$`f_{\mathrm{mix}}^{\mathrm{SN}}`$ $`=`$ $`M_{\mathrm{mix}}{\displaystyle \frac{f_\mathrm{G}^{\mathrm{SN}}}{M_{\mathrm{gas}}}}=(10^7\mathrm{yr})^1`$ (2)
$`\times `$ $`\left({\displaystyle \frac{M_{\mathrm{mix}}}{3\times 10^4M_{}}}\right)\left[{\displaystyle \frac{f_\mathrm{G}^{\mathrm{SN}}}{(30\mathrm{yr})^1}}\right]\left({\displaystyle \frac{10^{10}M_{}}{M_{\mathrm{gas}}}}\right),`$
where $`f_\mathrm{G}^{\mathrm{SN}}/M_{\mathrm{gas}}`$ is estimated using quantities for the present Galaxy.
Meteoritic data on <sup>182</sup>Hf (with a lifetime $`\overline{\tau }_{182}=1.3\times 10^7`$ yr) and <sup>129</sup>I ($`\overline{\tau }_{129}=2.3\times 10^7`$ yr) in the early solar system shed important light on the $`r`$-process and its association with SNe. Wasserburg, Busso, & Gallino (1996; see also QW) showed that the abundance ratio $`({}_{}{}^{182}\mathrm{Hf}/{}_{}{}^{180}\mathrm{Hf})_{\mathrm{SSF}}=2.4\times 10^4`$ (Harper & Jacobsen 1996; Lee & Halliday 1995, 1998) is consistent with common SNe injecting <sup>182</sup>Hf into the ISM at a high frequency $`f_{}f_{\mathrm{mix}}^{\mathrm{SN}}(10^7\mathrm{yr})^1`$ over a uniform production time $`T_{\mathrm{UP}}10^{10}`$ yr before solar system formation (SSF). However, the abundance ratio $`({}_{}{}^{129}\mathrm{I}/{}_{}{}^{127}\mathrm{I})_{\mathrm{SSF}}=10^4`$ (Reynolds 1960; see also Brazzle et al. 1999) must be accounted for by a different type of SNe occurring at a low frequency $`f_{}f_{}/10(10^8\mathrm{yr})^1`$ (Wasserburg et al. 1996; QW). Therefore, the meteoritic data require at least two distinct types of SN $`r`$-process events: the high-frequency $``$ events producing heavy nuclei with $`A>130`$ including <sup>182</sup>Hf and the low-frequency $``$ events producing light nuclei with $`A130`$ including <sup>129</sup>I. The $`r`$-process production in the SN environments associated with the $``$ and $``$ events was discussed in some detail by Qian, Vogel, & Wasserburg (1998a).
The above picture of $`r`$-process production and enrichment by SNe has clear predictions for $`r`$-process abundances resulting from a single event (QW; WQ). For example, with an average ISM enriched by the $``$ events at a frequency $`f_{}(10^7\mathrm{yr})^1`$, the solar $`r`$-process abundances of elements with $`A>130`$ such as Eu were provided by $`f_{}T_{\mathrm{UP}}10^3`$ $``$ events. This requires that the Eu abundance in an ISM enriched by a single $``$ event must be
$$\mathrm{log}ϵ_{}(\mathrm{Eu})=\mathrm{log}ϵ_{,r}(\mathrm{Eu})\mathrm{log}(f_{}T_{\mathrm{UP}})2.5,$$
(3)
where the $`\mathrm{log}ϵ`$ notation is defined as $`\mathrm{log}ϵ(\mathrm{Eu})\mathrm{log}(\mathrm{Eu}/\mathrm{H})+12`$ with Eu/H being the abundance ratio of Eu to hydrogen, and $`\mathrm{log}ϵ_{,r}(\mathrm{Eu})=0.51`$ is the value for solar $`r`$-process Eu (Käppeler et al. 1989). Similarly, the solar $`r`$-process abundances of elements with $`A130`$ such as Ag were provided by $`f_{}T_{\mathrm{UP}}10^2`$ $``$ events. This requires that the Ag abundance in an ISM enriched by a single $``$ event must be
$$\mathrm{log}ϵ_{}(\mathrm{Ag})=\mathrm{log}ϵ_{,r}(\mathrm{Ag})\mathrm{log}(f_{}T_{\mathrm{UP}})0.8,$$
(4)
where $`\mathrm{log}ϵ_{,r}(\mathrm{Ag})=1.19`$ (Käppeler et al. 1989) is used. The predictions in equations (3) and (4) are in good agreement with observations of very metal-poor stars which were formed when only a small number of SNe had contributed $`r`$-process elements to the ISM. The observed $`\mathrm{log}ϵ(\mathrm{Eu})`$ values for stars with $`[\mathrm{Fe}/\mathrm{H}]3`$ range from $`2.5`$ to $`0.9`$ (McWilliam et al. 1995; Sneden et al. 1996, 1998), which can be accounted for by $`1`$–40 $``$ events with $`\mathrm{log}ϵ_{}(\mathrm{Eu})2.5`$ from a single event (QW; WQ). In addition, Ag abundances at the level indicated by equation (4) were observed in HD 2665 and BD +371458 with $`[\mathrm{Fe}/\mathrm{H}]2`$ (Crawford et al. 1998) and in CS 22892–052 with $`[\mathrm{Fe}/\mathrm{H}]3`$ (Cowan et al. 1999).
In the above picture of $`r`$-process enrichment by SNe the total mass of $`r`$-process elements ejected in an $``$ event must be $`X_{,r}^{A>130}M_{\mathrm{mix}}/(f_{}T_{\mathrm{UP}})10^6M_{}`$, while that in an $``$ event must be $`X_{,r}^{100A130}M_{\mathrm{mix}}/(f_{}T_{\mathrm{UP}})10^5M_{}`$. A total $`10^6`$$`10^5M_{}`$ of material can be naturally provided by the neutrino-heated ejecta from the proto-neutron star in a SN over a period $`1`$–10 s (e.g., Qian & Woosley 1996). However, whether this neutrino-heated ejecta is composed of $`r`$-process material is yet to be shown. The difference by a factor $`10`$ in the total amount of $`r`$-process ejecta between the $``$ and $``$ events has been speculated to indicate that neutrino emission and hence, ejection of $`r`$-process material are terminated by the transition of the proto-neutron star into a BH in the $``$ events while prolonged ejection is sustained by neutrino emission from a stable NS in the $``$ events (Qian et al. 1998a).
In summary, despite the lack of a complete model for successful $`r`$-process production in SNe, there is a self-consistent picture of $`r`$-process enrichment by SNe that can account for the meteoritic data on <sup>182</sup>Hf and <sup>129</sup>I in the early solar system and the observed $`r`$-process abundances in metal-poor stars. Furthermore, the observed correlation between abundances of $`r`$-process elements and Fe at $`[\mathrm{Fe}/\mathrm{H}]2.5`$ (Gratton & Sneden 1994; Crawford et al. 1998; see also McWilliam et al. 1995) can be understood as the result of sufficient mixing of $`r`$-process products and Fe from multiple SN events in this picture. In fact, the same picture was used by Wasserburg & Qian (WQ) as the basic framework to explain the dispersion in abundances of heavy $`r`$-process elements such as Eu at $`[\mathrm{Fe}/\mathrm{H}]3`$.
## 4. Discussion and Conclusions
The amount of ISM to mix with the ejecta from an individual NSM event is limited by the total energy of the event. On the other hand, due to the very low rate of NSMs in the Galaxy, a large amount of $`r`$-process ejecta would be required from each event to account for the present Galactic $`r`$-process inventory. When mixed with the ISM, this required amount of ejecta would result in abundances of $`r`$-process elements with $`A130`$ (such as Ag) and $`A>130`$ (such as Eu) that are much too high (by factors $`10`$–10<sup>2</sup> for Ag and $`30`$$`10^3`$ for Eu) compared with the observed values in metal-poor stars. Furthermore, an average ISM received the ejecta from only $`1`$ NSM event over Galactic history. If $`r`$-process enrichment of the ISM was provided by NSMs in this way while Fe enrichment was provided by many SNe, there would be no correlation between abundances of $`r`$-process elements and Fe, in disagreement with the observed correlation at $`[\mathrm{Fe}/\mathrm{H}]2.5`$. While the complexities in mixing of the ejecta with the ISM can affect the above considerations in detail, it is unlikely that they can remove the difficulty of the NSM $`r`$-process model in explaining the observations of metal-poor stars (especially the rather early onset of the correlation between abundances of $`r`$-process elements and Fe at $`[\mathrm{Fe}/\mathrm{H}]2.5`$). Nevertheless, future numerical studies of $`r`$-process enrichment by NSMs accompanied by Fe enrichment by SNe should be interesting to pursue and may give a more definitive answer.
In contrast, a self-consistent picture of $`r`$-process enrichment by SNe can be obtained by considering mixing of the ejecta from an individual event with the ISM. Here an average ISM is enriched in $`r`$-process elements with $`A>130`$ by the $``$ events at a frequency $`f_{}(10^7\mathrm{yr})^1`$ and in those with $`A130`$ by the $``$ events at a frequency $`f_{}(10^8\mathrm{yr})^1`$. This picture can account for the meteoritic data on <sup>182</sup>Hf and <sup>129</sup>I in the early solar system and the observed $`r`$-process abundances in metal-poor stars. Furthermore, sufficient mixing of $`r`$-process products and Fe from multiple SN events in this picture would result in the observed correlation between abundances of $`r`$-process elements and Fe at $`[\mathrm{Fe}/\mathrm{H}]2.5`$. The same picture was also used by Wasserburg & Qian (WQ) as the basic framework to explain the dispersion in abundances of heavy $`r`$-process elements such as Eu at $`[\mathrm{Fe}/\mathrm{H}]3`$. However, a complete model of $`r`$-process and Fe production in the $``$ and $``$ events is still lacking and should be investigated by future theoretical studies.
In discussing $`r`$-process enrichment by NSMs I have assumed that the maximum amount of ISM to mix with the ejecta from an individual event is the same as swept up by a SN remnant. This is because the total energy of the NSM ejecta seen in numerical simulations (Rosswog et al. 1999) is at most comparable to the SN explosion energy ($`10^{51}`$ erg). For given conditions of the ISM, the expansion/evolution of the ejecta is essentially governed by its total energy. The large difference in the amount of ejecta between a NSM and a SN has no significant effect here as in both cases the mass of the swept-up ISM soon overwhelms that of the ejecta while the total energy remains more or less the same. I note that a small amount ($`10^5M_{}`$) of material might be ejected in highly-relativistic jets in a NS-BH merger event (Janka et al. 1999). However, the total energy of these jets is $`10^{51}`$ erg (Janka et al. 1999) and their existence would not increase significantly the amount of ISM that could mix with the entire ejecta from the event. It takes $`10^6`$ yr for the energy ($`10^{51}`$ erg) and the associated momentum of the ejecta to be dispersed in the ISM (e.g., Thornton et al. 1998), where the next NSM or SN would occur on a much longer timescale ($`10^{10}`$ yr for NSMs and $`10^7`$ yr for SNe). This leaves substantial time for mixing of the ejecta with the swept-up ISM. However, the details of the mixing process require further studies.
If, as argued here, SNe are the major sources for the $`r`$-process, then there are two possible direct observational tests: gamma rays from decay of $`r`$-process nuclei in SN remnants and surface contamination of the companion by SN $`r`$-process ejecta in binaries. Qian et al. (1998b, 1999) discussed gamma-ray signals from decay of long-lived $`r`$-process nuclei (with lifetimes $`10^3`$$`10^5`$ yr) in a nearby SN remnant and from decay of short-lived $`r`$-process nuclei (with lifetimes $`1`$–10 yr) produced in a Galactic SN that may occur in the future. The nuclide <sup>126</sup>Sn is of particular interest (Qian et al. 1998b) as its lifetime ($`10^5`$ yr) is much longer than the age ($`10^4`$ yr) of the Vela SN remnant at a distance $`250`$ pc. In the picture of $`r`$-process enrichment by SNe discussed here (see also QW; WQ) the solar $`r`$-process mass fraction of nuclei with $`A130`$ resulted from $`10^2`$ $``$ events (see §3). So a total mass $`X_{,r}^{A=126}M_{\mathrm{mix}}/10^25\times 10^7M_{}`$ of <sup>126</sup>Sn nuclei are produced in each $``$ event, where $`X_{,r}^{A=126}1.6\times 10^9`$ (Käppeler et al. 1989) is the solar $`r`$-process mass fraction of <sup>126</sup>Te, the stable daughter of <sup>126</sup>Sn. If the SN associated with the Vela remnant was an $``$ event, then decay of <sup>126</sup>Sn in this remnant would produce gamma-ray fluxes $`10^7\gamma \mathrm{cm}^2\mathrm{s}^1`$ at energies $`E_\gamma =415`$, 666, and 695 keV. Detection of these fluxes would require future gamma-ray experiments such as the proposed Advanced Telescope for High Energy Nuclear Astrophysics (ATHENA, Kurfess 1994). As the Vela remnant contains a pulsar, such detection would also provide evidence for the speculated association between $``$ events and SNe producing neutron stars (Qian et al. 1998a).
The other test mentioned above takes advantage of the occurrence of SNe in binaries. The $`r`$-process ejecta from the SN would contaminate the surface of its binary companion. Some binaries would survive the SN explosion and acquire a NS or a BH in place of the SN progenitor. Therefore, an ordinary star observed to be the binary companion of a NS or a BH might show $`r`$-process abundance anomalies on the surface. To estimate the plausible level of such anomalies, I assume that a fraction $`10^3`$ of the entire SN ejecta ($`10M_{}`$, mostly non-$`r`$-process material) would be intercepted by a low mass ($`1M_{}`$) companion and then mixed with $`10^2M_{}`$ of the surface material. If the SN was an $``$ event, $`10^9M_{}`$ of $`r`$-process elements with $`A>130`$ would be intercepted, while for an $``$ event $`10^8M_{}`$ of $`r`$-process elements with $`A130`$ would be intercepted (see §3). These quantities are to be compared with $`4\times 10^{10}M_{}`$ of the corresponding $`r`$-process elements in $`10^2M_{}`$ of the surface material in a companion star of solar $`r`$-process composition. So a SN in a binary could enhance significantly the surface $`r`$-process abundances in the companion star, especially if the SN was an $``$ event. In view of the large overabundance of O, Mg, Si, and S recently observed in the companion star of a probable BH (Israelian et al. 1999), detection of $`r`$-process enhancement in similar binary systems seems promising. Such detection may also test directly the speculation by Qian et al. (1998a) that $``$ events are associated with SNe producing BHs while $``$ events are associated with SNe producing neutron stars.
I thank Petr Vogel and Jerry Wasserburg for many discussions on the $`r`$-process. I am also grateful to the referee, Friedel Thielemann, for detailed criticisms that help improve the paper. This work was supported in part by the Department of Energy under grant DE-FG02-87ER40328.
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# Introduction
## Introduction
Dans la première partie de cet article, on étudie le schéma de Hilbert des courbes rationnelles lisses tracées sur une variété homogène. Soient $`G`$ un groupe algébrique simple et simplement connexe sur $`𝐂`$, $`P`$ un parabolique de $`G`$ et $`\alpha A_1(G/P)`$ une classe d’équivalence rationnelle de $`1`$-cycle. On suppose $`\alpha `$ positive (c’est à dire coupant positivement les classes de diviseurs effectifs de $`G/P`$), alors J.F. Thomsen \[T\], B. Kim et R. Pandharipande \[KP\] ont montré (en particulier) que le schéma de Hilbert des courbes lisses rationnelles de classe $`\alpha `$ tracées sur $`G/P`$ est lisse connexe mais pas nécessairement non vide. Dans cette partie, nous construisons de manière effective des courbes lisses sur les variétés homogènes. Nous retrouvons ainsi le résultat de J.F. Thomsen, B. Kim et R. Pandharipande par une méthode complètement différente et montrons l’existence de courbes lisses.
Théorème 1 : (ı) Le schéma de Hilbert des courbes rationnelles lisses tracées sur $`G/P`$ de classe $`\alpha `$ positive est irréductible et lisse.
(ıı) Si $`\alpha `$ est strictement positive (c’est à dire d’intersection strictement positive avec tous les diviseurs effectifs de $`G/P`$), alors il existe une courbe rationnelle lisse de classe $`\alpha `$ sur $`G/P`$ sauf pour $`𝐏^1`$, $`𝐏^2`$ et $`𝐏^1\times 𝐏^1`$.
On donnera également les cas pour lesquels le schéma de Hilbert des courbes lisses est non vide lorsque $`\alpha `$ est positive non strictement positive.
Pour démontrer ce théorème on étudie les orbites de $`G/P`$ sous un parabolique $`P^{}`$. Ces orbites que nous appelerons $`P^{}`$-orbites généralisent les cellules de Schubert classiques (définies par l’action d’un Borel). Elles n’apparaissent pas de façon systématique dans la littérature. Seuls quelques cas particuliers sont décrits par exemple dans \[BGG\], \[K1\], \[K2\], \[T\] et \[LMS\]. Les méthodes de J.F. Thomsen \[T\] et B. Kim et R. Pandharipande \[KP\] pour prouver la première partie du théorème $`1`$ sont totalement différentes de la notre. Ils montrent que le schéma des applications stables est connexe en reliant toute application à une application du bord par action d’un Borel ou d’un tore maximal, puis irréductible par lissité. La méthode présentée ici est plus directe (pas d’étude du bord du schéma des applications stables) et elle permet de montrer l’existence de morphismes non constants de $`𝐏^1`$ dans $`G/P`$ de classe $`\alpha `$ positive et même de courbes lisses ce qui n’était pas le cas des méthodes de \[T\] et \[KP\]. Pour les courbes de genre $`g>0`$, on n’a pour le moment que des résultats partiels. Citons E. Ballico \[B\] qui a montré l’irréductibilité du schéma de Hilbert des courbes de degré $`d`$ et de genre $`g`$ tracées sur une quadrique de $`𝐏^n`$ pour $`n7`$ et $`d2g1`$. Notre méthode permet de montrer, dans le cas où $`G`$ est $`SL_n`$ ou $`SO_{2n}`$ et si $`P`$ est un parabolique maximal, que le schéma des morphismes de degré $`d`$ d’une courbe de genre $`g`$ vers $`G/P`$ est irréductible dès que $`d`$ est grand devant $`g`$ (on en déduit le résultat pour le schéma de Hilbert associé). On peut également montrer ce résultat si $`G`$ est $`Sp_{2n}`$ et si $`P`$ est un parabolique maximal qui ne correspond pas aux sous espaces totalement isotropes maximaux. Enfin, la méthode de E. Ballico se généralise aux quadriques lisses de rang $`n`$ pour $`3n7`$. Dans tous ces cas excepté pour $`𝐏^2`$ et $`𝐏^1\times 𝐏^1`$, on sait également montrer l’existence de courbes lisses.
De plus, les $`P^{}`$-orbite que l’on introduit dans cette construction nous permettent, dans la seconde partie de cet article, de donner une désingularisation $`\pi `$ des variétés de Schubert plus fine que celle de M. Demazure \[D\]. En effet, elle sera bijective sur un plus grand ouvert et une désingularisation de Demazure se factorise par celle-ci. Cependant, je ne sais pas si elle est la plus fine possible c’est à dire si elle est bijective sur le lieu lisse des variétés de Schubert. On peut néanmoins vérifier que c’est le cas pour les groupes de petite dimension ($`SL_4`$ par exemple). On montre ainsi le résutat suivant :
Théorème 2 : (ı) Une désingularisation de Demazure se factorise par $`\pi `$.
(ıı) Le morphisme $`\pi `$ est une désingularisation des variétés de Schubert.
On donnera aussi une condition suffisante (non nécessaire) de lissité des variétés de Schubert et un critère pour qu’une variété de Schubert soit homogène sous l’action d’un sous groupe de $`G`$.
CONSTRUCTION DE COURBES LISSES SUR LES VARIÉTÉS HOMOGÈNES
On adoptera tout au long de cet article les notations de W. Fulton et J. Harris \[FH\] pour tout ce qui concerne les groupes, leurs algèbres de Lie et leurs systèmes de racines. Soit $`G`$ un groupe de Lie simple et simplement connexe et soit $`𝔤`$ son algèbre de Lie. Soit $`𝔥`$ une algèbre de Cartan de $`𝔤`$ et $`𝔟`$ un borel de $`𝔤`$ contenant $`𝔥`$. On note $`R`$, respectivement $`R^+`$ et $`R^{}`$ l’ensemble des racines, respectivement l’ensemble des racines positives et négatives. On a les décompositions de $`𝔤`$ et $`𝔟`$ suivant les poids :
$$𝔤=𝔥_{\alpha 𝔥^{}}𝔤_\alpha 𝔟=𝔥_{\alpha R^+}𝔤_\alpha $$
Le Groupe de Picard de $`G/B`$ est exactement le réseau des poids de $`𝔥^{}`$. Dans le groupe de Picard de $`G/B`$, le cône formé par les poids $`x`$ tels que : $`(x,\alpha )0`$ pour toute racine simple $`\alpha `$ (c’est la chambre de Weyl principale) est engendré par les poids fondamentaux (diviseurs engendrant le groupe de Picard des $`G/P`$$`P`$ est un parabolique maximal). On l’appelle cône ample de $`\mathrm{Pic}(G/B)`$. Le groupe $`A_1(G/B)`$ est le dual de $`\mathrm{Pic}(G/B)`$. C’est donc le réseau des racines dans $`𝔥`$. Le cône ample de $`\mathrm{Pic}(G/B)`$ définit par dualité le cône positif de $`A_1(G/B)`$. Ce cône est le cône engendré par les racines simples. On appelle cône strictement positif les éléments qui sont strictement positifs sur toutes les arêtes du cône ample du groupe de Picard, c’est le cône positif privé de ses faces de codimension $`1`$.
Un sous groupe parabolique $`P`$ de $`G`$ contenant $`B`$ est donné par un ensemble de racines simples négatives $`\mathrm{\Sigma }`$ et l’algèbre correspondante est alors :
$$𝔭(\mathrm{\Sigma })=𝔟_{\alpha T(\mathrm{\Sigma })}𝔤_\alpha $$
$`T(\mathrm{\Sigma })`$ est l’ensemble des racines obtenues comme somme de racines simples en dehors de $`\mathrm{\Sigma }`$. Le groupe de Picard de $`G/P(\mathrm{\Sigma })`$ est donné par le sous réseau des poids (c’est à dire le sous réseau de $`𝔥^{}`$) défini par les équations $`(x,\alpha )=0`$ pour toute les racines simples $`\alpha `$ en dehors de $`\mathrm{\Sigma }`$. On peut aussi le décrire comme le réseau engendré par les coracines $`\stackrel{ˇ}{\alpha }`$ des racines simples $`\alpha \mathrm{\Sigma }`$. On peut ainsi définir par restriction un cône ample dans $`\mathrm{Pic}(G/P)`$. Par dualité on peut définir des cônes positif et strictement positif dans $`A_1(G/P)`$ qui sont les quotients des cônes de $`A_1(G/B)`$.
Remarque 1 : Si $`C`$ est une courbe irréductible dans $`G/P`$ alors sa classe $`[C]`$ dans $`A_1(G/P)`$ est nécessairement dans le cône positif car tous les diviseurs $`D`$ qui forment une arête du cône ample de $`\mathrm{Pic}(G/P)`$ sont effectifs et par l’action du groupe on voit que $`(C,D)0`$ (voir par exemple \[Kl\]). Ainsi \[C\] est toujours dans le cône positif et la condition de la première partie du théorème $`1`$ est une condition nécessaire pour que le schéma de Hilbert soit non vide.
## 1 Les courbes rationnelles et leur schéma de Hilbert
On note $`\mathrm{𝐇𝐢𝐥𝐛}(\alpha ,X)`$ le schéma de Hilbert des courbes rationnelles lisses dans $`X`$ dont dans la classe dans $`A_1(X)`$ est $`\alpha `$.
Proposition 1 : Soit $`\alpha A_1(G/P)`$ dans le cône positif, alors $`\mathrm{𝐇𝐢𝐥𝐛}(\alpha ,G/P)`$ est lisse de dimension :
$$\underset{\alpha ^{}\mathrm{\Sigma }}{}(\stackrel{ˇ}{\alpha }^{},\alpha )+\mathrm{dim}(G/P)3$$
Démonstration: Il suffit de montrer que $`T_{G/P}`$ est engendré par ses sections pour avoir la lissité du schéma de Hilbert. Sa dimension au point $`C`$ est alors donnée par $`\mathrm{deg}(T_{G/P}|_C)+\mathrm{dim}(G/P)3`$ c’est ce qu’on cherche car le degré de $`T_{G/P}|_C`$ est $`_{\alpha ^{}\mathrm{\Sigma }}(\stackrel{ˇ}{\alpha }^{},\alpha )`$.
Si on a une action linéaire de $`P`$ sur un espace vectoriel $`N`$, alors $`P`$ agit sur le produit $`G\times N`$ par $`(p,(g,n))(gp^1,pn)`$. Le quotient $`G\times ^PN`$ est un fibré vectoriel au dessus de $`G/P`$. Ainsi, le faisceau $`T_{G/P}`$ est le fibré vectoriel associé à la représentation de $`P`$ suivante : $`𝔤/𝔭`$. Mais $`G`$ agit sur $`𝔤`$ donc le fibré vectoriel $`𝒢`$ associé à $`𝔤`$ est trivial. Comme quotient du fibré trivial $`𝒢`$, le faisceau $`T_{G/P}`$ est engendré par ses sections. Pour une autre démonstration voir \[Ko\] théorème $`1.4`$ page 241.
La difficulté de la première partie du théorème $`1`$ réside donc dans l’irréductibilité du schéma de Hilbert. Pour l’étudier, on va s’intéresser au schéma des morphismes $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,G/P)`$ qui paramétrise les morphismes $`f`$ de schémas de $`𝐏^1`$ dans $`G/P`$ tels que la classe de $`f(𝐏^1)`$ dans $`A_1(G/P)`$ est $`\alpha `$. L’ensemble des plongements de $`𝐏^1`$ dans $`G/P`$ de classe $`\alpha `$ est un ouvert de ce schéma. Cet ouvert domine $`\mathrm{𝐇𝐢𝐥𝐛}(\alpha ,G/P)`$ (c’est une fibration en $`SL_2`$) et l’irréductibilité du schéma de Hilbert se déduira de celle de $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,G/P)`$. La proposition suivante nous permet de ramener ce problème à celui d’un ouvert :
Proposition 2 : Soit $`X`$ une variété munie d’une action transitive d’un groupe $`G`$ et soit $`\alpha A_1(X)`$. Supposons qu’il existe un ouvert $`U`$ de $`X`$ dont le complémentaire $`Z`$ est de codimension supérieure ou égale à $`2`$ et tel que $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,U)`$ est irréductible, alors $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$ l’est aussi.
Démonstration: Les ouverts $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,g.U)`$ de $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$ recouvrent $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$. On utilise pour celà un résultat de Kleiman \[Kl\] : si $`C`$ est une courbe de $`X`$ et si $`Z`$ est un fermé de codimension au moins $`2`$, alors il existe un ouvert de $`G`$ tel que, pour tout point $`g`$ de cet ouvert, l’intersection de $`C`$ avec les translatés $`g.Z`$ du fermé $`Z`$ est vide. Donc si on a un point $`f`$ dans $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$ et $`C`$ son image, alors il existe un ouvert de $`G`$ tel que $`C`$ est contenue dans $`g.U`$ et donc $`f`$ est dans $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,g.U)`$ pour tout $`g`$ dans cet ouvert. Ceci impose l’irréductibilité. En effet, deux ouverts $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,g.U)`$ et $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,g^{}.U)`$ se coupent toujours : il suffit d’exhiber une courbe qui ne rencontre pas $`g.Zg^{}.Z`$. Soit donc $`f\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$ quelconque et $`C`$ son image, le même résultat de Kleiman \[Kl\] nous dit qu’il existe un ouvert de $`G`$ tel que pour tout élément $`g^{\prime \prime }`$ de cet ouvert la courbe $`g^{\prime \prime }.C`$ ne rencontre pas $`g.Zg^{}.Z`$. Toutes ces courbes $`g^{\prime \prime }.C`$ sont donc dans l’intersection. Supposons maintenant que $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$ a plusieures composantes irréductibles et soient $`𝐇`$ et $`𝐇^{}`$ deux telles composantes. Comme les ouverts $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,g.U)`$ recouvrent et sont irréductibles, il existe deux éléments $`g`$ et $`g^{}`$ de $`G`$ tels que $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,g.U)`$ est un ouvert non vide de $`𝐇`$ et $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,g^{}.U)`$ est un ouvert non vide de $`𝐇^{}`$. Mais alors on sait que $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,g.U)\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,g^{}.U)`$ est non vide, c’est donc un ouvert dense de $`𝐇`$ et de $`𝐇^{}`$ ce qui impose l’égalité de ces composantes.
Le principe de la démonstration sera le suivant : on commence par se rammener grace à la proposition précédente à des ouverts de $`G/P`$ dont le complémentaire est de codimension supérieure ou égale à $`2`$ : les $`P^{}`$-orbites (voir paragraphe suivant). On dévisse ensuite ces $`P^{}`$-orbites en variétés plus simples pour lesquelles on sait résoudre le problème. Deux cas se présenterons alors selon que $`P`$ est un Borel ou non :
\- Un fibré en droites projectives au dessus d’une variété homogène pour laquelle on sait résoudre le problème. Dans ce cas on aura besoin d’une condition sur le degré de la courbe par rapport à cette fibration.
\- Une tour de fibrés affines, de fibrés vectoriels direction engendrés par leurs sections, au dessus d’un produit de variétés homogènes sous des groupes dont le diagramme de Dynkin est de longueur strictement inférieure à celle du diagramme de Dynkin de $`G`$.
Pour chacun de ces cas on a une proposition qui permet d’obtenir le résultat. On commence par définir ce qu’on appelle une tour de fibrés affines.
Définition : Un morphisme $`X\stackrel{f}{}Y`$ est appelé tour de fibrés affines si $`f`$ se décompose en $`X\stackrel{f_1}{}X_1\mathrm{}\stackrel{f_n}{}X_n\stackrel{f_{n+1}}{}Y`$ où les $`f_i`$ sont des fibrés affines.
Proposition 3 : Si $`X\stackrel{\phi }{}Y`$ est une tour de fibrés affines de fibrés vectoriels direction engendrés par leurs sections, si $`\alpha A_1(X)`$, alors $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)\mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$ est une tour de fibrés affines. En particulier, si $`\mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$ est irréductible, alors $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$ l’est aussi.
Démonstration: Il suffit par récurrence de prouver le cas d’un fibré affine $``$ dont le fibré vectoriel direction $`F`$ est engendré par ses sections. On a le diagramme suivant :
$$\begin{array}{ccc}\mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)\times 𝐏^1& \stackrel{p}{}& Y\\ q& & \\ \mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)& & \end{array}$$
Le faisceau $`F`$ étant engendré par ses section le faisceau $`R^1q_{}p^{}F`$ est nul et $`q_{}p^{}F`$ est localement libre. On va montrer le résultat suivant : le schéma $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$ au dessus de $`\mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$ est le fibré affine de base $`q_{}p^{}F`$ associé à l’élément de $`H^1q_{}p^{}F`$ image de celui de $`H^1F`$ définissant $``$.
En effet, on a un morphisme universel $`𝐏^1\times \mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)Y\times \mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$ et si on fait le produit fibré avec $`X\times \mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$ on obtient une variété $`Z`$ qui est le fibré affine $`p^{}`$ au dessus de $`𝐏^1\times \mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$. La variété $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$ est alors donnée sur chaque fibre au dessus de $`\mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$ par les sections de ce fibré, c’est à dire par le faisceau d’espaces affines $`q_{}p^{}`$ (défini de la même façon que pour un fibré vectoriel). Comme $`q_{}p^{}F`$ est localement libre, $`q_{}p^{}`$ est un fibré affine de base $`q_{}p^{}F`$, il est associé à l’élément de $`H^1q_{}p^{}F=H^1p^{}F`$ image de celui de $`H^1F`$ définissant $``$.
Proposition 4 : Soit $`X\stackrel{\phi }{}Y`$ un fibré en droites projectives de fibré tangent relatif $`T_{X/Y}`$ et soit $`\alpha A_1(X)`$ tel que $`\alpha T_{X/Y}0`$, alors $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$ est un ouvert d’un fibré projectif au dessus de $`\mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$. En particulier, si $`\mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$ est irréductible alors $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$ l’est aussi.
Démonstration: Soit $`E`$ un faisceau localement libre de rang $`2`$ sur $`Y`$ tel que $`X=\mathrm{Proj}_Y(\mathrm{Sym}(E))`$. On a encore un morphisme de $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$ dans $`\mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$. Soit $`f\mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$ et soit $`S=\mathrm{Proj}_{𝐏^1}(\mathrm{Sym}(f^{}E))`$ la surface réglée obtenue comme produit fibré de $`𝐏^1`$ avec $`X`$ au dessus de $`Y`$. On note $`𝒪_S(\mathrm{0,1})`$ le quotient tautologique de cette surface réglée et $`𝒪_S(\mathrm{1,0})`$ le diviseur d’une fibre (le groupe de Picard de $`S`$ est $`𝐙^2`$). Les relèvements de $`f`$ dans $`X`$ de classe $`\alpha `$ correspondent aux sections de $`E(c)`$ (où $`c`$ est un entier dépendant de $`\alpha `$) qui sont partout injectives. Les sections de $`E(c)`$ correspondent exactement aux sections de $`𝒪_S(c\mathrm{,1})`$. L’entier $`c`$ vérifie la condition suivante :
$$\left(c\mathrm{,1}\right)\left(\begin{array}{cc}0& 1\\ 1& a\end{array}\right)(\begin{array}{c}a\\ 2\end{array})=b$$
où la matrice est celle de la forme d’intersection sur $`S`$, le vecteur de droite est la classe de $`T_{X/Y}`$, $`a=(\phi _{}\alpha )c_1(E)`$ et $`b=\alpha T_{X/Y}`$. On voit ainsi que $`c=\frac{1}{2}(ba)`$ (nécessairement $`ba[2]`$). On s’intéresse donc au faisceau $`F=q_{}(p^{}E𝒪_{𝐏^1}(c))`$ sur $`\mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$ (avec les notations de la proposition précédente pour l’incidence avec $`Y`$). Il est localement libre de rang $`b+2`$ au dessus de l’ouvert $`U`$ complémentaire de $`\mathrm{Supp}(R^1q_{}((p^{}E)𝒪_{𝐏^1}(c1)))`$. Cet ouvert contient l’image de $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$. En effet, si $`f\mathrm{Supp}(R^1q_{}((p^{}E)𝒪_{𝐏^1}(c1)))`$, alors le faisceau $`f^{}E(c)`$ a un facteur de degré inférieur à $`1`$ et ses sections s’annulent toutes et ne peuvent ainsi donner des courbes sur $`X`$. Au dessus d’un point de l’ouvert $`U`$, les éléments de la fibre sont donnés par les sections partout injectives qui forment un ouvert des sections. Sur $`\mathrm{Proj}_U(\mathrm{Sym}(\stackrel{ˇ}{F}))\times 𝐏^1`$ on a une section tautologique de $`p^{}E𝒪_{𝐏^1}(c)`$, soit $`Z`$ la projection dans $`\mathrm{Proj}_U(\mathrm{Sym}(\stackrel{ˇ}{F}))`$ du lieu des zéros de cette section. On peut maintenant affirmer que $`\mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$ s’identifie à l’ouvert de $`\mathrm{Proj}_U(\mathrm{Sym}(\stackrel{ˇ}{F}))`$ complémentaire $`Z`$.
Remarquons qu’il est possible que $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,X)`$ soit vide si l’ouvert image $`U`$ dans $`\mathrm{𝐇𝐨𝐦}_{\phi _{}\alpha }(𝐏^1,Y)`$ est vide.
Cette proposition est encore vraie si on remplace le fibré en droites projectives par un fibré en espaces projectifs de dimension supérieure.
## 2 Décomposition en $`P^{}`$-orbites
### 2.1 Définition et propriétés générales
On va ici introduire une classe de variétés plus générales que les cellules de Schubert classiques. Ces variétés n’apparaissent pas ou seulement partiellement dans la littérature : certains auteurs en ont décrit des cas particuliers (Kempf \[K1\] et \[K2\], Bernstein, Gelfand et Gelfand \[BGG\], Lakshmibai, Musili et Seshadri \[LMS\], Thomsen \[T\]). On fixe maintenant un tore $`T`$ et deux paraboliques $`P`$ et $`P^{}`$ contenant ce tore et dont l’intersection contient un Borel (que nous ne fixons pas à priori).
Définition : On appelle $`P^{}`$-orbite de $`G/P`$ les orbites de $`G/P`$ sous l’action de $`P^{}`$.
Remarques 2 : (ı) Les $`P^{}`$-orbites de $`G/P`$ forment une stratification de $`G/P`$, elles sont isomorphes à $`P^{}wP/P`$ ou à $`w(P^{})/(w(P^{})P)`$$`wW`$ ($`W`$ est le groupe de Weyl de $`G`$). Si on note pour tout parabolique $`Q`$, $`W(Q)`$ le sous groupe du groupe de Weyl qui laisse stable le parabolique $`Q`$ (si on fixe un Borel $`B_0`$ dans $`Q`$, ce groupe peut être vu comme le sous groupe de $`W`$ engendré par les réflexions par rapport aux racines de $`B_0`$ dont les opposés sont dans $`Q`$), alors les $`P^{}`$-orbites de $`G/P`$ sont paramétrées par les orbites de $`W`$ sous l’action de $`W(P)\times W(P^{})`$ donnée de la façon suivante : $`(g,(w,w^{}))wgw_{}^{}{}_{}{}^{1}`$. La $`P^{}`$-orbite $`P^{}wP/P`$ est réunion disjointe de cellules de Schubert, ceci vient de l’écriture
$$P^{}wP=\underset{(w_1,w_2)W(P^{})\times W(P)}{}Bw_1ww_2B$$
avec $`B`$ un Borel de $`PP^{}`$.
(ıı) Les cellules de Schubert classiques sont décrites par les variétés $`BwB/B`$$`B`$ est un Borel de $`G`$ et $`wW`$. Ce sont les $`B`$-orbites de $`G/B`$.
Des cas plus généraux ont étés décrits par exemple dans Kempf \[K1\] et \[K2\], Berstein, Gelfand et Gelfand \[BGG\], Lakshmibai, Musili, Seshadri \[LMS\] ou Thomsen \[T\] qui étudient les variétés $`B^{}wP/P`$ qui sont les cellules de Schubert de $`G/P`$ par rapport à $`B^{}`$ (cas où $`P^{}=B^{}`$ est un Borel contenu dans $`P`$) ou parfois leurs symétriques : $`P^{}wB/B`$ qui sont des $`P^{}`$-orbites plus grandes que les cellules classiques (cas où $`P=B`$ est un Borel contenu dans $`P^{}`$).
Exemple 1 : Les $`P^{}`$-orbites de $`G/P`$ sont des ouverts lisses de $`G/P`$. Le lemme suivant nous permet de montrer que pour toute $`P^{}`$-orbite $`P^{}\overline{w}P/P`$ de $`G/P`$ ($`\overline{w}`$ est une orbite de $`W`$ sous l’action de $`W(P)\times W(P^{})`$), il existe un Borel $`B`$ de $`P`$ et un élément $`w^{}\overline{w}`$ tel que $`P^{}\overline{w}P/P`$ contient la cellule de Schubert $`Bw^{}P/P`$ et est contenue dans l’adhérence de cette cellule. On voit ainsi que l’adérence de la $`P^{}`$-orbite $`P^{}\overline{w}P/P`$ est une variété de Schubert.
Lemme 1 : Soient $`P`$ et $`P^{}`$ deux paraboliques contenant un même Borel, soit $`\overline{w}W/(W(P)\times W(P^{}))`$, il existe un Borel $`BP`$ et un élément $`w^{}\overline{w}`$ tels que $`P^{}\overline{w}P\overline{Bw^{}B}`$.
Démonstration: On commence par montrer que l’espace tangent de $`P^{}\overline{w}P`$ en $`w`$ est $`𝔭+𝔭^{}`$$`𝔭`$ est l’espace tangent de $`P`$ en $`e`$ (l’identité) et $`𝔭^{}`$ est l’espace tangent de $`w(P^{})`$ en $`e`$. En effet, l’espace tangent de $`P^{}we=w.w(P^{})`$ en $`w`$ est le translaté à gauche par $`w`$ de $`𝔭^{}`$ et l’espace tangent de $`ewP`$ en $`w`$ est le translaté à gauche par $`w`$ de $`𝔭`$. Ainsi $`𝔭+𝔭^{}`$ est contenu dans l’espace tangent de $`P^{}wP`$ en $`w`$. Mais $`P^{}wP/P`$ s’identifie à $`w(P^{})/(w(P^{})P)`$ donc $`P^{}wP`$ est de dimension $`\mathrm{dim}(𝔭+𝔭^{})`$. La lissité de $`P^{}wP`$ nous permet de conclure que $`\mathrm{T}_e(P^{}wP)=𝔭+𝔭^{}`$.
Ainsi, notre lemme se ramène à montrer qu’il existe deux Borels $`BP`$ et $`B^{}w(P^{})`$ tels que $`𝔟+𝔟^{}=𝔭+𝔭^{}`$. En effet, le Borel $`w(B)`$ est contenu dans $`w(P^{})`$ donc il existe $`w^{\prime \prime }W(P^{})`$ tel que $`w^{\prime \prime }w(B)=B^{}`$ donc $`B^{}=w^{}(B)`$ pour $`w^{}\overline{w}`$. On aura alors $`Bw^{}B=w^{}B^{}BP^{}wP`$ mais comme ils ont le même espace tangent, on aura aussi $`P^{}wP\overline{Bw^{}B}`$.
On utilise l’abus de notation suivant : si $`\alpha `$ est une racine et $`𝔭`$ l’algèbre de Lie d’un parabolique, on dit que $`\alpha 𝔭`$ si $`\alpha `$ est une valeur propre pour l’action du tore $`𝔥`$ sur $`𝔭`$ ou encore si $`𝔤_\alpha `$ apparait dans la décomposition de $`𝔭`$. Soient $`𝔟`$ et $`𝔟^{}`$ deux Borels de $`𝔭`$ et $`𝔭^{}`$, on modifie ces Borels pour obtenir la propriété recherchée. Si il existe $`\alpha 𝔭+𝔭^{}`$ (disons dans $`𝔭`$) telle que $`\alpha 𝔟+𝔟^{}`$ alors il existe un racine simple $`\alpha _0`$ de $`𝔟`$ telle que $`\alpha _0𝔭`$ et $`\alpha _0𝔟+𝔟^{}`$. En effet sinon pour toute racine simple $`\alpha _i`$ de $`𝔟`$ telle que $`\alpha _i𝔭`$ on a $`\alpha _i𝔟+𝔟^{}`$ i.e. $`\alpha _i𝔟^{}`$. Mais alors $`\alpha `$ s’écrit $`\alpha =(\alpha _i)`$ où les $`\alpha _i`$ sont des racines simples de $`𝔟`$ telles que $`\alpha _i𝔟^{}`$ donc $`\alpha 𝔟^{}`$ ce qui est absurde.
On remplace maintenant $`𝔟`$ par $`s_{\alpha _0}(𝔟)`$ et on a :
$$s_{\alpha _0}(𝔟)+𝔟^{}=(𝔟+𝔟^{})\{\alpha _0\}$$
on se ramène ainsi à des Borels $`𝔟`$ et $`𝔟^{}`$ tels que $`𝔟+𝔟^{}=𝔭+𝔭^{}`$.
Ce lemme montre également l’existence d’un Borel $`B`$ et d’un élément $`w^{}\overline{w}`$ tel que la flèche $`\overline{Bw^{}B/B}\overline{P^{}wP/P}`$ entre variétés de Schubert est une fibration en $`P/B`$.
Supposons que $`P=B`$ est un Borel, les cellules de Schubert classiques de $`G/B`$ sont munies d’un ordre partiel (ordre de Bruhat). Les $`P^{}`$-orbites de $`G/B`$ décrites par $`P^{}\overline{w}B/B`$ avec $`\overline{w}W/W(P^{})`$ contiennent comme ouvert dense une cellule de Schubert classique $`Bw^{}B/B`$ qui est donnée par $`w^{}`$ (celui du lemme $`1`$) qui est l’élément maximal de $`\overline{w}`$ dans l’ordre de Bruhat. Cet élément est unique car les cellules associées à deux tels éléments ont la même adhérence : celle de $`P^{}\overline{w}B/B`$.
On va utiliser ces $`P^{}`$-orbites pour construire des ouverts de $`G/P`$ dont le complémentaire est de codimension supérieure ou égale à $`2`$. On commence par montrer un résultat sur ces $`P^{}`$-orbites qui au vu de la proposition $`3`$ nous permettra de faire fonctionner la récurrence.
Définition : Les composantes connexes du diagramme de Dynkin de $`G`$ privé de $`P^{}`$ sont les diagrammes de Dynkin obtenus à partir de celui de $`G`$ en enlevant les sommets correspondant à $`P^{}`$.
Exemple 2 : Dans $`SL_4`$, si $`P^{}`$ est le parabolique correspondant aux droites, il y a deux composantes connexes du diagramme de Dynkin de $`G`$ privé de $`P^{}`$. Ces deux composantes sont isomorphes au diagramme de Dynkin de $`SL_2`$.
Remarques 3 : (ı) Si $`N`$ est un espace vectoriel muni d’une action de $`P`$, alors on peut définir un fibré vectoriel $`𝒩`$ sur $`G/P`$ à partir de $`N`$. Ce fibré est le produit contracté $`G\times ^PN`$. De la même façon, si $`X`$ est une variété affine sur laquelle $`P`$ agit à gaughe, on peut définir le produit contracté $`G\times ^PX`$ qui est le quotient de $`G\times X`$ par $`P`$ (voir par exemple \[DG\]). Si l’action de $`P`$ sur $`X`$ se prolonge en une action de $`G`$, alors $`G\times ^PX`$ (et en particulier $`𝒩`$ dans le cas d’un fibré) est trivial sur $`G/P`$.
(ıı) Soit $`wW`$. On note $`N^{}`$ la partie unipotente de $`w(P^{})`$ et $`R^{}`$ le quotient réductif correspondant. La décomposition de Levi (voir \[Bo\] 11.22 et 14.17-19) nous permet de dire que $`w(P^{})`$ est le produit semi-direct de $`R^{}`$ et $`N^{}`$. Cette situation est rigidifiée par le choix du tore maximal $`w(T)`$ contenu dans $`w(P^{})`$ : la décomposition de $`𝔤`$ en sous espaces propes induit une section $`s`$ de $`R^{}`$ dans $`w(P^{})`$. On note $`N=N^{}P`$ et $`R=R^{}P`$ (c’est l’image de $`P`$ dans $`R^{}`$). Le groupe $`Pw(P^{})`$ est produit semi-direct de $`R`$ et $`N`$. On a ainsi deux fibrés vectoriels $`𝒩^{}`$ et $`𝒩`$ sur $`R^{}/R`$.
(ııı) Le groupe $`R^{}`$ est produit des groupes $`G_i`$ donnés par les composantes du diagramme de Dynkin de $`G`$ privé de $`P^{}`$ et de facteurs isomorphes à $`𝐂^{}`$ (il y en a $`\mathrm{rang}(\mathrm{Pic}(G/P^{}))`$). De même, $`R`$ est produit de paraboliques $`P_i`$ des $`G_i`$ et des mêmes facteurs isomorphes à $`𝐂^{}`$. Les $`P_i`$ sont donnés par les sommets de $`P`$ regardés dans le diagramme de Dynkin de $`G`$ prové de $`P^{}`$.
Proposition 5 : Le morphisme naturel $`f`$ de $`P^{}wP/P`$ vers $`R^{}/R`$ est une tour de fibrés affines dont les fibrés vectoriels direction correspondent à la décomposition naturelle de $`𝒩^{}/𝒩`$. Ils sont définis sur $`R^{}/R`$ et sont engendrés par leurs sections.
Démonstration: On regarde ici la $`P^{}`$-orbite $`P^{}wP/P`$ sous la forme $`w(P^{})/(w(P^{})P)`$. Soit $`1=N_0^{}N_1^{}\mathrm{}N_n^{}=N^{}`$ la suite centrale ascendante de l’unipotent $`N^{}`$. On note $`Z_i^{}=N_i^{}/N_{i1}^{}`$ qui est le centre de $`N^{}/N_{i1}^{}`$. C’est un espace vectoriel sur lequel $`R^{}`$ agit linéairement. Posons $`P_i^{}=w(P^{})/N_i^{}`$, $`P_i`$ l’image de $`P`$ dans $`P_i^{}`$ et $`Z_i=P_iZ_i^{}`$ qui est un espace vectoriel sur lequel $`R`$ agit linéairement. L’écriture de $`f`$ comme tour de fibré affine est $`w(P^{})/(w(P^{})P)=P_0^{}/P_0P_1^{}/P_1\mathrm{}P_i^{}/P_i\mathrm{}P_n^{}/P_n=R^{}/R`$. Les flèches sont des fibrations localement triviales pour la topologie de Zariski et la fibre de $`P_{i1}^{}/P_{i1}P_i^{}/P_i`$ au dessus du point $`1`$ est $`Z_i^{}/Z_i`$. Il sagit de montrer que cette fibration est affine de fibré vectoriel direction l’image réciproque (par la projection $`P_i^{}/P_iR^{}/R`$) du fibré vectoriel sur $`R^{}/R`$ déduit de la représentation $`Z_i^{}/Z_i`$ de $`R`$.
Notons $`A=P_{ii}^{}`$, $`B=P_{i1}`$ et $`U=Z_i^{}`$. La fibration $`P_{i1}^{}/P_{i1}P_i^{}/P_i`$ est donc $`A/B(A/U)/(B/(UB))=A/BU`$. On va montrer le lemme suivant :
Lemme 2 : Soient $`A`$ un groupe linéaire, $`B`$ et $`U`$ des sous groupes fermés de $`A`$. On suppose $`U`$ unipotent commutatif dans $`A`$, alors la projection $`A/B\stackrel{p}{}A/BU`$ est un fibré affine dont le fibré direction sur $`A/BU`$ est déduit de la représentation (par automorphismes intérieurs) de $`BU`$ sur $`U/(UB)`$.
Démonstration: On pose $`C=BU`$ et on considère la variété $`A\times (C/B)`$ qui est munie d’une action de $`C`$ donnée de la façon suivante : $`(c,(a,x))(ac^1,cx)`$. On peut alors regarder le quotient $`A\times ^C(C/B)`$ qui est muni d’une projection vers $`A/C`$. La fibration ainsi obtenue : $`A\times ^C(C/B)A/C`$ est exactement la fibration $`p`$ de l’énoncé. En effet, si on regarde l’action restreinte de $`B`$ sur $`A\times (C/B)`$, le quotient $`A\times ^B(C/B)`$ a une projection vers $`A/B`$ qui est scindée (il suffit de prendre la section qui à $`\widehat{a}`$ la classe de $`aA`$ dans $`A/B`$ associe la classe de $`(a,\overline{1})`$ dans $`A\times ^B(C/B)`$, où on a noté $`\overline{1}`$ la classe de l’identité dans $`C/B`$). Mais alors cette section et les morphismes naturels $`A\times ^B(C/B)A\times ^C(C/B)A/C`$ nous donnent la flèche de $`A/B`$ dans $`A/C`$ qui a $`\widehat{a}`$ associe $`\stackrel{~}{a}`$ pour tout $`aA`$. Ainsi, on a une flèche qui respecte les morphismes vers $`A/C`$ de $`A/B`$ dans $`A\times ^C(C/B)`$ et on peut construire sa réciproque : à la classe de $`(a,\overline{c})`$ dans $`A\times ^C(C/B)`$ on associe $`\widehat{ac}`$.
Ainsi $`p`$ est une fibration de groupe structural $`C`$ agissant sur $`C/B=U/(BU)`$. L’action de $`C`$ est l’action naturelle de $`C`$ sur $`C/B`$. Mais si $`bB`$ et $`uU`$ alors la classe de $`buC`$ dans $`U/(BU)`$ est celle de $`bub^1`$ donc l’action de $`C`$ sur $`U/(BU)`$ est donnée par $`(bu,\overline{u^{}})\overline{bub^1.bu^{}b^1}`$. On voit ainsi que cette action est affine : la partie $`bu^{}b^1`$ étant linéaire et on a la translation par $`bub^1`$. La fibration $`p`$ est donc affine et sa direction est donnée par la partie linéaire de l’action c’est à dire par l’action de $`C`$ sur $`U/(BU)`$ donnée par $`(bu,\overline{u^{}})\overline{bu^{}b^1}`$. C’est l’action restreinte de $`BC`$ sur $`U/(BU)`$ par automorphisme intérieur.
Il nous reste à voir que tous les fibrés vectoriels associés à ces fibrés affines sont définis sur $`R/R^{}`$ et qu’ils sont tous engendrés par leurs sections. Il suffit de montrer que dans le cas de la fibration $`P_{i1}^{}/P_{i1}P_i^{}/P_i`$ l’action de $`P_i`$ sur $`Z_i^{}/Z_i`$ est en fait donnée par une action de $`P_{i+1}`$. Mais l’action est donnée par $`(p,\overline{z^{}})\overline{pz^{}p^1}`$. Si on remplace $`p`$ par $`pz`$ avec $`zZ_{i+1}`$ alors on a $`\overline{pzz^{}z^1p^1}=\overline{pz^{}p^1}`$ car $`N`$ est abélien. Si on prend un élément $`\stackrel{~}{p}`$ de $`P_{i+1}`$, on peut définir son action sur $`Z_i^{}/Z_i`$ par l’action de $`p`$. On voit ainsi que tous les fibrés direction sont définis sur la base $`R^{}/R`$. Enfin, comme l’action de $`R`$ sur $`Z_i^{}`$ se prolonge à $`R^{}`$, les fibrés associés à la représentation $`Z_i^{}/Z_i`$ de $`R`$ sont quotient du fibré trivial associé à la représentation de $`R`$ (qui se prolonge à $`R^{}`$) dans $`Z_i^{}`$. Ils sont donc engendrés par leurs sections.
Si $`N^{}`$ est abélien, le morphisme $`P^{}wP/PR^{}/R`$ est un fibré vectoriel. En effet, la tour se réduit à un fibré affine et de plus la section de $`R^{}`$ dans $`P^{}`$ nous définit une section de $`R^{}/R`$ dans $`P^{}wP/P`$ qui nous dit que ce fibré affine est vectoriel.
### 2.2 Etude de la codimension du bord de la $`P^{}`$-orbite maximale
On appelle $`P^{}`$-orbite maximale de $`G/P`$ la $`P^{}`$-orbite $`P^{}wP/P`$ qui est dense dans $`G/P`$. Cette orbite est unique. On cherche à quelle condition cette orbite maximale a un complémentaire de codimension supérieure ou égale à $`2`$. On commence par donner une condition pour que $`P^{}wP/P`$ soit la $`P^{}`$-orbite maximale de $`G/P`$. On fixe un Borel $`B`$ dans $`PP^{}`$.
Lemme 3 : La $`P^{}`$-orbite $`P^{}wP/P`$ est dense dans $`G/P`$ si et seulement si $`𝔭(w(𝔭^{}))`$ contient l’algèbre de Lie d’un Borel, c’est à dire, si et seulement si $`w`$ est dans l’orbite de $`w_0`$ (l’élément de longueur maximale) sous $`W(P)\times W(P^{})`$.
Démonstration: Il suffit de démontrer que cette condition implique que l’orbite est maximale pour l’inclusion car par unicité de cette orbite toutes les orbites ainsi obtenues seront isomorphes. Les $`P^{}`$-orbites sont décrites par $`w(P^{})/(w(P^{})P)`$ pour $`wW`$. Cette $`P^{}`$-orbite est isomorphe à $`P^{}/(P^{}w^1(P))`$. Pour maximiser cette orbite on cherche à minimiser $`P^{}w^1(P)`$. Pour que cette intersection soit minimale il faut et il suffit que $`w^1(𝔭)`$ contienne le moins de racines de $`𝔭^{}`$ possible. C’est le cas si et seulement si $`w(𝔭^{})`$ contient toutes les racines qui ne sont pas dans $`𝔭`$. Donc si $`𝔭`$ et $`w(𝔭^{})`$ contiennent le même Borel l’orbite est maximale pour l’inclusion et réciproquement.
Soit $`B`$ un Borel de $`G`$ et soit $`P`$ un parabolique de $`G`$ contenant $`B`$, on note $`\mathrm{\Sigma }(𝔭,𝔟)`$ les racines simples de $`𝔟`$ correspondant aux sommets du diagramme de Dynkin qui définissent $`𝔭`$. Par exemple $`\mathrm{\Sigma }(𝔟,𝔟)`$ est l’ensemble des sommets du diagramme de Dynkin. On définit une involution $`i`$ du diagramme de Dynkin de la façon suivante : soit $`B`$ un Borel de $`G`$ et $`𝔟`$ son algèbre de Lie, soit $`w_0W`$ le seul élément du groupe de Weyl qui envoie $`𝔟`$ sur $`𝔟`$ (c’est l’élément de longueur maximale), soit $`\alpha `$ une racine simple de $`𝔟`$ (cette racine correspond à un sommet du diagramme de Dynkin), on définit $`i(\alpha )`$ comme étant la racine simple de $`𝔟`$ égale à $`w_0(\alpha )`$. Cette involution correspond à l’involution classique du diagramme de Dynkin de $`A_n`$, $`D_{2n}`$ et $`E_6`$ et à l’identité sur les autres diagrammes.
Proposition 6 : La $`P^{}`$-orbite maximale $`P^{}w_0P/P`$ a un complémentaire de codimension supérieure ou égale à $`2`$ dans $`G/P`$ si et seulement si dans le diagramme de Dynkin, $`\mathrm{\Sigma }(𝔭,𝔟)`$ et $`i(\mathrm{\Sigma }(𝔭^{},𝔟))`$ sont disjoints.
Démonstration: Il suffit de montrer que l’application naturelle $`p`$ de $`\mathrm{Pic}(G/P)`$ dans $`\mathrm{Pic}(P^{}w_0P/P)`$ est injective. Or le noyau de cette application est donné par $`\mathrm{Pic}(G/P)\mathrm{Pic}(G/w_0(P^{}))`$ dans $`\mathrm{Pic}(G/(Pw_0(P^{})))𝔥^{}`$ (car le noyau de l’application $`\mathrm{Pic}(G/(Pw_0(P^{})))\mathrm{Pic}(w_0(P^{})/(Pw_0(P^{})))`$ est $`\mathrm{Pic}(G/w_0(P^{}))`$).
Pour que $`p`$ soit injective, il faut et il suffit que cette intersection soit nulle. Or $`\mathrm{Pic}(G/P)`$ est l’orthogonal dans $`𝔥^{}`$ de l’ensemble $`\alpha (𝔭)`$ des racines $`\alpha 𝔭`$ telles que $`\alpha 𝔭`$. On voit que l’intersection $`\mathrm{Pic}(G/P)\mathrm{Pic}(G/w_0(P^{}))`$ est nulle si et seulement si $`\alpha (𝔭)\alpha (w_0(𝔭^{}))`$ engendre tout $`𝔥`$. Si $`RS`$ est l’ensemble des racines simples de $`𝔟`$, alors $`\mathrm{\Sigma }(𝔭,𝔟)=RS(\alpha (𝔭)RS)`$ et $`\mathrm{\Sigma }(w_0(𝔭^{}),𝔟)=RS(\alpha (w_0(𝔭^{}))RS)`$. On voit alors que $`\alpha (𝔭)\alpha (w_0(𝔭^{}))`$ engendre tout $`𝔥`$ si et seulement si $`(\alpha (𝔭)RS)(\alpha (w_0(𝔭^{}))RS)=RS`$ ce qui est équivalent à $`\mathrm{\Sigma }(𝔭,𝔟)`$ et $`\mathrm{\Sigma }(w_0(𝔭^{}),𝔟)=i(\mathrm{\Sigma }(𝔭^{},𝔟))`$ sont disjoints.
Remarques 4 : (ı) Cette condition nous permet de construire pour tout parabolique $`P`$ qui n’est pas un Borel un parabolique $`P^{}`$ tel que la $`P^{}`$-orbite $`P^{}w_0P/P`$ soit maximale et que son complémentaire soit de codimension au moins $`2`$. En effet, soit $`P`$ un parabolique qui n’est pas un Borel, alors il correspond dans le diagramme de Dynkin à un ensemble $`\mathrm{\Sigma }`$ de sommets qui ne contient pas tous les sommets, il suffit alors de prendre pour $`P^{}`$ un parabolique maximal dont le sommet $`s`$ dans le diagramme de Dynkin n’est pas dans $`i(\mathrm{\Sigma })`$. Par contre cette proposition nous montre que ceci ne sera jamais possible avec les Borels. On propose donc une autre méthode pour résoudre ce cas (voir proposition $`7`$).
(ıı) Soit $`P`$ un parabolique qui n’est pas un Borel et $`P^{}`$ un parabolique tel que $`P`$ et $`P^{}`$ vérifient les hypothèses de la proposition $`6`$. Alors la variété $`R^{}/R`$ obtenue à la proposition $`5`$ est donnée par le produit des variétés homogènes (sous les groupes $`G_i`$ definis par les composantes connexes du diagramme de Dynkin de $`G`$ privé de $`P^{}`$) de paraboliques $`P_i`$ définis par les points du diagramme de Dynkin de $`i(\mathrm{\Sigma }(𝔭))`$. Autrement dit, si on considère sur le diagramme de Dynkin les sommets de $`P`$ après involution et ceux de $`P^{}`$ alors ces diagrammes sont disjoints et les $`P_i`$ sont donnés dans les composantes connexes du diagramme de Dynkin de $`G`$ privé de $`P^{}`$ par les points du diagramme de $`i(\mathrm{\Sigma }(𝔭))`$ quand on a retiré ceux de $`P^{}`$. Par exemple si on considère la variété des droites de $`𝐏^4`$ et que l’on prend pour $`P^{}`$ un parabolique fixant une droite alors $`R^{}/R`$ est $`𝐏^2`$.
Proposition 7 : Soit $`B`$ un Borel et $`P`$ un parabolique contenant $`B`$ dont le diagramme de Dynkin (par rapport à $`B`$) a trois sommets consécutifs et que celui du milieu n’est rattaché qu’à deux sommets dans le diagramme de Dynkin de $`G`$ (respectivement a deux sommets consécutifs dont l’un est au bord du diagramme), soit $`P^{}`$ le parabolique obtenu en enlevant le sommet du milieu (respectivement celui du bord), alors $`G/P`$ est un $`𝐏^1`$-bundle au dessus de $`G/P^{}`$
Démonstration: Soient $`𝔭`$ et $`𝔭^{}`$ les algèbres de Lie de $`P`$ et $`P^{}`$, on voit que $`𝔭`$ est une sous algèbre de $`𝔭^{}`$ et que si $`\alpha `$ est la racine simple de $`B`$ qui correspond au sommet que l’on a retiré alors $`𝔭^{}`$ est l’algèbre de Lie engendrée par $`𝔭`$ et $`\alpha `$. Mais comme le sommet correpondant à $`\alpha `$ forme une composante connexe du diagramme de Dynkin de $`G`$ privé de $`P^{}`$, alors on voit que $`\alpha `$ est orthogonale à toutes les racines simples négatives de $`𝔭`$ ce qui impose $`𝔭^{}=𝔭𝔤_\alpha `$.
Remarque 5 : Soient $`P`$ et $`P^{}`$ comme dans la proposition précédente, si $`C`$ est une courbe tracée sur $`G/P`$ dont la classe est $`xA_1(G/P)`$, alors son degré par rapport à la fibration est donné par $`(x,\alpha )`$$`\alpha `$ est la racine simple (et donc le caractère) qui correspond au sommet que l’on a retiré. En d’autres termes, avec les notations de la proposition précédente, si $`\alpha `$ est le sommet (la racine) que l’on a retiré, alors le fibré tangent relatif $`T_{G/P/G/P^{}}`$ est de première classe de Chern $`\alpha `$ (voir par exemple \[D\]).
Si $`P`$ est un Borel, on peut appliquer cette proposition à tous les sommets du diagramme de Dynkin. Si de plus $`C`$ est une courbe tracée dans $`G/B`$ dont la classe $`xA_1(G/B)`$ est dans le cône positif, alors il existe au moins un sommet du diagramme pour lequel le degré de $`C`$ sera positif. En effet, $`x`$ est dans le cône positif si et seulement si pour toute racine simple $`\alpha `$ on a $`(x,\stackrel{ˇ}{\alpha })0`$ (où $`\stackrel{ˇ}{\alpha }`$ est la coracine de $`\alpha `$). Mais les coracines forment le cône ample et sont donc combinaisons linéaires à coefficients positifs des racines simples. Si pour toute racine simple $`\alpha `$ on a $`(x,\alpha )<0`$ alors $`x`$ ne peut pas être dans le cône positif. Ceci nous permettra donc d’appliquer la proposition $`4`$. De même, si $`x`$ est dans le cône strictement positif, il existe un sommet du diagramme de Dynkin tel que $`(x,\alpha )>0`$.
Application : démonstration de l’irréductibilité du schéma de Hilbert : On raisonne par récurrence sur la longueur du diagramme de Dynkin (nombre de sommets). Le cas de $`SL_2`$ est évident. Soit $`P`$ un parabolique de $`G`$ et soit $`\alpha A_1(G/P)`$ dans le cône positif. Si $`P`$ est un Borel, on a vu à la remarque $`6`$ qu’il existe un sommet du diagramme de Dynkin correspondant à la racine simple $`\alpha ^{}`$ tel que $`(\alpha ,\alpha ^{})0`$. Dans ce cas $`G/B`$ est une fibration en droites projectives au dessus de $`G/P^{}`$ tel que le degré de $`\alpha `$ par rapport à cette fibration est positif. La proposition $`4`$ nous permet donc de nous ramener au cas où $`P`$ n’est pas un Borel.
La remarque $`5`$ nous permet alors de construire un parabolique $`P^{}`$ tel que la $`P^{}`$-orbite $`P^{}w_0P/P`$ est maximale et que son complémentaire est de codimension au moins $`2`$. La proposition $`2`$ nous permet de nous ramener au problème d’irréductibilité du schéma des morphismes pour cette orbite.
Enfin, la proposition $`5`$ nous dit que $`P^{}w_0P/P`$ est une tour de fibrés affines associés à des fibrés vectoriels engendrés par leurs sections au dessus d’un produit de variétés homogènes sous des groupes dont le diagramme de Dynkin est de longueur strictement inférieure à celle du diagramme de Dynkin de $`G`$. On conclue par hypothèse de récurrence en utilisant la proposition $`3`$.
Remarque 6 : On sait de la même façon que tous les schémas $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,U)`$$`U`$ est une $`P^{}`$-orbite de $`G/P`$ sont irréductibles.
### 2.3 Existence de courbes lisses
On va dans ce paragraphe montrer l’existence de courbes lisses sur les variétés homogènes. On se restreindra au cas ou la classe de la courbe est dans le cône strictement positif et on verra (remarque $`8`$) comment se ramener à ce cas si la courbe est dans le cône positif. On appelle courbe nodale une courbe irréductible et génériquement réduite qui est lisse ou qui a pour seules singularités des points doubles ordinaires. Si $`F`$ est un fibré vectoriel sur un schéma $`X`$, on dit que $`F`$ sépare les points si pour tout couple de points $`(P,Q)`$ de $`X`$, on a $`h^0(F_P)>h^0(F_{PQ})`$. Si $`X`$ est une variété et $`\alpha A_1(X)`$, on note $`_{0,\alpha }(X)`$ le schéma de Hilbert des courbes rationnelles nodales de classe $`\alpha `$. On dit que $`X`$ vérifie la condition $`()`$ si la variété d’incidence $`\{(x,C)X\times _{0,\alpha }/x\mathrm{Sing}(C)\}`$ est irréductible.
Lemme 4 : Soit $`X`$ une variété munie d’un fibré affine $``$, de fibré vectoriel direction $`F`$ et défini par $`\eta H^1(X,F)`$. Soit $`C`$ une courbe rationnelle irréductible et génériquement réduite sur $`X`$ et $`f:𝐏^1C`$ une désingularisation de $`C`$. On note $`\overline{\eta }`$ la restriction de $`\eta `$ à $`H^1(C,F|_C)`$. Si on suppose que $`F`$ est engendré par ses sections, alors il existe un relèvement de $`f`$ dans $`\mathrm{Aff}(|_C)`$ le fibré affine associé à $``$ au dessus de $`C`$. Si on suppose de plus que $`C`$ est nodale et que l’une des conditions suivante est vérifiée :
(ı) $`\overline{\eta }0`$ et $`X`$ vérifie $`()`$
(ıı) $`f^{}F|_C`$ sépare les points
alors il existe un relèvement lisse de $`f`$ dans $`\mathrm{Aff}()|_C`$.
Démonstration: Comme $`F`$ est engendré par ses sections, $`H^1f^{}F|_C`$ est nul donc $`f^{}|_C`$ est le fibré vectoriel $`f^{}F|_C`$ qui est engendré par ses sections. Une telle section nous donne un relèvement de $`f`$ dans $`\mathrm{Aff}(|_C)`$. Si $`C`$ est lisse, alors un relèvement quelconque de $`C`$ est lisse.
Si $`\overline{\eta }0`$, le fibré affine n’est pas vectoriel et n’a donc pas de section. Prenons un relèvement $`f^{}`$ de $`f`$ donné par une section de $`f^{}F|_C`$. Ce relèvement est nécessairement non bijectif sur $`C`$ (sinon ce serait une section de $`|_C`$). C’est donc une désingularisation partielle de $`C`$. Ainsi, il existe un relèvement $`f^{}`$ de $`f`$ qui désingularise au moins un point de $`C`$. Mais par monodromie (le groupe de monodromie agit transitivement sur les points singulier grace à la condition $`()`$, cf. \[ACGH\] ou \[Har\]) on sait que pour chaque point double il existe une section de $`f^{}F|_C`$ qui désingularise ce point. En prenant une section générale on obtient une désingularisation en tout point.
Supposons que $`f^{}F|_C`$ sépare les points. On sait que $`H^1f^{}F|_C`$ est nul c’est à dire $`f^{}|_C`$ est le fibré vectoriel $`f^{}F|_C`$. Soit $`P`$ un point double de $`C`$ et $`x`$ et $`y`$ ses antécédents par $`f`$. Il existe une section de $`f^{}F|_C`$ qui sépare $`x`$ et $`y`$. Mais alors cette section nous donne un relèvement $`f^{}`$ de $`f`$ qui désingularise $`P`$. Ainsi, il existe un relèvement désingularisant chaque point double et une section générale de $`f^{}F|_C`$ nous donne un relèvement lisse.
Si $`\overline{\eta }=0`$ et qu’on ne suppose plus que $`f^{}F|_C`$ sépare les points alors il n’existe pas nécessairement de relèvement lisse de $`f`$. C’est le cas si $`F`$ est trivial sur une courbe nodale.
Remarques 7 : (ı) Avec les notations du lemme, supposons que $`C`$ est contenue dans une variété homogène $`X`$ et que sa classe dans $`A_1(X)`$ est dans le cône strictement positif. Supposons de plus que $`c_1(F)`$ est non nul dans le cône ample de $`\mathrm{Pic}(X)`$ alors il existe un relèvement lisse de $`f`$ dans $`\mathrm{Aff}()`$. En effet, il suffit de vérifier que l’une des condition du lemme est vérifiée. Il suffit donc de vérifier que $`f^{}F`$ sépare les points. Mais le degré de $`f^{}F`$ sur $`𝐏^1`$ est strictement positif donc $`f^{}F`$ sépare les points.
(ıı) On aura dans la suite besoin de savoir que $`𝐏^2`$ vérifie la condition $`()`$. Ceci est fait dans \[ACGH\].
Lemme 5 : Soient $`C`$ une courbe rationnelle irréductible et génériquement réduite et $`𝐏_C(E)\stackrel{\phi }{}C`$ une fibration en droites projectives au dessus de $`C`$ d’espace tangent relatif $`T`$. Soit $`f:𝐏^1C`$ une désingularisation de $`C`$. Supposons que $`f^{}E=𝒪_{𝐏^1}𝒪_{𝐏^1}(x)`$ avec $`x0`$. Soit $`d0`$ un entier tel que $`dx[2]`$ et $`dx`$, alors il existe un relèvement $`f^{}`$ de $`f`$ tel que $`f_{}^{}[𝐏^1]T=d`$. Supposons de plus que $`C`$ est nodale et $`d>0`$, alors on peut choisir $`f^{}`$ d’image lisse.
Démonstration: Les relèvements de degré relatif $`d`$ de $`f`$ sont donnés par les quotients isomorphes à $`𝒪_{𝐏^1}(\frac{x+d}{2})`$ de $`f^{}E`$.Un tel quotient existe si et seulement si $`dx[2]`$ et $`\frac{x+d}{2}x`$ c’est à dire $`dx`$ ce qui est le cas. En effet, si les conditions sont vérifiées, il existe un tel quotient. Si $`d<x`$ alors pour avoir un tel quotient il faut que $`\frac{x+d}{2}=0`$ c’est à dire $`d=x=0`$. C’est absurde.
Dans le cas où $`C`$ est nodale , il suffit de vérifier que l’on peut séparer les points. La donnée d’un quotient de $`f^{}E`$ isomorphe à $`𝒪_{𝐏^1}(\frac{x+d}{2})`$ étant équivalente à la donnée d’une section partout non nulle de $`f^{}E(\frac{dx}{2})`$, on est ramené à montrer que pour tout couple de points $`(x,y)`$ de $`𝐏^1`$ il existe une telle section $`s`$ telle que $`s(x)`$ et $`s(y)`$ sont linéairement indépendants. Mais $`f^{}E(\frac{dx}{2})=𝒪_{𝐏^1}(\frac{x+d}{2})𝒪_{𝐏^1}(\frac{dx}{2})`$ donc ceci est possible dès que $`\frac{dx}{2}0`$ et $`\frac{x+d}{2}>0`$, ce qui est vérifié sous nous hypothèses.
Si on ne suppose plus $`d>0`$, il n’existe pas nécessairement de relèvement lisse de $`f`$. En effet, si $`E`$ est trivial et $`C`$ nodale alors $`𝐏_C(E)=C\times 𝐏^1`$ n’a pas de relèvement lisse de $`C`$.
La condition $`dx[2]`$ ne dépend que de la classe $`\alpha `$ du relèvement $`f^{}`$ et de $`c_1(E)`$. En effet, on a $`d=\alpha T`$ et $`x=(\phi _{}\alpha )c_1(E)`$. La condition $`dx`$ est équivalente à $`h^1f^{}E(\frac{dx}{2}1)=0`$. C’est donc une condition ouverte.
Démonstration de l’existence de courbes lisses : On procèdera de la façon suivante : on commence par supposer que pour $`\alpha `$ positive dans $`A_1(G/P)`$, le schéma $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,G/P)`$ est non vide. On ramène le cas d’un Borel à celui d’un parabolique qui n’est pas un Borel. On montre ensuite le résultat par récurrence sur la longueur du diagramme de Dynkin. On initialise la récurrence en montrant les cas des groupes dont le diagramme est de longueur au plus $`3`$. On montrera le cas de $`G_2`$ par une méthode différente. Enfin on montre que $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,G/P)`$ est non vide si $`\alpha `$ est positive.
On suppose connu les résultats suivants : il existe des courbes rationnelles lisses sur $`𝐏^n`$ dès que $`n3`$, sur $`𝐏^1\times 𝐏^n`$ dès que $`n2`$ (lemme $`5`$) et sur $`𝐏^1\times 𝐏^1\times 𝐏^1`$ (lemme $`5`$). Sur $`𝐏^2`$ et $`𝐏^1\times 𝐏^1`$ il existe des courbes nodales. Soit $`\alpha `$ dans le cône strictement positif de $`A_1(G/P)`$.
LE CAS DES BORELS : Si $`P=B`$ est un Borel, la remarque $`5`$ nous permet de trouver un parabolique $`P^{}`$ tel que $`G/BG/P^{}`$ est une fibration en droites projectives de fibré tangent relatif $`T`$ vérifiant $`\alpha T>0`$. Comme on a supposé $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,G/B)`$ non vide, il existe une courbe de $`G/P^{}`$ telle que la fibration vérifie les conditions du lemme $`5`$. Comme ces conditions sont ouvertes, la courbe générale vérifie ces conditions. On est donc ramené à prouver l’exitence de courbes lisses (et même seulement à points doubles ordinaires) sur $`G/P^{}`$. Ceci nous permet notamment de dire que sur $`SL_3/B`$ il existe des courbes lisses dont la classe est quelconque dans le cône strictement positif.
LE CAS GÉNÉRAL : On procède par récurrence sur la longueur du diagramme de Dynkin. On choisit un parabolique $`P^{}`$ tel que la $`P^{}`$-orbite $`P^{}w_0P/P`$ est maximale. La proposition $`5`$ et le lemme $`4`$ nous permettent de construire des courbes lisses sur les ouverts $`P^{}w_0P/P`$ : dès qu’il y a des courbes lisses sur $`R^{}/R`$, il y en a sur la $`P^{}`$-orbite (on prend une section du fibré affine). Par hypothèse de récurrence, les seuls cas qui posent problème sont donc ceux où $`R^{}/R`$ est $`𝐏^1`$, $`𝐏^2`$ ou $`𝐏^1\times 𝐏^1`$. Dans ce cas le rang du groupe de Picard est au plus $`2`$ et $`P`$ est donné par au plus deux points dans le diagramme de Dynkin. On choisit pour $`P^{}`$ le parabolique maximal qui correspond à un point extrême du diagramme de Dynkin. Ainsi dans le cas de $`SL_n`$ on se ramène à $`SL_{n1}`$ ; dans le cas de $`Sp_{2n}`$ on se ramène à $`SL_n`$ ou $`Sp_{2n2}`$ ; dans le cas de $`SO_{2n+1}`$ on se ramène à $`SL_n`$ ou $`SO_{2n1}`$ ; dans le cas de $`SO_{2n}`$ on se ramène à $`SL_n`$ ou $`SO_{2n2}`$ ; dans le cas de $`F_4`$ on se ramène à $`Sp_6`$ ou $`SO_7`$ ; dans le cas de $`E_6`$ on se ramène à $`SL_6`$ ou $`SO_{10}`$ ; dans le cas de $`E_7`$ on se ramène à $`SL_7`$, $`SO_{12}`$ ou $`E_6`$ et dans le cas de $`E_8`$ on se ramène à $`SL_8`$, $`SO_{14}`$ ou $`E_7`$.
Ce choix de $`P^{}`$ n’est pas toujours possible pour $`SL_{n+1}`$, $`Sp_{2n}`$, $`SO_{2n+1}`$ et $`F_4`$ qui n’ont que deux points extrêmes. Ceci ne se produit que si $`P`$ est donné par les deux points extrêmes. Dans ce cas, on choisit l’avant dernier point du diagramme. La variété $`R^{}/R`$ est alors $`𝐏^{n2}\times 𝐏^1`$ pour les trois premiers groupes et $`𝐏^2\times 𝐏^1`$ pour $`F_4`$. Il existe donc des courbes lisses sur $`R^{}/R`$ dès que $`n4`$.
Pour initialiser la récurrence, il nous faut donc montrer les cas de $`SL_4`$, $`Sp_4=SO_5`$, $`Sp_6`$ et $`SO_7`$.
On se ramène, dans le cas de $`SL_4`$, pour les incidences point/droite et droite/plan, à $`SL_3/B`$ sur laquelle il y a des courbes lisses. Il existe une $`P^{}`$-orbite de l’incidence point/plan qui est un fibré vectoriel de degré $`(\mathrm{1,1})`$ au dessus de $`𝐏^1\times 𝐏^1`$. Le lemme $`4`$ (et la remarque $`7`$) nous permet de construire des courbes lisses sur cette $`P^{}`$-orbite. De même il existe une $`P^{}`$-orbite de la Grassmannienne des droites qui est un fibré vectoriel de degré $`1`$ au dessus de $`𝐏^2`$. Le lemme $`4`$ nous permet de conclure.
Pour $`Sp_4`$, les variétés homogènes associées à des paraboliques maximaux sont $`𝐏^3`$ et $`Q_3`$ la quadrique de $`𝐏^4`$. On va traiter le cas de $`Q_3`$ en fin de démonstration et sur $`𝐏^3`$ il y a des courbes lisses. Sur $`Sp_4/B`$ on sait tracer des courbes lisses en se ramenant soit à $`𝐏^3`$ soit à $`Q_3`$.
Pour $`Sp_6`$ : si le groupe de Picard est de rang $`1`$, soit le point de $`P`$ est le dernier du diagramme alors en prenant le premier point on a un morphisme vers $`Q_3`$ qui a des courbes lisses, soit ce n’est pas le cas et en choisissant le dernier point on arrive dans $`𝐏^2`$ avec un fibré vectoriel de degré $`2`$ ou $`3`$. On conclue grace au lemme $`3`$ (et à la remarque $`7`$). Si le groupe de Picard est de rang $`2`$, on a deux cas selon que les points sont aux deux extrémités ou non. Dans le premier cas on arrive, en prenant le point du milieu, dans $`𝐏^1\times 𝐏^1`$ avec un fibré vectoriel de degré $`(\mathrm{1,1})`$ on conclue grace au lemme $`3`$ (et à la remarque $`7`$), dans le second on prend un des deux points extrêmes et on arrive dans $`SL_3/B`$ ou $`Sp_4/B`$.
Pour $`SO_7`$ : si le groupe de Picard est de rang $`1`$, soit le point de $`P`$ est le dernier du diagramme alors en prenant le premier point on a un morphisme vers $`𝐏^3`$, soit ce n’est pas le cas et en choisissant le dernier point on arrive dans $`𝐏^2`$ avec un fibré vectoriel de degré $`1`$ ou $`2`$. On conclue grace au lemme $`3`$ (et à la remarque $`7`$). Si le groupe de Picard est de rang $`2`$, on a deux cas selon que les points sont aux deux extrémités ou non. Dans le premier cas on arrive en prenant le point du milieu dans $`𝐏^1\times 𝐏^1`$ avec un fibré vectoriel de degré $`(\mathrm{1,2})`$ on conclue grace au lemme $`3`$ (et à la remarque $`7`$), dans le second on prend un des deux points extrêmes et on arrive dans $`SL_3/B`$ ou $`Sp_4/B`$.
LE CAS DE $`G_2`$ : Le cas de $`G_2`$ pose un problème et on ne peut montrer l’existence de courbes lisses par cette méthode. On va utiliser une seconde méthode. L’une des variétés homogènes de $`G_2`$ est $`Q_5`$ la quadrique de $`𝐏^6`$ pour laquelle le problème est déjà résolu il faut donc le vérifier pour l’autre variété homogène associée à un parabolique maximal et l’incidence qui est un $`𝐏^1`$-bundle au dessus de chacune des deux.
On a la situation suivante : soient $`P_1`$ et $`P_2`$ les paraboliques maximaux de $`G_2`$ et $`B`$ un Borel les contenant. On a alors l’incidence donnée par les flèches $`f:G_2/BG_2/P_1`$ et $`g:G_2/BG_2/P_2`$. On sait de plus que $`f`$ est une fibration en $`𝐏^1`$. On considère $`G_2/P_2`$ plongé dans son plongement minimal qui est alors une quadrique de dimension $`5`$ (dont on sait qu’elle a des courbes lisses). On regarde une section hyperplane générale $`H`$ (sur laquelle il y a aussi des courbes lisses). On obtient ainsi des morphismes $`f^{}`$ et $`g^{}`$ par restriction. Les conditions suivantes sont réalisées :
-Le morphisme $`f^{}`$ est birationnel de $`g^1((G_2/P_2)H)`$ vers $`G_2/P_1`$ et c’est un isomorphisme en dehors d’un fermé $`E`$ de codimension au moins $`1`$ de $`g^1((G_2/P_2)H)`$. Donc $`f^{}(E)`$ est de codimension au moins $`2`$.
\- Le morphisme $`g`$ permet de réaliser $`g^1((G_2/P_2)H)E`$ comme un fibré vectoriel engendré par ses sections au dessus de $`(G_2/P_2)HZ`$$`Z`$ est un fermé de codimension au moins $`2`$ contenu dans $`g(E)`$.
Ainsi, comme il existe des courbes lisses sur $`(G_2/P_2)H`$ et que $`g^{}`$ est un fibré vectoriel engendré par ses sections, il existe des courbes lisses sur $`g^1((G_2/P_2)H)E`$ et par $`f^{}`$ des courbes lisses sur $`G_2/P_1`$.
LE CAS DE $`Q_3`$ : Il nous reste à montrer qu’il existe des courbes rationnelles lisses de tous les degrés sur $`Q_3`$. On considère cette variété comme les droites isotropes pour une forme symplectique dans un espace de dimension $`4`$ (c’est à dire comme une variété homogène sous $`Sp_4`$). La variété d’incidence avec $`𝐏^3`$ est donné par le fibré projectif associé au faisceau de nulle corellation $`E`$ défini par la forme symplectique : $`0𝒪_{𝐏^3}(1)\mathrm{\Omega }^1(1)E0`$. On fixe un point $`P_0`$ de $`𝐏^3`$. Soit $`H_0`$ l’orthogonal de ce point et $`L_0`$ la droite de $`Q_3`$ formée par les droites isotropes de $`H_0`$. La restriction de $`E`$ à $`H_0`$ est donnée par l’extension non triviale $`0𝒪_{H_0}E_{P_0}0`$. Au dessus de $`H_0P_0`$ le fibré $`E`$ est une extension non triviale de $`𝒪_{H_0P_0}`$ par lui même qui a une unique section $`s`$ donnée par $`P(P,(PP_0))`$. Soit $`Z`$ l’image de cette section.
On s’intéresse maintenant à $`Q_3L_0`$ et la variété d’incidence $`X`$ (qui est au dessus de $`𝐏^3P_0`$). On a une section de $`Q_3L_0`$ vers $`X`$ donnée par $`L(LH_0,L)`$. Soit $`Y`$ l’image de cette section, si on reprojette $`Y`$ vers $`𝐏^3`$, son image est $`H_0P_0`$. Plus précisément $`Y`$ est le fibré $`𝐏_{H_0P_0}(E)`$ privé de la section $`Z`$. $`Y`$ est donc un fibré affine non vectoriel au dessus de $`H_0P_0`$, de fibré vectoriel direction $`𝒪_{H_0p_0}`$, donné par $`\eta H^1𝒪_{H_0p_0}`$. On cherche à tracer des courbes lisses grace au lemme $`4`$. Soit $`C`$ une courbe rationnelle nodale de degré $`d`$ sur $`H_0P_0`$, on cherche à la relever en une courbe lisse dans $`Y`$ qui est isomorphe à $`Q_3L_0`$. Le lemme $`4`$ nous dit (qu’il faut et) qu’il suffit que l’image $`\overline{\eta }`$ de $`\eta `$ dans $`H^1𝒪_C`$ soit non nulle (ou que $`C`$ soit lisse). Ceci est vrai dès que $`d3`$ (si $`d2`$ la courbe $`C`$ est lisse).
EXISTENCE DE COURBES RATIONNELLES PARAMÉTRÉES : Pour terminer la démonstration du théorème $`1`$, il nous reste à prouver la non vacuité de $`\mathrm{𝐇𝐨𝐦}_\alpha (𝐏^1,G/P)`$ pour $`\alpha `$ dans le cône positif. Si $`P`$ n’est pas un Borel, en prenant tous les points du diagramme distincts de ceux de $`P`$ après involution, on construit une $`P^{}`$-orbite maximale qui est une tour de fibrés affines au dessus de $`R^{}/R`$ dont les fibrés vectoriels direction sont engendrés par leurs sections. De plus le choix de $`P^{}`$ nous permet de dire que $`R^{}/R`$ est un produit de variétés de la forme $`G^{}/B^{}`$$`B^{}`$ est un Borel de $`G^{}`$. Le lemme $`4`$ nous permet de nous ramener au cas des Borels.
Il reste donc à traiter ce cas. Pour cela on procède une fois encore par récurrence sur la longueur du diagramme de Dynkin. On utilise la construction suivante en supposant que la longueur du diagramme est au moins $`3`$. Soit $`B`$ un Borel de $`G`$, soit $`\alpha A_1(G/B)`$ dans le cône positif, soient $`x`$ et $`y`$ deux points du diagramme de Dynkin qui ne sont pas sur la même arête, et notons $`P_x`$ (resp. $`P_y`$) et $`P_{x,y}`$ les paraboliques donnés par tous les points du diagramme sauf $`x`$ (resp. sauf $`y`$) et par tous les points du diagramme sauf $`x`$ et $`y`$. On a alors le diagramme suivant :
$$\begin{array}{ccc}G/B& \stackrel{p^{}}{}G/P_x& \\ q^{}& q& \\ G/P_y& \stackrel{p}{}G/P_{x,y}& \end{array}$$
Toutes les flèches de ce diagramme sont des fibrations en droites projectives et si $`p`$ et $`q`$ sont donnés par les fibrés $`E`$ et $`F`$ de rang $`2`$, alors $`p^{}`$ et $`q^{}`$ sont donnés par les fibrés $`q^{}E`$ et $`p^{}F`$.
Sur $`G/P_x`$ (resp. $`G/P_y`$), on sait tracer des courbes de classe $`p_{}^{}\alpha `$ (resp. $`q_{}^{}\alpha `$). En effet, la remarque préliminaire permet de se ramener au cas des Borels d’un groupe dont le diagramme de Dynkin est de longueur strictement inférieure à celle de $`G`$ et on conclue par hypothèse de récurrence. Ceci nous permet d’affirmer que dans $`G/P_{x,y}`$ il existe des courbes vérifiant la condition du lemme $`5`$ pour le fibré $`F`$. Ces courbes forment un ouvert $`U`$ non vide de $`\mathrm{𝐇𝐨𝐦}_{p_{}p_{}^{}\alpha }(𝐏^1,G/P_{x,y})`$.
De la même façon on voit que $`\mathrm{𝐇𝐨𝐦}_{q_{}^{}\alpha }(𝐏^1,G/P_y)`$ est non vide. Son image dans $`\mathrm{𝐇𝐨𝐦}_{q_{}q_{}^{}\alpha }(𝐏^1,G/P_{x,y})`$ ($`q_{}q_{}^{}\alpha =p_{}p_{}^{}\alpha `$) est un ouvert (proposition $`4`$) non vide. Elle rencontre donc l’ouvert $`U`$ précédent. Ainsi, il existe dans $`\mathrm{𝐇𝐨𝐦}_{p_{}\alpha }(𝐏^1,G/P_y)`$ un élément $`f:𝐏^1G/P_y`$ tel que $`f^{}p^{}F=𝒪_{𝐏^1}𝒪_{𝐏^1}(x)`$ avec $`x0`$, $`xd[2]`$ et $`xd`$$`d=p_{}^{}\alpha T_q`$ est le degré de $`p_{}^{}\alpha `$ par rapport à la fibration $`q`$ ($`T_q`$ est le fibré tangent relatif de $`q`$). Mais le degré de $`\alpha `$ par rapport à la fibration $`q^{}`$ est $`\alpha T_q^{}=\alpha p_{}^{}{}_{}{}^{}T_q=p_{}^{}{}_{}{}^{}(p_{}^{}\alpha T_q)=d`$. Ainsi $`f`$ vérifie les conditions du lemme $`5`$ pour $`p^{}F`$ et ceci nous permet de construire un relèvement dans $`G/B`$.
On s’est ainsi ramené aux cas des groupes dont le diagramme est de longueur au plus $`2`$. Pour $`SL_2`$ et $`SO_4=SL2\times SL_2`$ c’est évident. Il nous reste donc à montrer le cas variétés de drapeaux complets des groupes $`SL_3`$, $`SO_5=Sp_4`$ et $`G_2`$. On procède pour les trois de la même façon et on note $`G`$ l’un de ces trois groupes. Soit $`X`$ la variétés homogène qui correspond aux droites (pour le second groupe on le considère comme $`Sp_4`$) et $`Y`$ l’autre variété homogène (qui correspond aux points). Soit $`B`$ un Borel de $`G`$, soit $`\phi `$ la fibration en droites projectives $`G/BX`$ et de fibré tangent relatif $`T_\phi `$ et soit $`\alpha A_1(G/B)`$ de degrés $`d_1`$ et $`d_2`$ par rapport à $`Y`$ et $`X`$, alors $`\alpha T_\phi =2d_1d_2`$. Soit $`f:𝐏^1X`$ un morphisme de degré $`d_2`$. Soit $`E`$ la restriction du fibré tautologique de la Grassmannienne $`𝐆(2,m)`$ ($`m=\mathrm{3,4,7}`$ selon les cas) à $`X`$, c’est un fibré vectoriel associé à la fibration $`\phi `$. On a nécessairement $`f^{}E=𝒪_{𝐏^1}(a)𝒪_{𝐏^1}(b)`$ avec $`0ab`$ et $`a+b=d_2`$. Ainsi $`f`$ vérifie les hypothèses du lemme $`5`$ dès que $`2d_1d_2d_2ba`$ ie $`d_1d_2`$ ce qui dans ce cas nous permet de relever $`f`$ dans $`G/B`$.
Il nous reste donc à tracer des courbes sur $`G/B`$ pour $`d_2d_1`$. Pour ceci on trace une courbe de bidegré $`(d_1,d_2^{})`$ dans $`G/B`$ avec $`d_1d_2^{}`$. Ceci nous permet de dire qu’il existe un faisceau $`F`$ associé à la fibration de $`G/B`$ au dessus de $`Y`$ tel que si $`f:𝐏^1Y`$ est le morphisme induit, alors $`f^{}F`$ vérifie les hypothèse du lemme $`5`$ (le degré relatif $`d`$ est alors $`2d_2^{}d_1`$ pour $`SL_3`$, $`2d_2^{}2d_1`$ pour $`Sp_4`$ et $`2d_2^{}3d_1`$ pour $`G_2`$). On sait ainsi que $`f^{}F=𝒪_{𝐏^1}𝒪_{𝐏^1}(x)`$ avec $`x0`$, $`dx[2]`$ et $`dx`$. Mais alors si on veut une courbe de bidegré $`(d_1,d_2)`$ avec $`d_2d_2^{}`$ quelconque, il suffit de relever $`f`$ dans $`G/B`$ ce qui est possible car le degré relatif est alors plus grand que $`d`$ (et donc plus grand que $`x`$) et que la parité ne change pas.
Remarque 8 : Soit $`C`$ un courbe dont la classe dans $`A_1(G/P)`$ est dans le cône positif mais pas dans le cône strictement positif. On distingue deux types de points parmis les points définissant $`P`$ dans le diagramme de Dynkin, ceux pour lesquels le degré de $`C`$ est strictement positif et ceux pour lesquels le degré de $`C`$ est nul. On note $`P^{}`$ le parabolique obtenu en prenant les points du deuxième type. On voit que $`PP^{}`$ et on a ainsi un morphisme $`G/PG/P^{}`$. La courbe $`C`$ est tracée dans une fibre de ce morphisme. Si on regarde les sommets du diagramme de Dynkin de $`P`$ dans les composantes connexes du diagramme de Dynkin de $`G`$ privé de $`P^{}`$, le produit de variétés homogènes ainsi défini est isomorphe à la fibre. Pour savoir si il existe des courbes lisses, on est ainsi ramené aux cônes strictements positifs de ce produit de variétés homogènes.
UNE DÉSINGULARISATION DES VARIÉTÉS DE SCHUBERT
On va proposer dans cette partie une application des $`P^{}`$-orbites définies dans la partie précédente. En s’inspirant en grande partie de la désingularisation des variétés de Schubert donnée par M. Demazure \[D\], on construit une désingularisation plus fine des variétés de Schubert.
Soit $`B`$ un Borel et $`wW`$, dans son article \[D\], M. Demazure construit une suite de paraboliques $`Q_i`$ minimaux (associés à un unique sommet du diagramme) tels que $`wQ_i\overline{BwB}`$ (ce qui revient à dire $`𝔮_i𝔟+𝔟^{}`$), $`Q_1`$ contient $`B`$, $`Q_iQ_{i+1}`$ contient un Borel, la suite des $`Q_iB`$ est décroissante et $`𝔮_i=𝔟+𝔟^{}`$ (on a noté $`B^{}=w^1Bw`$). La désingularisation est alors donnée par le quotient par $`B`$ de la variété $`X^{}=Y^{}/G^{\prime \prime }`$$`Y^{}=_iQ_i`$ et $`G^{\prime \prime }=_i(Q_iQ_{i+1})`$. Le produit dans $`G`$ donne un morphisme de $`Y^{}`$ dans $`\overline{BwB}`$ qui est invariant sous $`G^{\prime \prime }`$ et par passage au quotient sous $`B`$ on a le morphisme $`D:X^{}\overline{BwB/B}`$.
Lorsque $`Y^{}`$ est le produit d’un ou deux groupes, le morphisme de $`X^{}`$ dans $`G`$ est bijectif sur son image. Lorsque $`Y^{}`$ a plus de facteurs, le morphisme de $`X^{}`$ dans $`\overline{BwB}`$ est seulement birationnel. Plus $`Y^{}`$ a de facteurs, plus le morphisme est succeptible d’avoir des contractions. Ce que l’on va faire ici est réduire, de façon canonique, le nombre de facteurs de $`Y^{}`$. On va ainsi regrouper les paraboliques $`Q_i`$ pour les remplacer par des paraboliques plus grand. On cherche donc des paraboliques $`P_i`$ les plus grand possible tels que $`wP_i\overline{BwB}`$, $`P_1`$ contient $`B`$, $`P_iP_{i+1}`$ contient un Borel, la suite des $`P_iB`$ est décroissante et $`𝔭_i=𝔟+𝔟^{}`$.
Cette consruction nous permettra de donner un critère pour qu’une variété de Schubert soit une $`P^{\prime \prime }`$-orbite pour un parabolique $`P^{\prime \prime }`$ de $`G`$. On donnera ainsi une condition suffisante (non nécessaire) de lissité des variété de Schubert.
On utilise l’abus de notation suivant : si $`\alpha `$ est une racine et $`𝔭`$ l’algèbre de Lie d’un parabolique, on dit que $`\alpha 𝔭`$ si $`\alpha `$ est une valeur propre pour l’action du tore $`𝔥`$ sur $`𝔭`$ ou encore si $`𝔤_\alpha `$ apparait dans la décomposition de $`𝔭`$. En d’autres termes, on identifie l’algèbre de Lie $`𝔭`$ et la partie parabolique de $`P`$ qui est l’ensemble des racines de $`P`$.
## 1 Construction des paraboliques
On commence par montrer (lemme $`1`$) qu’il existe un parabolique $`P_1`$ contenu dans $`w^1\overline{BwB}`$, contenant $`B`$ et qui est maximal pour cette propriété. Ceci nous permet de construire par récurrence une suite de paraboliques plus gros que ceux de M. Demazure qui donnent la désingularisation annoncée.
Lemme 1 : Soient $`𝔟`$ et $`𝔟^{}`$ deux Borels de $`𝔤`$ contenant le tore $`𝔥`$. Il existe un unique parabolique $`𝔭_1`$ et un unique parabolique $`𝔭_1^{}`$ qui sont maximaux pour les propriétés suivantes : $`𝔟𝔭_1𝔟+𝔟^{}`$ et $`𝔟^{}𝔭_1^{}𝔟+𝔟^{}`$.
Démonstration: Il suffit de montrer l’existence et l’unicité de $`𝔭_1`$, par symétrie celle de $`𝔭_1^{}`$ en découlera. On prend alors pour $`𝔭_1`$ le parabolique contenant $`𝔟`$ défini par les opposées des racines simples de $`𝔟`$ qui sont dant $`𝔟^{}`$. Vérifions que ce parabolique est contenu dans $`𝔟+𝔟^{}`$. Soit $`\alpha 𝔭`$. Si $`\alpha 𝔟`$, on a terminé. Sinon, $`\alpha `$ s’écrit $`(\alpha _i)`$$`\alpha _i`$ est une racine simple de $`𝔟`$ telle que $`\alpha _i𝔟^{}`$. Mais alors $`(\alpha _i)𝔟^{}`$ donc $`\alpha 𝔟^{}`$.
Soit maintenant $`𝔭_2`$ un parabolique vérifiant ces conditions. Soit $`\alpha `$ une racine simple de $`𝔟`$ telle que $`a𝔭_2`$. Alors, $`\alpha 𝔟^{}`$ et donc $`\alpha 𝔭_1`$. Ainsi $`𝔭_2𝔭_1`$.
Ce lemme permet de construire une suite de paraboliques qui serviront à la désingularisation. Soient $`B`$ et $`B^{}`$ deux Borels. On note $`P_1`$ et $`P_1^{}`$ les paraboliques construits à partir de $`B`$ et $`B^{}`$ et du lemme $`1`$. On construit ainsi par récurrence deux suites de Borels $`B_n`$ et $`B_n^{}`$ ($`B_1=B`$ et $`B_1^{}=B^{}`$) et deux suites de paraboliques $`P_n`$ et $`P_n^{}`$ tels que $`B_n`$ est contenu dans $`P_{n1}`$ et $`P_n`$ et de même $`B_n^{}`$ est contenu dans $`P_{n1}^{}`$ et $`P_n^{}`$. En effet, supposons $`B_n`$, $`B_n^{}`$, $`P_n`$ et $`P_n^{}`$ construits, alors on construit $`B_{n+1}`$ (et par symétrie $`B_{n+1}^{}`$) de la façon suivante : on décrit les $`\alpha 𝔭_n`$ qui sont dans $`𝔟_{n+1}`$ : si $`\alpha 𝔭_n`$ est telle que $`\alpha 𝔭_n`$ alors $`\alpha 𝔟_{n+1}`$. Si $`\alpha 𝔭_n𝔭_n^{}`$ et $`\alpha 𝔭_n`$ mais $`\alpha 𝔭_n^{}`$ alors $`\alpha 𝔟_{n+1}`$. Si $`\alpha 𝔭_n`$ mais $`\alpha 𝔭_n^{}`$ et $`\alpha 𝔭_n𝔭_n^{}`$ alors $`\alpha 𝔟_{n+1}`$. Enfin, si $`\alpha 𝔭_n𝔭_n^{}`$ et $`\alpha 𝔭_n𝔭_n^{}`$ alors $`\alpha 𝔟_{n+1}\alpha 𝔟_n`$. Une fois les Borels $`B_{n+1}`$ et $`B_{n+1}^{}`$ définis, on définit $`P_{n+1}`$ et $`P_{n+1}^{}`$ comme étant les paraboliques obtenus à partir de $`B_{n+1}`$ et $`B_{n+1}^{}`$ et du lemme $`1`$.
Pour le lemme suivant, on a besoin du :
Fait 1 : Soit $`𝔭`$ un parabolique et soient $`\alpha 𝔭`$ et $`\alpha ^{}𝔭`$ deux racines, on a l’implication :
$$\alpha +\alpha ^{}𝔭\alpha 𝔭$$
Démonstration: On va en fait montrer que la première condition implique la condition suivante (qui est équivalente à celle que l’on cherche) : $`\alpha `$ appartient à tous les Borels de $`𝔭`$.
Soit $`𝔟`$ un Borel de $`𝔭`$, $`\alpha ^{}`$ s’écrit $`(\alpha _j^{})`$$`\alpha _j^{}`$ est une racine simple de $`𝔟`$. Comme $`\alpha ^{}𝔭`$, il existe au moins un $`j`$ pour lequel $`\alpha _j^{}𝔭`$. Supposons que $`\alpha 𝔟`$, alors $`\alpha `$ s’écrit $`(\alpha _i)`$$`\alpha _i`$ est une racine simple de $`𝔟`$. Comme $`\alpha 𝔭`$, on a pour tout $`i`$ : $`\alpha _i𝔭`$. Mais alors $`\alpha +\alpha ^{}𝔭`$ impose que $`\alpha +\alpha ^{}`$ s’écrive $`(\alpha _k^{\prime \prime })`$, où pour tout $`k`$ la racine $`\alpha _k^{\prime \prime }`$ est simple dans $`𝔟`$ et telle que $`\alpha _k^{\prime \prime }𝔭`$. Mais alors les $`\alpha _j^{}`$ sont contenus dans les $`\alpha _k^{\prime \prime }`$ ce qui impose que pour tout $`j`$ on ait $`\alpha _j^{}𝔭`$ ce qui est une contradiction.
Lemme 2 : Les parties $`𝔟_{n+1}`$ et $`𝔟_{n+1}^{}`$ sont les algèbres de Lie de Borels et on a pour tout $`n`$ : $`𝔟_{n+1}+𝔟_{n+1}^{}𝔟_n+𝔟_n^{}`$ et $`𝔟_n𝔟_1^{}𝔟_{n+1}𝔟_1^{}𝔟_{n+1}^{}𝔟_1^{}𝔟_n^{}𝔟_1^{}`$.
Démonstration: On commence par montrer que les parties $`𝔟_{n+1}`$ et $`𝔟_{n+1}^{}`$ sont les algèbres de Lie de Borels. Il suffit par symétrie de le faire pour $`𝔟_{n+1}`$. On doit donc montrer que $`𝔟_{n+1}`$ est stable, que pour tout $`\alpha 𝔤`$, l’une des deux racines $`\alpha `$ ou $`\alpha `$ est dans $`𝔟_{n+1}`$ et que l’on a jamais les deux en même temps. On commence par la stabilité. Soit $`\alpha `$ et $`\alpha ^{}`$ des racines de $`𝔟_{n+1}`$.
$``$ Si $`\alpha `$ et $`\alpha ^{}`$ ne sont pas dans $`𝔭_{n+1}`$, alors c’est aussi le cas de $`(\alpha +\alpha ^{})`$ et donc $`\alpha +\alpha ^{}𝔟_{n+1}`$.
$``$ Si $`\alpha `$ n’est pas dans $`𝔭_n`$ et que l’on a $`\alpha ^{}𝔭_n𝔭_n^{}`$ et $`\alpha ^{}𝔭_n`$ mais $`\alpha ^{}𝔭_n^{}`$, alors si $`(\alpha +\alpha ^{})𝔭_n`$, on a (Fait $`1`$) $`\alpha ^{}𝔭_n`$ ce qui est absurde. Donc $`(\alpha +\alpha ^{})𝔭_n`$ et donc $`\alpha +\alpha ^{}𝔟_{n+1}`$.
$``$ Si $`\alpha `$ n’est pas dans $`𝔭_n`$ et que l’on a $`\alpha ^{}𝔭_n𝔭_n^{}`$ et $`\alpha ^{}𝔭_n𝔭_n^{}`$, alors de la même façon on doit avoir $`(\alpha +\alpha ^{})𝔭_n`$ et donc $`\alpha +\alpha ^{}𝔟_{n+1}`$.
$``$ Si $`\{\alpha ,\alpha ^{}\}𝔭_n𝔭_n^{}`$ et $`\{\alpha ,\alpha ^{}\}𝔭_n`$ mais $`\alpha 𝔭_n^{}`$, $`\alpha ^{}𝔭^{}`$ alors $`\alpha +\alpha ^{}𝔭_n𝔭_n^{}`$, $`(\alpha +\alpha ^{})𝔭_n`$ et $`(\alpha +\alpha ^{})𝔭_n^{}`$ donc $`\alpha +\alpha ^{}𝔟_{n+1}`$.
$``$ Si $`\alpha 𝔭_n𝔭_n^{}`$ et $`\alpha 𝔭_n`$ mais $`\alpha 𝔭_n^{}`$ et que l’on a $`\alpha ^{}𝔭_n𝔭_n^{}`$ et $`\alpha ^{}𝔭_n𝔭_n^{}`$, alors $`\alpha +\alpha ^{}𝔭_n𝔭_n^{}`$ et $`(\alpha +\alpha ^{})𝔭_n`$. Si de plus $`(\alpha +\alpha ^{})𝔭_n^{}`$, on a (Fait $`1`$) $`\alpha ^{}𝔭_n^{}`$ ce qui est absurde donc $`(\alpha +\alpha ^{})𝔭_n^{}`$ et donc $`\alpha +\alpha ^{}𝔟_{n+1}`$.
$``$ Enfin, si $`\{\alpha ,\alpha ^{}\}𝔭_n𝔭_n^{}`$, $`\{\alpha ,\alpha ^{}\}𝔭_n𝔭_n^{}`$ et $`\{\alpha ,\alpha ^{}\}𝔟_n`$, alors $`\alpha +\alpha ^{}𝔭_n𝔭_n^{}`$, $`(\alpha +\alpha ^{})𝔭_n𝔭_n^{}`$ et $`\alpha +\alpha ^{}𝔟_n`$ donc $`\alpha +\alpha ^{}𝔟_{n+1}`$.
Soit maintenant $`\alpha 𝔤`$. Si $`\alpha 𝔭_n`$ alors $`\alpha 𝔭_n`$ et donc $`\alpha 𝔟_{n+1}`$. De même si $`\alpha 𝔭_n`$ alors $`\alpha 𝔭_n`$ et donc $`\alpha 𝔟_{n+1}`$. Il reste donc les racines $`\alpha 𝔭_n`$ telles que $`\alpha 𝔭_n`$. On sait que $`\alpha `$ ou $`\alpha `$ est dans $`𝔭_n^{}`$, on peut donc supposer (quitter à échanger $`\alpha `$ et $`\alpha `$) que $`\alpha 𝔭_n^{}`$. On a alors deux cas : $`\alpha 𝔭_n^{}`$ ou $`\alpha 𝔭_n^{}`$. Dans le premier cas on sait que $`\alpha 𝔟_{n+1}`$, dans le second on a $`\alpha 𝔟_{n+1}\alpha 𝔟_n`$. Or on sait que $`𝔟_n`$ est un Borel donc $`\alpha `$ ou $`\alpha `$ est dans $`𝔟_n`$ et ainsi $`\alpha `$ ou $`\alpha `$ est dans $`𝔟_{n+1}`$. Il reste à voir que l’on a pas $`\alpha `$ et $`\alpha `$ dans $`𝔟_{n+1}`$. Si c’est le cas on sait que $`\alpha `$ et $`\alpha `$ sont dans $`𝔭_n`$. Si $`\alpha 𝔭_n^{}`$ alors $`\alpha 𝔭_n^{}`$ et ceci impose que $`\alpha 𝔟_{n+1}`$ ce qui est absurde. Par symétrie on peut donc supposer que $`\alpha `$ et $`\alpha `$ sont dans $`𝔭_n^{}`$, mais alors $`\alpha 𝔟_{n+1}\alpha 𝔟_n`$ et comme $`𝔟_n`$ est un Borel on ne peut avoir $`\alpha `$ et $`\alpha `$ dans $`𝔟_{n+1}`$.
Par construction on sait que $`𝔟_{n+1}𝔭_n`$ et $`𝔟_{n+1}^{}𝔭_n`$ ce qui nous donne que $`𝔟_{n+1}+𝔟_{n+1}^{}𝔭_n+𝔭_n^{}=𝔟_n+𝔟_n^{}`$.
On procède par récurrence en supposant que $`𝔟_n𝔟_1^{}𝔟_n^{}𝔟_1^{}`$ cette propriété étant évidement vraie pour $`n=1`$. On a $`𝔟_{n+1}^{}𝔟_1^{}(𝔟_n+𝔟_n^{})𝔟_1^{}𝔟_n^{}𝔟_1^{}`$ par hypothèse de récurrence.
De même, on a $`𝔟_{n+1}𝔟_1^{}𝔟_n^{}𝔟_1^{}`$. Soit alors $`\alpha 𝔟_{n+1}𝔟_1^{}`$. On sait alors que $`\alpha 𝔭_n𝔟_n^{}𝔭_n𝔭_n^{}`$ et on a les cas suivants :
$``$ $`\alpha 𝔭_n^{}`$ alors $`\alpha `$ est dans tous les Borels de $`𝔭_n^{}`$ et donc $`\alpha 𝔟_{n+1}^{}`$.
$``$ $`\alpha 𝔭_n^{}`$ mais $`\alpha 𝔭_n`$ alors $`\alpha 𝔟_n^{}`$ (sinon $`\alpha 𝔟_n+𝔟_n^{}`$ alors que $`\alpha 𝔭_n^{}`$) et donc $`\alpha 𝔟_n^{}`$ et ce cas ne se produit pas.
$``$ $`\alpha 𝔭_n𝔭_n^{}`$ alors on a $`\alpha 𝔟_n^{}`$ et donc $`\alpha 𝔟_{n+1}^{}`$.
On conclue ainsi que $`𝔟_{n+1}𝔟_1^{}𝔟_{n+1}^{}𝔟_1^{}`$.
On sait que $`𝔟_n𝔟_1^{}𝔟_n^{}𝔟_1^{}`$. Soit $`\alpha 𝔟_n𝔟_1^{}`$. Si $`\alpha 𝔭_n`$ alors $`\alpha `$ est dans tous les Borels de $`𝔭_n`$ et donc $`\alpha 𝔟_{n+1}`$. Supposons $`\alpha 𝔭_n`$. Comme $`\alpha 𝔟_n^{}`$ alors $`\alpha 𝔭_n^{}`$. Mais alors on a les deux cas suivants :
$``$ $`\alpha 𝔭_n^{}`$ alors $`\alpha 𝔟_n`$ (sinon $`\alpha 𝔟_n+𝔟_n^{}`$ alors que $`\alpha 𝔭_n`$) ce qui impose $`\alpha 𝔟_n`$ et ce cas est exclu.
$``$ $`\alpha 𝔭_n^{}`$ et on a $`\alpha 𝔭_n𝔭_n^{}`$ et $`\alpha 𝔭_n𝔭_n^{}`$ donc comme $`\alpha 𝔟_n`$ on $`\alpha 𝔟_{n+1}`$.
On conclue ainsi que $`𝔟_n𝔟_1^{}𝔟_{n+1}𝔟_1^{}`$.
On a ainsi construit deux suites de Borels et deux suites de paraboliques. On s’arrête dès que $`P_nP_n^{}`$ contient un Borel (il est nécessaire d’avoir cette propriété pour que le morphisme $`\pi `$ que l’on construit dans la suite et qui est la désingularisation soit propre). On a le lemme suivant qui nous permet de dire que notre construction s’arrête.
Lemme 3 : Si $`𝔭_n𝔭_n^{}`$ ne contient pas de Borel alors l’inclusion $`𝔟_{n+1}+𝔟_{n+1}^{}𝔟_n+𝔟_n^{}`$ est stricte.
Démonstration: Supposons que $`𝔭_n𝔭_n^{}`$ ne contient pas de Borel et que $`𝔟_{n+1}+𝔟_{n+1}^{}=𝔟_n+𝔟_n^{}`$. Si il existe $`\alpha 𝔭_n`$ telle que $`\alpha 𝔭_n^{}`$ et que $`\alpha 𝔭_n𝔭_n^{}`$ alors $`\alpha 𝔟_n+𝔟_n^{}`$ et $`\alpha 𝔟_{n+1}𝔟_{n+1}^{}`$ donc $`\alpha 𝔟_{n+1}+𝔟_{n+1}^{}`$ ce qui est impossible. De même si il existe $`\alpha 𝔭_n𝔭_n^{}`$ telle que $`\alpha 𝔭_n^{}`$ mais $`\alpha 𝔭_n`$ alors $`\alpha 𝔟_n+𝔟_n^{}`$ et $`\alpha 𝔟_{n+1}+𝔟_{n+1}^{}`$ ce qui est impossible. Les racines de $`𝔟_n`$ sont donc d’un des trois types suivants : $`\alpha 𝔭_n𝔭_n^{}`$ et $`\alpha 𝔭_n𝔭_n^{}`$ ou $`\alpha 𝔭_n`$, $`\alpha 𝔭_n^{}`$, $`\alpha 𝔭_n^{}`$ et $`\alpha 𝔭_n`$ ou $`\alpha 𝔭_n𝔭_n^{}`$ et $`\alpha 𝔭_n+𝔭_n^{}`$.
Soit maintenant $`𝔟`$ un Borel de $`𝔭_n`$ tel que $`\mathrm{Card}(𝔟𝔭_n𝔭_n^{})`$ est maximal (ici on appelle $`\mathrm{Card}(𝔭)`$ le nombre de racines qui apparaissent dans $`𝔭`$). Alors il existe une racine simple $`\alpha `$ de $`𝔟`$ telle que $`\alpha 𝔭_n𝔭_n^{}`$ et $`\alpha 𝔭_n𝔭_n^{}`$ (pour $`\alpha `$ c’est clair sinon $`𝔟`$ serait contenu dans $`𝔭_n𝔭_n^{}`$, si $`\alpha 𝔭_n𝔭_n^{}`$, alors $`s_\alpha (𝔟)`$ est un Borel de $`𝔭_n`$ tel que $`s_\alpha (𝔟)𝔭_n𝔭_n^{}=(𝔟𝔭_n𝔭_n^{})\{\alpha \}`$ ce qui contredit la maximalité). On sait alors que $`\alpha 𝔭_n`$ mais $`\alpha 𝔭_n^{}`$ et $`\alpha 𝔭_n^{}`$ mais $`\alpha 𝔭_n`$. On montre maintenant que $`\alpha 𝔭_n`$ ce qui sera une contradiction. Pour cela il suffit de montrer que pour tout $`\beta 𝔟_n`$ on a $`\beta \alpha 𝔟_n+𝔟_n^{}`$. Mais la remarque faite au debut nous permet de dire que l’on les trois cas suivants :
$``$ $`\beta 𝔭_n𝔭_n^{}`$ et $`\beta 𝔭_n𝔭_n^{}`$, alors $`\beta 𝔭_n^{}`$ et $`\alpha 𝔭_n^{}`$ donc $`\beta \alpha 𝔭_n^{}𝔟_n+𝔟_n^{}`$.
$``$ $`\beta 𝔭_n𝔭_n^{}`$ et $`\beta 𝔭_n+𝔭_n^{}`$, alors $`\beta 𝔭_n^{}`$ et $`\alpha 𝔭_n^{}`$ donc $`\beta \alpha 𝔭_n^{}𝔟_n+𝔟_n^{}`$.
$``$ $`\beta 𝔭_n`$ mais $`\beta 𝔭_n^{}`$ et $`\beta 𝔭_n^{}`$ mais $`\beta 𝔭_n`$, alors $`\beta `$ est dans tous les Borels de $`𝔭_n`$ et en particulier dans $`𝔟`$. On peut donc écrire $`\beta =\alpha _i`$ où les $`\alpha _i`$ sont des racines simples de $`𝔟`$. Mais alors $`\beta \alpha =\alpha _i\alpha `$ est une racine de $`𝔤`$ si et seulement si $`\alpha \{\alpha _i\}`$ (car $`\alpha `$ est une racine simple) et donc $`\beta \alpha 𝔟𝔭_n𝔟_n+𝔟_n^{}`$.
## 2 La désingularisation
Soient $`P`$ et $`P^{}`$ deux paraboliques contenant un même Borel et soit $`wW`$. On peut maintenant construire la désingularisation de la variété $`\overline{P^{}wP/P}`$ qui est l’adhérence de la $`P^{}`$-orbite $`P^{}wP/P`$ dans $`G/P`$. Le lemme $`1`$ de la partie précédente nous permet de construire un Borel $`BP`$ et $`w^{}\overline{w}`$ tels que $`𝔟+𝔟^{}=𝔭+𝔭^{}`$ ($`𝔟^{}=w^{}(𝔟)`$). On peut donc se contenter de construire la suite de paraboliques pour $`B`$ et $`w^{}`$. La désingularisation de $`\overline{P^{}wP/P}`$ sera alors le quotient de celle de $`\overline{Bw^{}B/B}`$ par $`P`$. Ce sera une fibration en $`P/B`$. On note $`B^{}=w^{}(B)`$. Pour cette construction, on s’inspire directement de celle de Demazure \[D\]. Les lemmes précédents nous ont permis de construire des suites de paraboliques et de Borels. De plus le lemme $`3`$ nous permet de dire qu’a partir d’un certain rang $`𝔭_n𝔭_n^{}`$ contiendra un Borel : tant que ce n’est pas le cas la suite des $`𝔟_n+𝔟_n^{}`$ est strictement décroissante (en dimension) et sera donc constante à partir d’un certain rang.
On note $`P_i`$, $`P_i^{}`$, $`B_i`$ et $`B_i^{}`$ les paraboliques associés aux algèbres de Lies $`𝔭_i`$, $`𝔭_i^{}`$, $`𝔟_i`$ et $`𝔟_i^{}`$. On construit la variété suivante : $`X=P_1^{}\times ^{P_1^{}P_2^{}}P_2^{}\mathrm{}P_2\times ^{P_1P_2}P_1`$ qui est le quotient de $`Y=P_1^{}\times P_2^{}\mathrm{}P_2\times P_1`$ par $`G^{}=(P_1^{}P_2^{})\times \mathrm{}\times (P_1P_2)`$. On a un morphisme de $`Y`$ vers $`G`$ donné par le produit (et par multiplication par $`w^{}`$ pour arriver dans $`\overline{Bw^{}B}`$) qui est invariant sous l’action de $`G^{}`$ ce qui nous donne un morphisme de $`X`$ dans $`G`$. On voit alors que $`B`$ (et même $`P`$) agit à droite sur ces variétés et on obtient ainsi un morphisme $`\pi `$ de $`X/B`$ (ou $`X/P`$) vers $`G/B`$ (ou $`G/P`$).
On va comparer cette construction avec celle donnée par M. Demazure dans \[D\] : on a, pour former $`P_1`$, regroupé les paraboliques que choisit M. Demazure. On voit ainsi que $`\pi `$ factorise la désingularisation de Demazure (cf. proposition $`1`$).
L’image de $`\pi `$ est l’adhérence de la cellule $`BwB/B`$ et $`X/B`$ est le quotient de $`P_1^{}\times P_2^{}\times \mathrm{}\times P_2\times P_1/B`$ par $`G^{}`$ et est donc lisse. On peut aussi le voir en considérant la filtration suivante :
$$X/BP_1^{}\times ^{P_1^{}P_2^{}}P_2^{}\times \mathrm{}\times P_2/(P_1P_2)\mathrm{}P_1^{}\times ^{P_1^{}P_2^{}}P_2^{}/(P_2^{}P_3^{})P_1^{}/(P_1^{}P_2^{})$$
où la première flèche est une fibration en $`P_1/B`$,la seconde en $`P_2/(P_1P_2)`$ et ainsi de suite chacune des flèches est un fibration en $`P_{n+1}/(P_nP_{n+1})`$ ou en $`P_n^{}/(P_n^{}P_{n+1}^{})`$. Ce qui prouve que $`X/P`$ est lisse.
Remarque 1 : On va voir que cette désingularisation est plus fine que celle de \[D\] : la désingularisation de Demazure se factorise par $`\pi `$. Notre désingularisation est alors un isomorphisme sur un ouvert plus grand. En fait notre désingularisation est bijective sur tout l’ouvert $`w^{}P_1^{}B/B`$ qui contient $`w^{}B^{}B/B=Bw^{}B/B`$.
Exemple 1 : Si $`Bw^{}B/B`$ est la cellule maximale de $`G/B`$ alors ceci signifie que $`𝔟+𝔟^{}=𝔤`$ et donc $`𝔭_1=𝔭_1^{}=𝔤`$. Ainsi notre désingularisation est $`G/B`$ (qui était déjà lisse) alors que celle de \[D\] était plus compliquée et notamment pas un isomorphisme.
Pour montrer la factorisation annoncée, on construits des paraboliques minimaux contenus dans les $`P_i`$ et qui redonne la désingularisation de M. Demazure.
Lemme 4 : Soit $`n`$ le plus petit entier tel que $`𝔭_n𝔭_n^{}`$ contienne un Borel. Il existe $`𝔟_{n+1}`$ un Borel de $`𝔭_n𝔭_n^{}`$ tel que $`𝔟_n𝔟_1^{}𝔟_{n+1}𝔟_1^{}𝔟_n^{}𝔟_1^{}`$.
Démonstration: On modifie $`𝔟_n`$ pour le faire appartenir à $`𝔭_n^{}`$ : si $`𝔟_n𝔭_n^{}`$, on pose $`𝔟_{n+1}=𝔟_n`$. Sinon, il existe $`\alpha `$ racine simple de $`𝔟_n`$ telle que $`\alpha 𝔭_n^{}`$ mais $`\alpha 𝔭_n`$. En effet, soit $`\alpha `$ une racine simple de $`𝔟_n`$ qui n’est pas dans $`𝔭_n^{}`$ (qui existe car $`𝔟_n𝔭_n^{}`$), alors si $`\alpha 𝔭_n`$, on a $`\alpha 𝔭_n𝔭_n^{}`$ et $`\alpha 𝔭_n𝔭_n^{}`$ ce qui contredit le fait que $`𝔭_n𝔭_n^{}`$ contient un Borel. On regarde alors $`s_\alpha (𝔟_n)`$. On a $`s_\alpha (𝔟_n)𝔭_n`$ et $`s_\alpha (𝔟_n)𝔭_n^{}=(𝔟_n𝔭_n^{})\{\alpha \}`$. De plus, si $`\alpha 𝔟_1^{}`$, alors $`\alpha 𝔟_n𝔭_1^{}𝔟_n^{}𝔟_1^{}𝔭_n^{}`$ ce qui est absurde. Donc $`\alpha 𝔟_1^{}`$ et $`s_\alpha (𝔟_n)𝔟_1^{}=(𝔟_n𝔟_1^{})\{\alpha \}`$. On continue ce processus tant que le Borel n’est pas contenu dans $`𝔭_n^{}`$. On obtient ainsi le Borel $`𝔟_{n+1}`$ qui vérifie $`𝔟_{n+1}𝔭_n𝔭_n^{}`$, $`𝔟_n𝔟_1^{}𝔟_{n+1}𝔟_1^{}`$. Enfin, $`𝔟_{n+1}𝔟_1^{}(𝔟_n+𝔟_n^{})𝔟_1^{}𝔟_n^{}𝔟_1^{}`$.
Proposition 1 : Soit $`S`$ une variété de Schubert, il existe une désingularisation de Demazure $`D:X^{}S`$ et un morphisme $`X^{}X`$ qui s’insère dans le diagramme suivant :
$$\begin{array}{ccc}X^{}& \stackrel{D}{}& S\\ & & \mathrm{Id}\\ X& \stackrel{\pi }{}& S\end{array}$$
Démonstration: On commence par le lemme suivant :
Lemme 5 : Soit $`𝔭`$ un parabolique, $`𝔟`$ un Borel, $`𝔟^{}`$ et $`𝔟^{\prime \prime }`$ des Borels de $`𝔭`$ tels que $`𝔟𝔟^{}𝔟𝔟^{\prime \prime }`$, alors il existe une suite de Borels $`(𝔟_i)_{1in}`$ de $`𝔭`$ tels que $`𝔟_1=𝔟^{}`$, $`𝔟_n=𝔟^{\prime \prime }`$, $`𝔟_i𝔟𝔟_{i+1}𝔟`$ et $`\mathrm{Card}(𝔟_{i+1}𝔟)=\mathrm{Card}(𝔟_i𝔟)+1`$.
Démonstration: On peut se placer dans le cas où l’inclusion $`𝔟𝔟^{}𝔟𝔟^{\prime \prime }`$ est stricte. Alors il existe $`\alpha `$ une racine simple de $`𝔟^{}`$ telle que $`\alpha 𝔟𝔟^{\prime \prime }`$. En effet, sinon toute les racines simples de $`𝔟^{}`$ sont telles que $`\alpha 𝔟𝔟^{\prime \prime }`$ c’est à dire $`\alpha 𝔟`$ ou $`\alpha 𝔟^{\prime \prime }`$. Mais si $`\alpha 𝔟`$ alors $`\alpha 𝔟𝔟^{}𝔟𝔟^{\prime \prime }`$ et donc dans tous les cas $`\alpha 𝔟^{\prime \prime }`$. Ceci impose que $`𝔟^{}𝔟^{\prime \prime }`$ et donc $`𝔟^{}=𝔟^{\prime \prime }`$ ce qui est impossible en raison de l’inclusion stricte.
On pose alors $`𝔟_2=s_\alpha (𝔟^{})`$ et on a ($`\alpha 𝔟^{\prime \prime }`$ donc $`\alpha 𝔭`$) $`𝔟_2𝔭`$ et $`𝔟_2𝔟=(𝔟^{}𝔟)\{\alpha \}𝔟𝔟^{\prime \prime }`$. On recommence le processus tant que l’inclusion de $`𝔟_i𝔟𝔟𝔟^{\prime \prime }`$ est stricte.
Ce lemme nous permet de construire pour tout $`1kn`$ deux suites de Borels $`(𝔟_{k,i})_{1ir_k}`$ (resp. $`(𝔟_{k,i}^{})_{1ir_k^{}}`$) de $`𝔭_i`$ (resp. $`𝔭_i^{}`$) tels que les $`𝔟_{k,i}𝔟_1^{}`$ (resp. $`𝔟_{k,i}^{}𝔟_1^{}`$) forment une suite croissante (resp. décroissante) pour l’ordre lexicographique et que leur dimension augente (resp. descende) de exactement $`1`$ à chaque. De même le lemme $`5`$ nous permet de construire deux suites de Borels $`(𝔟_{n+1,i})_{1ir_{n+1}}`$ de $`𝔭_n`$ (entre $`𝔟_n`$ et $`𝔟_{n+1}`$) et $`(𝔟_{n+1,i}^{})_{1ir_{n+1}^{}}`$ de $`𝔭_n^{}`$ (entre $`𝔟_{n+1}`$ et $`𝔟_n^{}`$) telles que les dimensions de leurs intersections avec $`𝔟_1^{}`$ forment une suite croissante pour la première et décroissante pour la seconde. Ces deux suites complètent les deux premières suites en une seule qui est telle que les intersections avec $`𝔟_1^{}`$ forment une suite dont les dimensions prennent une seule fois toutes les valeurs entre $`\mathrm{dim}(𝔟_1^{})`$ et $`\mathrm{dim}(𝔟_1𝔟_1^{})`$. Cette suite de Borels nous permet de construire une désingularisation de Demazure correspondant à $`𝔟_1`$ et $`𝔟_1^{}`$ en prenant, si $`𝔟_x`$ et $`𝔟_y`$ sont deux termes consécutifs de la suite le parabolique $`𝔟_x+𝔟_y`$ qui est minimal mais différent d’un Borel. On construit ainsi pour tout $`1kn+1`$ et tout $`1ir_k1`$ les paraboliques $`𝔭_{k,i}=𝔟_{k,i}+𝔟_{k,i+1}`$, pour tout $`1kn`$ les paraboliques $`𝔭_{k,r_k}=𝔟_{k,r_k}+𝔟_{k+\mathrm{1,1}}`$ et le parabolique $`𝔭_{n+1,r_{n+1}}=𝔟_{n+1,r_{n+1}}+𝔟_{n+1,r_{n+1}^{}}^{}`$ et de même pour tout $`1kn+1`$ et tout $`1ir_k^{}1`$ les paraboliques $`𝔭_{k,i}^{}=𝔟_{k,i}^{}+𝔟_{k,i+1}^{}`$ et pour tout $`1kn`$ les paraboliques $`𝔭_{k,r_k^{}}^{}=𝔟_{k,r_k^{}}^{}+𝔟_{k+\mathrm{1,1}}^{}`$. On pose alors $`Y^{}=P_{k,i}^{}\times P_{k,i}`$ le premier produit est effectué dans l’ordre lexicographique et le second dans l’ordre lexicographique inversé, on pose $`G^{\prime \prime }=B_{k,i}^{}\times B_{k,i}`$ (avec les mêmes ordres) et on pose $`X^{}=Y^{}/G^{\prime \prime }`$. La désingularisation de Demazure de $`\overline{B_1^{}B_1/B_1}`$ est alors donnée par $`X^{}/B_1`$ et un morphisme $`D`$ de cette variété vers $`G/B_1`$. Le morphisme $`D`$ est obtenu à partir de la multiplication de $`Y^{}`$ dans $`G`$. Or cette multiplication se factorise par $`Y`$ car les groupes $`P_{k,i}`$ sont contenus dans $`P_k`$ et par passage au quotient on voit que $`D`$ se factorise par $`\pi `$.
Corollaire 1 : Le morphisme $`\pi `$ est une désingularisation
Démonstration: On commence par montrer le cas des variétés de Schubert : $`\overline{Bw^{}B/B}`$ est exactement $`\overline{w^{}P_1^{}B/B}`$. On obtient de cette façon un ouvert (la $`P_1^{}`$-orbite $`w^{}P_1^{}B/B`$) lisse plus grand que la cellule de Schubert. Le morphisme $`\pi `$ est un isomorphisme au dessus de cet ouvert. En effet, cet ouvert contient la cellule $`Bw^{}B/B`$ au dessus de laquelle la désingularisation de Demazure est un isomorphisme donc $`\pi `$ est aussi un isomorphisme au dessus de cette cellule. De plus, le morphisme $`\pi `$ est invariant sous l’action de $`P_1`$ donc l’orbite de la cellule lisse précédente (c’est exactement $`w^{}P_1^{}B/B`$) est une $`P_1^{}`$-orbite lisse et $`\pi `$ est encore un isomorphisme au dessus de cette orbite.
Pour le cas général on considère le diagramme commutatif suivant :
$$\begin{array}{ccc}X/B& & \overline{Bw^{}B/B}\\ & & \\ X/P& & \overline{w^{}P_1^{}P/P}\end{array}$$
dont les flèches verticales sont des fibrations en $`P/B`$ (pour la seconde ceci vient du fait que comme $`𝔭_1+𝔭_1^{}=𝔟+𝔟^{}=𝔭+𝔭^{}`$, la variété $`\overline{Bw^{}B/B}`$ est l’image réciproque de $`\overline{w^{}P_1^{}B/B}`$ par le morphisme de $`G/BG/P_1`$) ainsi comme le morphisme $`\pi `$ pour les Borels est une désingularisation et est donc birationnel, alors le morphisme $`\pi `$ de $`X/P`$ vers la variété de Schubert correspondante est birationnel et c’est bien une désingularisation.
Remarque 2 : La $`P_1^{}`$-orbite $`w^{}P_1^{}B/B`$ est le plus grand ouvert lisse de $`\overline{Bw^{}B/B}`$ qui est une orbite sous l’action d’un sous groupe de $`G`$ laissant stable $`\overline{Bw^{}B/B}`$ (ceci vient du fait que $`𝔭_1^{}`$ a été choisit maximal pour sa propriété). Par ailleurs, pour tout groupe $`P`$ contenant $`B`$ et contenu dans $`P_1`$, le morphisme $`\overline{Bw^{}P/P}\overline{Bw^{}P_1/P_1}`$ est une fibration en $`P_1/P`$. Ainsi, les singularités des variétés $`\overline{Bw^{}P/P}`$ pour de tels $`P`$ sont identiques à celle de $`\overline{Bw^{}P_1/P_1}`$. Les variétés de Schubert du type de $`\overline{Bw^{}P_1/P_1}`$ seront dits minimales : il n’existe pas de parabolique $`P`$ contenant $`P_1`$ tel que la flèche $`\overline{Bw^{}P_1/P_1}\overline{Bw^{}P/P}`$ est une fibration en $`P/P_1`$.
Notre construction nous permet de donner une condition nécessaire et suffisante pour qu’une variété de Schubert minimale soit homogène sous l’action d’un sous groupe de $`G`$ et on en déduit une condition suffisante de lissité des variété de Schubert. Notons $`S`$ la variété de Schubert $`\overline{Bw^{}P/P}`$ et $`S_1`$ la variété de Schubert minimale associée $`\overline{Bw^{}P_1/P_1}`$ on a alors le :
Corollaire 2 : La variété de Schubert minimale $`S_1`$ est homogène sous l’action d’un sous groupe de $`G`$ si et seulement si $`𝔭_1𝔭_1^{}`$ contient un Borel. Dans ce cas $`S`$ est lisse.
Démonstration: Si $`𝔭_1𝔭_1^{}`$ contient un Borel on a $`X=P_1^{}\times ^{P_1^{}P_1}P_1`$. Si on quotiente par $`P_1`$ qui contient $`P`$ on a alors $`P_1^{}/(P_1^{}P_1)=X/P_1S_1`$ est la désingularisation de $`S_1`$ qui était donc homogène. Mais alors le morphisme $`G/PG/P_1`$ nous donne le diagramme commutatif suivant :
$$\begin{array}{ccc}X/P& & S\\ & & \\ X/P_1& & S_1\end{array}$$
dont les flèches verticales sont des fibrations en $`P_1/P`$ et donc, comme la flèche du bas est un isomorphisme, celle du haut est également un isomorphisme. Ainsi on voit que $`S`$ est lisse car $`X/P`$ l’est.
Réciproquement, si $`S_1`$ est homogène sous l’action d’un sous groupe de $`G`$ alors le plus grand de ces sous groupes est $`P_1^{}`$ par construction. Mais alors $`P_1^{}/(P_1^{}P_1)`$ doit être $`S_1`$ toute entière. Ceci signifie que la cellule $`P_1^{}/(P_1^{}P_1)`$ est propre. Ceci n’est possible que si $`P_1^{}P_1`$ contient un Borel. Remarquons que pour que $`S`$ soit homogène sous $`P_1^{}`$, il faut et il suffit que $`\mathrm{}P_1^{}`$ contienne un Borel.
Exemple 2 : On étudie notre désingularisation pour certaines variétés de Schubert de $`SL_4`$. Soit $`P_0L_0H_0`$ un drapeau de $`𝐏^3`$. Considérons les variétés de Schubert suivantes : $`S=\{(P,L,H)𝐏^3\times 𝐆\times \stackrel{ˇ}{𝐏}^3/PLH\mathrm{et}\mathrm{dim}(LL_0)1\}`$, $`S^{}=\{\{(P,L,H)𝐏^3\times 𝐆\times \stackrel{ˇ}{𝐏}^3/PLH,PL_0\mathrm{et}LH_0\}`$ et $`S^{\prime \prime }=\{\{(P,L,H)𝐏^3\times 𝐆\times \stackrel{ˇ}{𝐏}^3/PLH,PH_0\mathrm{et}P_0H\}`$, cette dernière est la seule variété de Schubert de $`SL_4`$ qui ne correspond pas à une permutation vexillaire. Notre désingularisation est alors donnée dans chacun des cas par $`X=\{(P,P^{},L,H,H^{})𝐏^3\times 𝐏^3\times 𝐆\times \stackrel{ˇ}{𝐏}^3\times \stackrel{ˇ}{𝐏}^3/PLH,P^{}LH^{}\mathrm{et}\mathrm{}^{}L_0H^{}\}`$, $`X^{}=S^{}`$ et $`X^{\prime \prime }=\{\{(P,L,L^{},H)𝐏^3\times 𝐆\times 𝐆\times \stackrel{ˇ}{𝐏}^3/PLH,PL^{}H\mathrm{et}P_0L^{}H_0\}`$. Ces désingularisations sont toutes bijectives sur le lieu lisse. Notons que l’exemple de $`S^{}`$ montre que la condition suffisante de lissité du corollaire $`2`$ n’est pas nécessaire.
Remarque 3 : Si le lieu lisse d’une variété de Schubert $`S`$ est homogène sous l’action d’un sous groupe de $`G`$ alors ce sous groupe est $`P_1`$ et notre désingularisation est bijective sur le lieu lisse. Par ailleurs, on peut vérifier que si le groupe $`G`$ est $`SL_n`$ pour $`n4`$, $`SO_n`$ pour $`n6`$ ou $`Sp_n`$ pour $`n4`$ alors notre désingularisation est bijective sur le lieu lisse.
On peut encore affiner notre désingularisation de la façon suivante : on remplace les paraboliques $`P_i`$ par un parabolique contenant $`P_i`$ et contenu dans $`w^1\overline{BwB}`$. Ceci ne change pas les paraboliques $`P_1`$ et $`P_1^{}`$. Cette technique permet ainsi de construire les désingularisations $`X_1=\{(P,P^{},L,H)𝐏^3\times 𝐏^3\times 𝐆\times \stackrel{ˇ}{𝐏}^3/PLH,P^{}L\mathrm{et}P^{}L_0\}`$ et $`X_2=\{(P,L,H,H^{})𝐏^3\times 𝐆\times \stackrel{ˇ}{𝐏}^3\times \stackrel{ˇ}{𝐏}^3/PLH,LH^{}\mathrm{et}L_0H^{}\}`$ de la variété de Schubert $`S`$ de l’exemple $`2`$. On voit que cette construction n’est plus canonique, il faut choisir un parabolique contenant $`P_i`$ et contenu dans $`w^1\overline{BwB}`$.
Remerciements : Je tiens ici à remercier mon directeur de thèse Laurent Gruson pour toute l’aide qu’il m’a apportée durant la préparation de ce travail et également Patrick Polo pour m’avoir signaler une erreur dans une première version de cet article.
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\[D\] Demazure M. : Désingularisation des variétés de Schubert généralisées, Ann. Sci. ENS (4) 7 (1974).
\[DG\] Demazure M. Gabriel P : Groupes algébriques, Tome I : Géométrie algébrique, généralités, groupes commutatifs. Masson & Cie, Editeur, Paris; North-Holland Publishing Co., Amsterdam (1970).
\[FH\] Fulton W. Harris J. : Representation theory, GTM 129 Springer Verlag, New-York (1991).
\[Har\] Harris J. : The genus of space curves, Math. Ann. 249 (1980) no. 3.
\[K1\] Kempf G.R. : Vanishing theorems for flag maniflods, Amer. j. math. 98 (1976).
\[K2\] Kempf G.R. : Linear systems, Ann. of Math. 103 (1976).
\[Kl\] Kleiman S.L. : The transversality of a general translate, Compositio Math. 28 (1974).
\[Ko\] Kollár J. : Rational curves on algebraic varieties, Ergebnisse der Mathematik und ihrer Grenzgebiete 32, Springer Verlag, Berlin (1996).
\[KP\] Kim B., Pandharipande R. : The connectedness of the moduli space of maps to homogeneous spaces, preprint AG 0003168.
\[LMS\] Lakshmibai V., Musili C., Seshadri C.S. : Cohomology of line bundles on $`G/B`$, Ann. Sci. ENS (4) 7 (1974).
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# References
Rearrangement of the Fermi Surface of Dense Neutron Matter
and Direct Urca Cooling of Neutron Stars
D. N. Voskresensky<sup>1</sup>
Moscow Institute for Physics and Engineering, Moscow 115409, Russia
V. A. Khodel<sup>2</sup> and M. V. Zverev
Russian Research Center Kurchatov Institute, Moscow 123182, Russia
and
J. W. Clark
McDonnell Center for the Space Sciences and Department of Physics, Washington University,
St. Louis, MO 63130 USA
<sup>1</sup>Gesellschaft für Schwerionenforschung GSI, D-64220 Darmstadt, Germany
<sup>2</sup>McDonnell Center for the Space Sciences and Department of Physics, Washington University,
St. Louis, MO 63130 USA
## Abstract
It is proposed that a rearrangement of single-particle degrees of freedom may occur in a portion of the quantum fluid interior of a neutron star. Such a rearrangement is associated with the pronounced softening of the spin-isospin collective mode which, under increasing density, leads to pion condensation. Arguments and estimates based on fundamental relations of many-body theory show that one realization of this phenomenon could produce very rapid cooling of the star via a direct nucleon Urca process displaying a $`T^5`$ dependence on temperature.
1. INTRODUCTION
The EINSTEIN, EXOSAT, and ROSAT orbiting X-ray observatories have measured surface temperatures of certain neutron stars and set upper limits on surface temperatures of others (). The data for the supernova remnants in 3C58, the Crab, and RCW103 indicate relatively slow cooling, while that for Vela, PSR2334+61, PSR0656+14, and Geminga point to substantially more rapid cooling. In the so-called standard scenario for neutron-star cooling, the primary role is played by the modified Urca process $`(nnnpe^{}\overline{\nu }_e;npe^{}nn\nu _e)`$, first considered by Bahcall & Wolf (1965) and later reexamined by Friman & Maxwell (1979) in terms of an in-vacuum one-pion exchange model. Cooling simulations based on the results of these works play out the slow scenario of thermal evolution and fail to explain the rapid cooling of some stars. This picture is profoundly altered when in-medium effects are taken into account. Modification of $`NN`$, $`\pi N`$, and $`KN`$ interactions with increasing density may be so strong that pion () and kaon (\[Brown 1994\]) condensates form in the interior region of a high-mass neutron star, leading to a dramatic increase of the neutrino luminosity ().
Here we shall demonstrate that softening of the spin-isospin (pion) collective mode in dense neutron matter as predicted by Migdal (1978) could give rise to a rearrangement of single-particle degrees of freedom prior to the onset of pion condensation and open a new channel of neutrino cooling of neutron stars, by giving access to the direct Urca process. This process does not involve a neutron (or other) spectator and, if allowed, is an extremely efficient cooling mechanism ().
2. BUBBLE REARRANGEMENT OF THE FERMI SPHERE
A rearrangement of single-particle degrees of freedom takes place if the necessary condition for stability of the normal state of a Fermi liquid is violated. At $`T=0`$ this condition requires that the change of the ground-state energy $`E_0`$ remain positive for any admissible variation $`\delta n(p)`$ of the Landau quasiparticle distribution $`n(p)`$ away from the normal-state step-function distribution $`\theta (pp_F)`$. Formally,
$$\delta E_0=\xi (p,n(p))\delta n(p)\frac{\mathrm{d}^3p}{(2\pi )^3}>0,$$
(1)
where $`\xi (p,n(p))\epsilon (p,n(p))\mu `$ is the energy of a quasiparticle relative to the chemical potential $`\mu `$. The condition (1) fails if a depression with $`\xi <0`$ forms in the spectrum $`\xi (p)`$ at $`p>p_F`$; it likewise fails if there arises an elevation with $`\xi >0`$ at $`p<p_F`$. The rearrangement is precipitated when the density $`\rho `$ reaches a critical value $`\rho _{cF}`$ at which there emerges a bifurcation and a new root $`p=p_0`$ of the relation
$$\xi (p,n(p);\rho _{cF})=0,$$
(2)
which ordinarily serves merely to specify the Fermi momentum $`p_F`$.
The simplest kind of rearrangement of the momentum distribution $`n(p)`$ of quasiparticles of given spin projection retains the property that its values are restricted to 0 and 1, but the Fermi sea becomes doubly connected (). At densities exceeding the critical value $`\rho _{cF}`$, the normal-state distribution $`\theta (pp_F)`$ is modified by the presence of a “bubble,” or absence of particles, over the range $`p_i<p<p_f<p_F`$, with the inner surface $`p_i`$ located relatively close to the origin and the Fermi momentum $`p_F`$ readjusted to maintain the prescribed neutron density. The distance $`p_fp_i`$ between the two new Fermi surfaces lying interior to $`p_F`$ can be estimated using the formula $`\xi (pp_0,n(p);\rho )=\xi _0(\rho \rho _{cF})A(pp_0)^2`$, where $`A`$ and $`\xi _0`$ are positive constants. This formula embodies the essential properties that $`\xi (p)`$ is negative for any $`p<p_F`$ at $`\rho <\rho _{cF}`$ and that its maximum value first reaches zero at $`\rho =\rho _{cF}`$. Employing this parametrization in the relation (2), it is found that no bifurcation point exists for $`\rho <\rho _{cF}`$, whereas two solutions arise for $`\rho >\rho _{cF}`$, with the distance between the two new Fermi surfaces growing in proportion to $`\sqrt{\rho \rho _{cF}}`$.
Were such a rearrangement to occur in the neutron subsystem of neutron-star matter, the emergence of new Fermi surfaces situated at lower momenta would permit the direct nucleon Urca process to operate at a much lower density than hitherto considered possible (). In this process, the tandem reactions $`npe^{}\overline{\nu }_e`$ and $`pe^{}n\nu _e`$ are driven by thermal excitations. The condition of high degeneracy prevailing even in young neutron stars implies that these excitations remain close to the Fermi surfaces of the participants. For the direct Urca mechanism to contribute appreciably to neutron-star cooling, momentum conservation then demands satisfaction of the triangle inequalities $`|p_{Fp}p_{Fe}|p_{Fn}p_{Fp}+p_{Fe}`$ among the proton, electron, and neutron Fermi momenta. With the conventional singly-connected neutron Fermi sphere, these conditions are met only at very high baryon densities where the proton Fermi momentum $`k_{Fp}`$ reaches sufficiently large values. The best current estimates yield a threshold baryon density of at least five times the saturation density $`\rho _0`$ of symmetrical nuclear matter (). On the other hand, if the neutron subsystem undergoes a rearrangement that allows for thermal excitations of neutron quasiparticles at much lower momenta than $`p_{Fn}`$, specifically at $`p_i`$ and $`p_f`$ in the bubble rearrangement scenario, the triangle inequalities are much more easily satisfied and the direct Urca process is greatly facilitated.
We next provide a quantitative basis for this qualitative idea, by appealing to established methods of microscopic many-body theory. For a bifurcation point to arise in the solution of equation (2), both the spectrum $`\xi (p)`$ and the scalar component $`f(𝐩_1,𝐩_2;𝐤=0)`$ of the Landau amplitude of the quasiparticle interaction must depend strongly on momentum. The connection (, p. 37)
$$\frac{\xi _(p)}{𝐩}=\frac{\xi ^0(p)}{𝐩}+\frac{1}{2}\mathrm{Tr}_\sigma \mathrm{Tr}_{\sigma _1}f(𝐩,\sigma ,𝐩_1;\sigma _1,𝐤=0)\frac{n(p_1)}{𝐩_1}\frac{\mathrm{d}^3p_1}{(2\pi )^3}$$
(3)
between these two quantities, wherein $`\xi _0(p)`$ is the free spectrum and spin dependence is made explicit, is usually employed to find the effective mass $`M^{}`$ in terms of the first harmonic of the expansion of the amplitude $`f(\widehat{𝐩}_1\widehat{𝐩}_2;0)`$ in Legendre polynomials. We exploit this relation in the search for new roots of equation (2).
Within the integral in (3), the quasiparticle interaction amplitude $`f(𝐩_1,𝐩_2;𝐤=0)`$ may be replaced by $`(M^{}/M)\mathrm{\Gamma }^k(𝐩_1,𝐩_2)`$, the scattering amplitude $`\mathrm{\Gamma }^k(𝐩_1,𝐩_2)\mathrm{\Gamma }(𝐩_1,𝐩_2;𝐤=0)`$ modified by an effective-mass factor (). The requisite strong momentum dependences can arise if the system approaches a second-order phase transition that occurs at some critical density $`\rho _c`$ where a collective mode of frequency $`\omega _s(k)`$ collapses at wave number $`k=k_0`$ and the corresponding susceptibility diverges along with the scattering amplitude. Near such a soft-mode critical point, the singular part $`\mathrm{\Gamma }^s`$ of the scattering amplitude $`\mathrm{\Gamma }(𝐩_1,𝐩_2;𝐤)`$ at momentum transfer $`𝐤`$ may be expressed quite generally by ()
$$\frac{M^{}}{M}\mathrm{\Gamma }_{\alpha \kappa ;\beta \lambda }^s(𝐩_1,𝐩_2;𝐤,\rho \rho _c)=O_{\alpha \kappa }O_{\beta \lambda }D(k)+O_{\alpha \lambda }O_{\beta \kappa }D(|𝐩_1𝐩_2+𝐤|),$$
(4)
where $`O`$ denotes the vertex determining the structure of the collective-mode operator (e.g. $`O=(\sigma 𝐧)\tau `$ for the spin-isospin mode with pion quantum numbers, where $`𝐧`$ is a unit vector along the relevant momentum). In deriving (4), antisymmetry of the two-particle wave function under interchange of the coordinates and spins of the two particles has been invoked. The propagator $`D`$ may be parametrized as $`D(k)=\left[\beta ^2+\gamma ^2(k^2/k_0^21)^2\right]^1`$, where $`\beta (\rho )`$ measures the proximity to the phase transition point, with $`\beta (\rho \rho _c)\rho _c\rho `$ (cf. Dyugaev (1976)).
The essential messages of the preceding development are that the singular part $`\mathrm{\Gamma }^s(𝐩_1,𝐩_2;𝐤=0)D(|𝐩_1𝐩_2|)`$ of the scattering amplitude depends on the difference $`𝐩_1𝐩_2`$ and that as one approaches the soft-mode phase transition point this dependence becomes quite strong. We assume that the remaining contributions to $`\mathrm{\Gamma }(𝐩_1,𝐩_2;𝐤=0)`$ can be adequately incorporated by renormalization of the chemical potential $`\mu `$. Equation (3) is then easily integrated to produce an explicit expression suitable for calculation of the single-particle spectrum,
$$\xi (p)=\xi ^0(p)+\frac{1}{2}\mathrm{Tr}(O_{\alpha \lambda }O_{\lambda \alpha })D(|𝐩𝐩_1|)n(p_1)\frac{\mathrm{d}^3p_1}{(2\pi )^3}.$$
(5)
We now apply this equation to dense neutron matter in the vicinity of the second-order phase transition associated with neutral pion condensation, which is engendered by the softening of the spin-isospin mode having $`\pi ^0`$ quantum numbers (). It has been predicted that the collapse of this mode will take place at a neutron density $`\rho _c=\rho _{c\pi }`$ in the range (0.2 – 0.5) fm<sup>-3</sup> (roughly 1–3 times $`\rho _0`$), depending on theoretical assumptions (). Unfortunately, there is as yet no definitive microscopic treatment of neutron-star matter from which one can extract or derive quantitatively reliable values for the input parameters $`\beta `$, $`\gamma `$, and $`k_0`$ of our model.
In this situation, a reasonable strategy is to perform calculations based on expression (5) for several choices of the parameters of the microscopic model. Substituting (5) into relation (2), one finds the critical density $`\rho _{cF}`$ for the onset of a bifurcation of the latter equation. For $`\rho >\rho _{cF}`$ this equation then determines two new momenta $`p_i`$ and $`p_f`$ where $`\xi (p)`$ vanishes, which delimit the bubble region of $`n(p)`$ and between which $`\xi (p)`$ is positive. Representative numerical results for the spectrum $`\xi (p)`$ are plotted in Fig. 1. Results for the phase diagram of dense neutron matter in the $`\rho /\rho _0`$ versus $`\beta ^2/m_\pi ^2`$ plane are displayed in Fig. 2. Different values of $`\gamma `$ are considered, while keeping the parameter $`k_0`$ fixed at the value $`0.9p_{Fn}`$ suggested by earlier numerical calculations ().
It is evident that variation of the parameters $`\beta `$, $`\gamma `$, and $`k_0`$ within sensible bounds can have strong effects on the phase diagram and therefore on the extent of the phase with rearranged quasiparticle occupation. Nevertheless, our numerical study has documented three salient features of the bubble rearrangement. First, the critical density $`\rho _{cF}`$ for the rearrangement is less than the critical density $`\rho _{c\pi }`$ for pion condensation. Since both phenomena stem from the strong momentum dependence of the Landau amplitude $`f(𝐩_1,𝐩_2;𝐤0)`$, rearrangement of the quasiparticle distribution may be regarded as a precursor of pion condensation. Second, the bifurcation point corresponding to formation of a hole bubble in the neutron momentum distribution is positioned at small momenta, $`p_0<0.2p_F`$, irrespective of the applicable value of $`\rho _{c\pi }`$. Third, the spectrum $`\xi (p)`$ shows a deep depression for $`p(0.50.6)p_F`$. And fourth, the ratios $`\rho _{cF}/\rho _{c\pi }`$ and $`p_0/p_F`$ are insensitive to the actual value taken by $`\rho _{c\pi }`$ within the usual range of theoretical predictions.
Analogous considerations apply to the proton subsystem of neutron-star matter, in which case one is dealing with the charged pion mode. Estimates () indicate that this mode is also softened in dense matter, with a critical density not far from that for neutral pion condensation. One may then argue that, under the influence of the strongly momentum-dependent external field provided by the neutron medium, protons will leave the old Fermi sphere and occupy states of relatively large momentum, $`p0.5p_{Fn}`$. The impact of this further rearrangement on the proton-neutron ratio and on the rate of neutrino cooling requires a separate analysis, which we defer.
Having laid the microscopic basis for a rearrangement of the neutron Fermi surface that creates a bubble at low momenta in the Fermi sea, we return to its most striking astrophysical implication. Beyond the bifurcation point, the triangle inequalities can now be satisfied without the conventional requirement ()) that the proton fraction exceed some 11–14%. Accordingly, the direct Urca process becomes active in the density regime just short of the threshold for pion condensation. At temperatures $`T`$ above the anticipated superfluid phase transition, the $`T`$-dependence of the resulting neutrino emissivity is determined through the usual expression
$`ϵ(npe^{}\overline{\nu })`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}}}2G_F^2(1+3g_A^2){\displaystyle \underset{i}{}}\epsilon _\nu n_n(1n_p)(1n_e)`$ (6)
$`\times \delta (𝐩_n𝐩_p𝐩_e𝐩_\nu )\delta (\epsilon _n\epsilon _p\epsilon _e\epsilon _\nu ),`$
derived in quasiparticle approximation in terms of the occupations $`n_i`$ and energies $`\epsilon _i`$ of the reacting particles. We do not consider a renormalization of the coupling constants $`G_F^2`$ and $`g_A^2`$ due to medium effects, and the emissivity takes the customary form ()
$$ϵ(npe^{}\overline{\nu })1.2\times 10^{27}\frac{M_n^{}}{M_n}\frac{M_p^{}}{M_p}\left(\frac{\mu _e}{100\text{MeV}}\right)T_9^6\text{erg}\text{cm}^3\text{sec}^1,$$
(7)
where $`\mu _e`$ is the electron chemical potential and the temperature $`T`$ is measured in multiples of $`10^9`$ K. Since the neutron and proton effective masses $`M_n^{}`$ and $`M_p^{}`$ remain $`T`$-independent, the bubble rearrangement serves only to turn on the direct Urca process at a lower density, without altering the familiar $`T^6`$ power-law behavior of the emissivity.
3. FERMION CONDENSATION
Interestingly enough, there exists a more radical scenario for rearrangement of the quasiparticle distribution, known as fermion condensation (). In this case, the occupancy $`n(p)`$ may be partial, i.e., it may lie between 0 and 1. At $`T=0`$, the new quasiparticle distribution $`n(p)`$ is to be found from the variational condition
$$\frac{\delta E_0[n(p)]}{\delta n(p)}=\mu ,p\mathrm{\Omega }.$$
(8)
The left-hand side of equation (8) is just the quasiparticle energy $`\epsilon (p)`$. Hence its coincidence with the chemical potential $`\mu `$ means that the spectrum of single-particle excitations will be dispersionless (i.e., the quasiparticle group velocity will vanish) throughout the entire momentum domain $`\mathrm{\Omega }`$ where $`\xi (p)=\epsilon (p)\mu =0`$, and not merely at isolated points as implied by equation (2). The family of quasiparticles having momenta $`p\mathrm{\Omega }`$ is called the fermion condensate because of a conspicuous analogy with the low-temperature Bose gas, in which the energy of condensate particles is also equal to the chemical potential $`\mu `$. A key signature of fermion condensation has been observed in strongly correlated electron systems (): flat portions of the single-particle spectrum have been seen experimentally in a number of high-temperature superconductors ().
The two types of rearrangement – bubble formation and fermion condensation – can compete with each other in the density regime just below the soft-mode phase transition. Numerical studies () demonstrate that fermion condensation wins the contest at nonzero $`T`$. Let us suppose this is the case in the neutron-star medium, while continuing to disregard nucleonic pairing phenomena. Consistency of the Fermi-Dirac form $`n(p,T)=\left\{\mathrm{exp}\left[\xi (p;n(p,T))/T\right]+1\right\}^1`$ for the quasiparticle distribution with the variational condition (8) requires that the spectrum $`\xi (p,T)`$ of the fermion-condensate phase grows linearly with $`T`$ at low temperature, implying an effective mass inversely proportional to $`T`$ ().
Although the rearranged momentum distribution derived from equation (8) differs from the bubble configuration, its structure will also admit thermal excitations at low neutron momenta. Hence we may again expect most neutron stars to contain a region of relatively moderate density, bounded below by $`\rho _{cF}`$ and above by $`\rho _{c\pi }`$, in which the direct Urca process operates vigorously. However, due to the new feature of a $`T`$-dependent neutron effective mass, $`M_n1/T`$, we may anticipate an enhancement of the neutrino emissivity relative to the standard result (), corresponding to a $`T^5`$ rather than a $`T^6`$ dependence on the temperature.
4. CONCLUSIONS
We have explored the possibility that an effective interaction with strong momentum dependence gives rise to a rearrangement of the neutron momentum distribution in neutron-star matter. Two plausible manifestations of this phase transformation – creation of a doubly-connected Fermi surface and fermion condensation – have been considered. Both open the prospect that direct nucleon Urca cooling is present in a density regime just below the threshold for pion condensation and consequently at a density much lower than previously estimated. If a fermion condensate is formed, the resulting neutrino emissivity is significantly larger than that generated by the direct Urca process in normal matter. Within the affected density range, it would therefore dominate all other proposed neutrino cooling mechanisms (). Future studies along this line will focus on temperatures below the superfluid transition and on the effect of the dramatically increased emissivity on neutrino opacity.
This research was supported by NSF Grants PHY-9602127 and PHY-9900713 (JWC and VAK) and by the McDonnell Center for the Space Sciences (VAK). We thank M. Baldo, M. Di Toro, and E. E. Kolomeitsev for fruitful discussions. DNV expresses his appreciation for hospitality and support provided by GSI Darmstadt. MVZ acknowledges the hospitality of INFN (Sezione di Catania).
FIGURE CAPTIONS
Fig. 1. The dimensionless neutron spectrum $`y_n(p)=\xi _n(p)/(p_F^2/2M)`$ at the critical densities $`\rho _{cF}`$ corresponding to three different sets of model parameters: (a) $`\gamma =1.25m_\pi `$, $`k_0=0.9p_{Fn}`$, $`\beta ^2=0.22m_\pi ^2`$ ($`\rho _{cF}1.19\rho _0)`$, (b) $`\gamma =1.25m_\pi `$, $`k_0=0.9p_{Fn}`$, $`\beta ^2=0.25m_\pi ^2`$ ($`\rho _{cF}1.76\rho _0)`$, (c) $`\gamma =1.25m_\pi `$, $`k_0=p_{Fn}`$, $`\beta ^2=0.13m_\pi ^2`$ ($`\rho _{cF}1.88\rho _0)`$. Two different positions of the bifurcation point, namely $`p_0=0`$ (for parameter sets (a) and (b)) and $`p_00.12p_{Fn}`$ (for set (c)), are indicated by arrows.
Fig. 2. Phase diagram of neutron matter in the variables $`\rho `$ (measured in $`\rho _0`$) and $`\beta ^2`$ (measured in $`m_\pi ^2`$), as calculated for $`k_0=0.9p_{Fn}`$ and four different values of $`\gamma `$, which (in $`m_\pi `$ units) label the corresponding the phase boundaries separating the bubble phase (upper left) from the normal phase (lower right).
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# 1 Introduction
## 1 Introduction
Accretion disks are inferred to exist in objects of very different scales: km to millions of km in low Mass X-ray Binaries (LMXB) and Cataclysmic Variables (CV), solar radius-to-AU scale in protostellar disks, and AU-to-parsec scales for the disks in Active Galactic Nuclei (AGN).
An interesting observational connection exists between accretion disks and jets (such as the spectacular jets from AGN and protostars), and outflows (the ‘CO-outflows’ from protostars and possibly the ‘broad-line-regions’ in AGN). Lacking direct (i.e. spatially resolved) observations of disks, theory has tried to provide models, with varying degrees of success. Uncertainty still exists with respect to some basic questions. In this situation, progress made by observations or modeling of a particular class of objects is likely to have direct impact for the understanding of other objects, including the enigmatic connection with jets.
In this lecture I concentrate on the more basic aspects of accretion disks, but an attempt is made to mention topics of current interest, such as magnetic viscosity, as well. Emphasis is on those aspects of accretion disk theory that connect to the observations of LMXB and CV’s. For other reviews on the basics of accretion disks, see Pringle (1981), Treves et al. (1988). For a more in-depth treatment, see the textbook by Frank et al. (1992).
## 2 Accretion: general
Gas falling into a point mass potential
$$\mathrm{\Phi }=\frac{GM}{r}$$
from a distance $`r_0`$ to a distance $`r`$ converts gravitational into kinetic energy, by an amount $`\mathrm{\Delta }\mathrm{\Phi }=GM(1/r1/r_0)`$. For simplicity, assuming that the starting distance is large, $`\mathrm{\Delta }\mathrm{\Phi }=GM/r`$. If the gas is then brought to rest, for example at the surface of a star, the amount of energy $`e`$ dissipated per unit mass is
$$e=\frac{GM}{r}(\mathrm{r}est)$$
or, if it goes into a circular Kelper orbit at distance $`r`$:
$$e=\frac{1}{2}\frac{GM}{r}(\mathrm{o}rbit).$$
The dissipated energy goes into internal energy of the gas, and into radiation which escapes to infinity (usually in the form of photons, but neutrino losses can also play a role).
### 2.1 Adiabatic accretion
Consider first the case when radiation losses are neglected. This is adiabatic accretion. For an ideal gas with constant ratio of specific heats $`\gamma `$, the internal energy per unit mass is
$$e=\frac{P}{(\gamma 1)\rho }.$$
With the equation of state
$$P=\rho T/\mu $$
(1)
where $``$ is the gas constant, $`\mu `$ the mean atomic weight per particle, we find the temperature of the gas after the dissipation has taken place (assuming that the gas goes into a circular orbit):
$$T=\frac{1}{2}(\gamma 1)T_{\mathrm{vir}},$$
(2)
where $`T_{\mathrm{vir}}`$, the virial temperature is given by
$$T_{\mathrm{vir}}=\frac{GM\mu }{r}.$$
In an atmosphere with temperature near $`T_{\mathrm{vir}}`$, the sound speed is close to the escape speed from the system, the hydrostatic pressure scale height is of the order of $`r`$, and such an atmosphere may evaporate on a relatively short time scale in the form of a stellar wind.
A simple example is spherical adiabatic accretion (Bondi, 1952). An important result is that such accretion is possible only if $`\gamma 5/3`$. The larger $`\gamma `$, the larger the temperature in the accreted gas (eq. 2), and beyond a critical value the temperature is too high for the gas to stay bound in the potential. A classical situation where adiabatic and roughly spherical accretion takes place is a supernova implosion: when the central temperature becomes high enough for the radiation field to start desintegrating nuclei, $`\gamma `$ drops and the envelope collapses onto the forming neutron star via a nearly static accretion shock. Another case are Thorne-Zytkow objects (e.g. Cannon et al. 1992), where $`\gamma `$ can drop to low values due to pair creation, initiating an adiabatic accretion onto the black hole.
Adiabatic spherical accretion is fast, taking place on the dynamical or free fall time scale
$$\tau _\mathrm{d}=r/v_\mathrm{K}=(r^3/GM)^{1/2},$$
(3)
where $`v_\mathrm{K}`$ is the Kepler orbital velocity.
When radiative loss becomes important, the accreting gas can stay cool irrespective of the value of $`\gamma `$, and Bondi’s critical value $`\gamma =5/3`$ plays no role. With such losses, the temperatures of accretion disks are usually much lower than the virial temperature. The optical depth of the accreting flow increases with the accretion rate $`\dot{M}`$. When the optical depth becomes large enough so that the photons are ‘trapped’ in the flow, the accretion just carries them in, together with the gas (Rees 1978, Begelman 1979). Above a certain critical rate $`\dot{M}_\mathrm{c}`$, accretion is therefore adiabatic.
### 2.2 The Eddington Limit
Objects of high luminosity have a tendency to blow their atmospheres away due to the radiative force exerted when the outward traveling photons are scattered or absorbed. Consider a volume of gas on which a flux of photons is incident from one side. Per gram of matter, the gas presents a scattering (or absorbing) surface area of $`\kappa `$ cm<sup>2</sup>. The force exerted by the radiative flux $`F`$ on one gram is $`F\kappa /c`$. The force of gravity pulling back on this one gram of mass is $`GM/r^2`$. The critical flux at which the two forces balance is
$$F_\mathrm{E}=\frac{c}{\kappa }\frac{GM}{r^2}$$
(4)
Assuming that the flux is spherically symmetric, this can be converted into a critical luminosity
$$L_\mathrm{E}=4\pi GMc/\kappa ,$$
(5)
the Eddington luminosity (e.g. Rybicki and Lightman, 1979). If the gas is fully ionized, its opacity is dominated by electron scattering, and for solar composition $`\kappa `$ is then of the order $`0.3`$ cm<sup>2</sup>/g (about a factor 2 lower for fully ionized helium, a factor up to $`10^3`$ higher for partially ionized gases). With these assumptions,
$$L_\mathrm{E}\mathrm{1.7\hspace{0.17em}10}^{38}\frac{M}{M_{}}\mathrm{e}rg/s\mathrm{4\hspace{0.17em}10}^4\frac{M}{M_{}}L_{}$$
If this luminosity results from accretion, it corresponds to the Eddington accretion rate $`\dot{M}_\mathrm{E}`$:
$$\frac{GM}{r}\dot{M}_\mathrm{E}=L_\mathrm{E}\dot{M}_\mathrm{E}=4\pi rc/\kappa .$$
(6)
Whereas $`L_\mathrm{E}`$ is a true limit that can not be exceeded by a static radiating object except by geometrical factors of order unity (see chapter 10 in Frank et al, 1992), no maximum exists on the accretion rate. For $`\dot{M}>\dot{M}_\mathrm{E}`$ the plasma is just swallowed whole, including the radiation energy in it (cf. discussion in the preceding section). With $`\kappa =0.3`$:
$$\dot{M}_\mathrm{E}\mathrm{1.3\hspace{0.17em}10}^{18}r_6\mathrm{g}/s\mathrm{2\hspace{0.17em}10}^8r_6M_{}\mathrm{y}r^1,$$
where $`r_6`$ is the radius of the accreting object in units of 10 km.
## 3 Accretion with Angular Momentum
When the accreting gas has a zonzero angular momentum with respect to the accreting object, it can not accrete directly. A new time scale appears, the time scale for outward transport of angular momentum. Since this is in general much longer than the dynamical time scale, much of what was said about spherical accretion needs modification for accretion with angular momentum.
Consider the accretion in a close binary consisting of a compact (white dwarf, neutron star or black hole) primary of mass $`M_1`$ and a main sequence companion of mass $`M_2`$. The mass ratio is defined as $`q=M_2/M_1`$ (note: $`q`$ is just as often defined the other way around).
If $`M_1`$ and $`M_2`$ orbit each other in a circular orbit and their separation is $`a`$, the orbital frequency $`\mathrm{\Omega }`$ is
$$\mathrm{\Omega }^2=G(M_1+M_2)/a^3.$$
The accretion process is most easily described in a coordinate frame that corotates with this orbit, and with its origin in the center of mass. Matter that is stationary in this frame experiences an effective potential, the Roche potential (Ch. 4 in Frank, King and Raine, 1992), given by
$$\varphi _\mathrm{R}(𝐫)=\frac{GM}{r_1}\frac{GM}{r_2}\frac{1}{2}\mathrm{\Omega }^2r^2$$
(7)
where $`r_{1,2}`$ are the distances of point $`𝐫`$ to stars 1,2. Matter that does not corotate experiences a very different force (due to the Coriolis force). The Roche potential is therefore useful only in a rather limited sense. For non-corotating gas intuition based on the Roche geometry is usually confusing. Keeping in mind this limitation, consider the equipotential surfaces of (7). The surfaces of stars $`M_{1,2}`$, assumed to corotate with the orbit, are equipotential surfaces of (7). Near the centers of mass (at low values of $`\varphi _\mathrm{R}`$) they are unaffected by the other star, at higher $`\mathrm{\Phi }`$ they are distorted and at a critical value $`\mathrm{\Phi }_1`$ the two parts of the surface touch. This is the critical Roche surface $`S_1`$ whose two parts are called the Roche lobes. Binaries lose angular momentum through gravitational radiation and a magnetic wind from the secondary (if it has a convective envelope). Through this loss the separation between the components decreases and both Roche lobes decrease in size. Mass transfer starts when $`M_2`$ fills its Roche lobe, and continues as long as the angular momentum loss from the system lasts. A stream of gas then flows through the point of contact of the two parts of $`S_1`$, the inner Lagrange point $`L_1`$. If the force acting on it were derivable entirely from (7) the gas would just fall in radially onto $`M_1`$. As soon as it moves however, it does not corotate any more and its orbit under the influence of the Coriolis force is different (Fig. 1).
Since the gas at $`L_1`$ is very cold compared with the virial temperature, its sound speed is small compared with the velocity it gets after only a small distance from $`L_1`$. The flow into the Roche lobe of $`M_1`$ is therefore highly supersonic. Such hypersonic flow is essentially ballistic, that is, the stream flows along the path taken by free particles.
Though the gas stream on the whole follows an orbit close to that of a free particle, a strong shock develops at the point where the orbit intersects itself. \[In practice shocks already develop shortly after passing the pericenter at $`M_1`$, when the gas is decelerated again. Supersonic flows that are decelerated by whatever means in general develop shocks (e.g. Courant and Friedrichs 1948, Massey, 1968). The effect can be seen in action in the movie published in Różyczka and Spruit, 1993\]. After this, the gas settles into a ring, into which the stream continues to feed mass. If the mass ratio $`q`$ is not too small this ring forms fairly close to $`M_1`$. An approximate value for its radius is found by noting that near $`M_1`$ the tidal force due to the secondary is small, so that the angular momentum of the gas with respect to $`M_1`$ is approximately conserved. If the gas continues to conserve angular momentum while dissipating energy, it settles into the minimum energy orbit with the specific angular momentum $`j`$ of the incoming stream. The radius of this orbit, the circularization radius $`r_\mathrm{c}`$ is determined from
$$(GM_1r_\mathrm{c})^{1/2}=j.$$
The value of $`j`$ is found by a simple integration of the orbit starting at $`L_1`$ and measuring $`j`$ at some point near pericenter. In units of the orbital separation $`a`$, $`r_\mathrm{c}`$ and the distance $`r_{\mathrm{L1}}`$ from $`M_1`$ to $`L_1`$ are functions of the mass ratio only. As an example for $`q=0.2`$, $`r_{\mathrm{L1}}0.66a`$ and the circularization radius $`r_\mathrm{c}0.16a`$. In practice the ring forms somewhat outside $`r_\mathrm{c}`$, because there is some angular momentum redistribution in the shocks that form at the impact of the stream on the ring.
The evolution of the ring depends critically on nature and strength of the angular momentum transport processes. If sufficient ‘viscosity’ is present it spreads inward and outward to form a disk.
At the point of impact of the stream on the disk the energy dissipated is a significant fraction of the orbital kinetic energy, hence the gas heats up to a significant fraction of the virial temperature. For a typical system with $`M_1=1M_{}`$, $`M_2=0.2M_{}`$ having an orbital period of 2 hrs, the observed size of the disk (e.g. Wood et al. 1989b, Rutten et al. 1992) $`r_\mathrm{d}/a0.3`$, the orbital velocity at $`r_\mathrm{d}`$ about 900 km/s, the virial temperature at $`r_\mathrm{d}`$ is $`10^8`$K. The actual temperatures at the impact point are much lower, due to rapid cooling of the shocked gas. Nevertheless the impact gives rise to a prominent ‘hot spot’ in many systems, and an overall heating of the outermost part of the disk.
## 4 Thin disks: properties
### 4.1 Flow in a cool disk is supersonic
Ignoring viscosity, the equation of motion in the potential of a point mass is
$$\frac{𝐯}{t}+𝐯𝐯=\frac{1}{\rho }P\frac{GM}{r^2}\widehat{𝐫},$$
(8)
where $`\widehat{𝐫}`$ is a unit vector in the spherical radial direction $`r`$. To compare the order of magnitude of the terms, choose a position $`r_0`$ in the disk, and choose as typical time and velocity scales the orbital time scale $`\mathrm{\Omega }_0^1=(r_0^3/GM)^{1/2}`$ and velocity $`\mathrm{\Omega }_0r_0`$. The pressure gradient term is
$$\frac{1}{\rho }P=\frac{}{\mu }T\mathrm{ln}P.$$
In terms of the dimensionless quantities
$$\stackrel{~}{r}=r/r_0,\stackrel{~}{v}=v/(\mathrm{\Omega }_0r_0),$$
$$\stackrel{~}{t}=\mathrm{\Omega }_0t,\stackrel{~}{}=r_0,$$
the equation of motion is then
$$\frac{\stackrel{~}{𝐯}}{\stackrel{~}{t}}+\stackrel{~}{𝐯}\stackrel{~}{}\stackrel{~}{𝐯}=\frac{T}{T_{\mathrm{vir}}}\stackrel{~}{}\mathrm{ln}P\frac{1}{\stackrel{~}{r}^2}\widehat{𝐫}.$$
(9)
All terms and quantities in this equation are of order unity by the assumptions made, except the pressure gradient term which has the coefficient $`T/T_{\mathrm{vir}}`$. If cooling is important, so that $`T/T_{\mathrm{vir}}1`$, the pressure term is negligible to first approximation, and vice versa. Equivalent statements are also that the gas moves hypersonically on nearly Keplerian orbits, and that the disk is thin, as is shown next.
### 4.2 Disk thickness
The thickness of the disk is found by considering its equilibrium in the direction perpendicular to the disk plane. In an axisymmetric disk, using cylindrical coordinates ($`\varpi ,\varphi ,z`$), measure the forces at a point $`𝐫_0`$ ($`\varpi ,\varphi ,0`$) in the midplane, in a frame rotating with the Kepler rate $`\mathrm{\Omega }_0`$ at that point. The gravitational acceleration $`GM/r^2\widehat{𝐫}`$ balances the centrifugal acceleration $`\mathrm{\Omega }_0^2`$$`\varpi `$$`\varpi `$$`\varpi `$at this point, but not at some distance $`z`$ above it because gravity and centrifugal acceleration work in different directions. Expanding both accelerations near $`𝐫_0`$, one finds a residual acceleration toward the midplane of magnitude
$$g_z=\mathrm{\Omega }_0^2z.$$
Assuming again an isothermal gas, the condition for equilibrium in the $`z`$ direction under this acceleration yields a hydrostatic density distribution
$$\rho =\rho _0(\varpi )\mathrm{exp}\left(\frac{z^2}{2H^2}\right).$$
$`H`$, the scale height of the disk, is given in terms of the isothermal sound speed $`c_\mathrm{s}=(T/\mu )^{1/2}`$ by
$$H=c_\mathrm{s}/\mathrm{\Omega }_0.$$
We define $`\delta H/r`$, the aspect ratio of the disk, and find that it can be expressed in several equivalent ways:
$$\delta =\frac{H}{r}=\frac{c_\mathrm{s}}{\mathrm{\Omega }r}=M^1=\left(\frac{T}{T_{\mathrm{vir}}}\right)^{1/2},$$
where $`M`$ is the Mach number of the orbital motion.
### 4.3 Viscous spreading
The shear flow between neighboring Kepler orbits in the disk causes friction due to viscosity. The frictional torque is equivalent to exchange of angular momentum between these orbits. But since the orbits are close to Keplerian, a change in angular momentum of a ring of gas also means it must change its disctance from the central mass. If the angular momentum is increased, the ring moves to a larger radius. In a thin disk angular momentum transport (more precisely a nonzero divergence of the angular momentum flux) therefore automatically implies redistribution of mass in the disk.
A simple example (Lüst 1952, see also Lynden-Bell and Pringle 1974) is a narrow ring of gas at some distance $`r_0`$. If at $`t=0`$ this ring is released to evolve under the viscous torques, one finds that it first spreads into an asymmetric hump with a long tail to large distances. As $`t\mathrm{}`$ the hump flattens in such a way that almost all the mass of the ring is accreted onto the center, while a vanishingly small fraction of the gas carries almost all the angular momentum to infinity. As a result of this asymmetric behavior essentially all the mass of a disk can accrete, even if there is no external torque to remove the angular momentum.
### 4.4 Observations of disk viscosity
Evidence for the strength of the angular momentum transport processes in disks comes from observations of variability time scales. This evidence is not good enough to determine whether the processes really have the same effect as a viscosity, but if this is assumed, estimates can be made of the magnitude of the viscosity.
Cataclysmic Variables give the most detailed information. These are binaries with white dwarf (WD) primaries and (usually) main sequence companions (for reviews see Meyer-Hofmeister and Ritter 1993, Cordova 1995, Warner 1995). A subclass of these systems, the Dwarf Novae, show semiregular outbursts. In the currently most developed theory, these outbursts are due to an instability in the disk (Smak 1971, Meyer and Meyer-Hofmeister 1981, for recent references see, King 1995, Hameury et al. 1998). The outbursts are episodes of enhanced mass transfer of the disk onto the primary, involving a significant part of the whole disk. The decay time of the burst is thus a measure of the viscous time scale of the disk (the quantitative details depend on the model, see Cannizzo et al. 1988, Hameury et al. 1998):
$$t_{\mathrm{visc}}=r_\mathrm{d}^2/\nu ,$$
where $`r_\mathrm{d}`$ is the size of the disk. With decay times on the order of days, this yields viscosities of the order $`10^{15}`$ cm<sup>2</sup>/s, about 14 orders of magnitude above the microscopic viscosity of the gas.
Other evidence comes from the inferred time scale on which disks around protostars disappear, which is of the order of $`10^7`$ years (Strom et al, 1993).
### 4.5 $`\alpha `$-parametrization
The process responsible for such a large viscosity has not been identified with certainty yet. Many processes have been proposed, some of which demonstrably work, though often not with an efficiency as high as the observations of CV outbursts seem to indicate. Other ideas, such as certain turbulence models, do not have much predictive power and are based on ad-hoc assumptions about hydrodynamic instabilities in disks. In order to compare the viscosities in disks under different conditions, one introduces a dimensionless vsicosity $`\alpha `$:
$$\nu =\alpha \frac{c_\mathrm{s}^2}{\mathrm{\Omega }},$$
(10)
where $`c_\mathrm{s}`$ is the isothermal sound speed as before. The quantity $`\alpha `$ was introduced by Shakura and Sunyaev (1973), as a way of parametrizing our ignorance of the angular momentum transport process (their definition is based on a different formula however, and differs by a constant of order unity).
How large can the value of $`\alpha `$ be, on theoretical grounds? As a simple model, let’s assume that the shear flow between Kepler orbits is unstable to the same kind of shear instabilities found for flows in tubes, channels, near walls and in jets. These instabilities occur so ubiquitously that the fluid mechanics community considers them a natural and automatic consequence (e.g. DiPrima and Swinney 1981, p144 2nd paragraph) of a high Reynolds number:
$$\mathrm{R}e=\frac{LV}{\nu }$$
where $`L`$ and $`V`$ are characteristic length and velocity scales of the flow. If this number exceeds about 1000 (for some forms of instability much less), instability and turbulence are generally observed. It has been argued (e.g. Zel’dovich 1981) that for this reason hydrodynamic turbulence is the cause of disk viscosity. Let’s look at the consequences of this assumption. If an eddy of radial length scale $`l`$ develops due to shear instability, it will rotate at a rate given by the rate of shear, $`\sigma `$, in the flow, here
$$\sigma =r\frac{\mathrm{\Omega }}{r}\frac{3}{2}\mathrm{\Omega }.$$
The velocity amplitude of the eddy is $`V=\sigma l`$, and a field of such eddies produces a turbulent viscosity of the order (leaving out numerical factors of order unity)
$$\nu _{\mathrm{turb}}=l^2\mathrm{\Omega }.$$
(11)
In compressible flows, there is a maximum to the size of the eddy set by causality considerations. The force that allows an instability to form an overturning eddy is the pressure, which transports information about the flow at the sound speed. The eddies formed by a shear instability can therefore not move faster than $`c_\mathrm{s}`$, hence their size does not exceed $`c_\mathrm{s}/\sigma H`$. At the same time, the largest eddies formed also have the largest contribution to the turbulent viscosity. Thus we should expect that the turbulent viscosity is given by eddies with size of the order $`H`$:
$$\nu H^2\mathrm{\Omega },$$
or
$$\alpha 1.$$
Does hydrodynamical turbulence along these lines exist in disks? Unfortunately, this question is still open, but current opinion is leaning toward the view that the angular momentum transport in sufficiently ionized disks is due a a small scale magnetic field (Shakura and Sunyaev 1973). This is discussed briefly in section 8.
## 5 Thin Disks: equations
Consider a thin (= cool, nearly Keplerian, cf. section 4.2) disk, axisymmetric but not stationary. Using cylindrical coordinates ($`r,\varphi ,z`$), (note that we have changed notation from $`\varpi `$ to $`r`$ compared with section 4.2) we define the surface density $`\mathrm{\Sigma }`$ of the disk as
$$\mathrm{\Sigma }=_{\mathrm{}}^{\mathrm{}}\rho dz2H_0\rho _0,$$
(12)
where $`\rho _0`$, $`H_0`$ are the density and scaleheight at the midplane. The approximate sign is used to indicate that the coefficient in front of $`H`$ in the last expression actually depends on details of the vertical structure of the disk. Conservation of mass, in terms of $`\mathrm{\Sigma }`$ is given by
$$\frac{}{t}(r\mathrm{\Sigma })+\frac{}{r}(r\mathrm{\Sigma }v_r)=0.$$
(13)
(derived by integrating the continuity equation over $`z`$). Since the disk is axisymmetric and nearly Keplerian, the radial equation of motion reduces to
$$v_\varphi ^2=GM/r.$$
(14)
The $`\varphi `$-equation of motion is
$$\frac{v_\varphi }{t}+v_r\frac{v_\varphi }{r}+\frac{v_rv_\varphi }{r}=F_\varphi ,$$
(15)
where $`F_\varphi `$ is the azimuthal component of the viscous force. By integrating this over height $`z`$ and using (13), one gets an equation for the angular momentum balance:
$$\frac{}{t}(r\mathrm{\Sigma }\mathrm{\Omega }r^2)+\frac{}{r}(r\mathrm{\Sigma }v_r\mathrm{\Omega }r^2)=\frac{}{r}(Sr^3\frac{\mathrm{\Omega }}{r}),$$
(16)
where $`\mathrm{\Omega }=v_\varphi /r`$, and
$$S=_{\mathrm{}}^{\mathrm{}}\rho \nu dz\mathrm{\Sigma }\nu .$$
(17)
The second approximate equality in (17) holds if $`\nu `$ can be considered independent of $`z`$. The right hand side of (16) is the divergence of the viscous angular momentum flux, and is derived most easily with a physical argument, as described in, e.g. Pringle (1981) or Frank et al. (1992)<sup>1</sup><sup>1</sup>1If you prefer a more formal derivation, the fastest way is to consult Landau and Lifshitz (1959) chapter 15 (hereafter LL). Noting that the divergence of the flow vanishes for a thin axisymmetric disk, the viscous stress $`\sigma `$ becomes (LL eq. 15.3)
$$\sigma _{ik}=\eta \left(\frac{v_i}{x_k}+\frac{v_k}{x_i}\right),$$
where $`\eta =\rho \nu `$. This can be written in cylindrical or spherical coordinates using LL eqs. (15.15-15.18). The viscous force is
$$F_i=\frac{\sigma _{ik}}{x_k}=\frac{1}{\eta }\frac{\eta }{x_k}\sigma _{ik}+\eta ^2v_i,$$
Writing the Laplacian in cylindrical coordinates, the viscous torque is then computed from the $`\varphi `$-component of the viscous force by multiplying by $`r`$, and is then integrated over $`z`$..
Assume now that $`\nu `$ can be taken constant with height. For an isothermal disk ($`T`$ independent of $`z`$), this is equivalent to taking the viscosity parameter $`\alpha `$ independent of $`z`$. As long as we are not sure what causes the viscosity this is a reasonable simplification. Note, however, that recent numerical simulations of magnetic turbulence suggest that the effective $`\alpha `$, and the rate of viscous dissipation per unit mass, are higher near the disk surface than near the midplane. See the discussion in section 8. While eq (16) is still valid for rotation rates $`\mathrm{\Omega }`$ deviating from Keplerian (only the integration over disk thickness must be justifiable), we now use the fact that $`\mathrm{\Omega }r^{3/2}`$. Then Eqs. (13-16) can then be combined into a single equation for $`\mathrm{\Sigma }`$:
$$r\frac{\mathrm{\Sigma }}{t}=3\frac{}{r}[r^{1/2}\frac{}{r}(\nu \mathrm{\Sigma }r^{1/2})].$$
(18)
Under the same assumptions, eq. (15) yields the mass flux $`\dot{M}`$ at any point in the disk:
$$\dot{M}=2\pi r\mathrm{\Sigma }v_r=6\pi r^{1/2}\frac{}{r}(\nu \mathrm{\Sigma }r^{1/2}).$$
(19)
Eq. (18) is the standard form of the thin disk diffusion equation. An important conclusion from this equation is: in the thin disk limit, all the physics which determines the time dependent behavior of the disk enters through one quantitity only, the viscosity $`\nu `$. This is the main attraction of the thin disk approximation.
### 5.1 Steady Thin Disks
In a steady disk ($`/t=0`$) the mass flux $`\dot{M}`$ is constant through the disk and equal to the accretion rate onto the central object. From (19) we get the surface density distribution:
$$\nu \mathrm{\Sigma }=\frac{1}{3\pi }\dot{M}\left[1\beta \left(\frac{r_i}{r}\right)^{1/2}\right],$$
(20)
where $`r_i`$ is the inner radius of the disk and $`\beta `$ is a parameter appearing through the integration constant. It is related to the flux of angular momentum $`F_J`$ through the disk:
$$F_J=\dot{M}\beta \mathrm{\Omega }_ir_i^2,$$
(21)
where $`\mathrm{\Omega }_i`$ is the Kepler rotation rate at the inner edge of the disk. If the disk accretes onto an object with a rotation rate $`\mathrm{\Omega }_{}`$ less than $`\mathrm{\Omega }_i`$, one finds (Shakura and Sunyaev, 1973, Lynden-Bell and Pringle, 1974) that $`\beta =1`$, independent of $`\mathrm{\Omega }_{}`$. Thus the angular momentum flux (torque on the accreting star) is inward (spin-up) and equal to the accretion rate times the specific angular momentum at the inner edge of the disk. For stars rotating near their maximum rate ($`\mathrm{\Omega }_{}\mathrm{\Omega }_i`$) and for accretion onto magnetospheres, which can rotate faster than the disk, the situation is different (Sunyaev and Shakura 1977, Popham and Narayan 1991, Paczyński 1991, Bisnovatyi-Kogan 1993).
Accreting magnetospheres, for example, can spin down by interaction with the disk. This case has a surface density distribution (20) with $`\beta <1`$ (see also Spruit and Taam 1993). The angular momentum flux is then outward, and the accreting star spins down. This is possible even when the interaction between the disk and the magetosphere takes place only at the inner edge of the disk. Magnetic torques due interaction with the magetosphere may exist at larger distances in the disk as well, but are not necessary for creating an outward angular momentum flux. Recent numerical simulations of disk-magnetosphere interaction (Miller and Stone 1997) give an interesting new view of how such interaction may take place, and suggests it happens very differently from what is assumed in previous ‘standard’ models.
### 5.2 Disk Temperature
In this section I assume accretion onto not-too-rapidly rotating objects, so that $`\beta =1`$. The surface temperature of the disk, which determines how much energy it loses by radiation, is governed primarily by the energy dissipation rate in the disk, which in turn is given by the accretion rate. From the first law of thermodynamics we have
$$\rho T\frac{\mathrm{d}S}{\mathrm{d}t}=\mathrm{d}iv𝐅+Q_\mathrm{v},$$
(22)
where $`S`$ the entropy per unit mass, $`𝐅`$ the heat flux (including radiation and any form of ‘turbulent’ heat transport), and $`Q_\mathrm{v}`$ the viscous dissipation rate. For changes which happen on time scales longer than the dynamical time $`\mathrm{\Omega }^1`$, the left hand side is small compared with the terms on the right hand side. Integrating over $`z`$, the divergence turns into a surface term and we get
$$2\sigma _\mathrm{r}T_\mathrm{s}^4=_{\mathrm{}}^{\mathrm{}}Q_\mathrm{v}dz,$$
(23)
where $`T_\mathrm{s}`$ is the surface temperature of the disk, $`\sigma _\mathrm{r}`$ is Stefan-Boltzmann’s radiation constant $`\sigma _\mathrm{r}=a_\mathrm{r}c/4`$, and the factor 2 comes about because the disk has 2 radiating surfaces (assumed to radiate like black bodies). Thus the energy balance is local (for such slow changes): what is generated by viscous dissipation inside the disk at any radius $`r`$ is also radiated away from the surface at that position. The viscous dissipation rate is equal to $`Q_\mathrm{v}=\sigma _{ij}v_i/x_j`$, where $`\sigma _{ij}`$ is the viscous stress tensor (see footnote in section 5), and this works out<sup>2</sup><sup>2</sup>2using, e.g. LL eq. 16.3 to be
$$Q_\mathrm{v}=9/4\mathrm{\Omega }^2\nu \rho .$$
(24)
Eq. (23), using (20) then gives the surface temperature in terms of the accretion rate:
$$\sigma _\mathrm{r}T_\mathrm{s}^4=\frac{9}{8}\mathrm{\Omega }^2\nu \mathrm{\Sigma }=\frac{GM}{r^3}\frac{3\dot{M}}{8\pi }\left[1\left(\frac{r_i}{r}\right)^{1/2}\right].$$
(25)
This shows that the surface temperature of the disk, at a given distance $`r`$ from a steady accreter, depends only on the product $`M\dot{M}`$, and not on the highly uncertain value of the viscosity. For $`rr_i`$ we have
$$T_\mathrm{s}r^{3/4}.$$
(26)
These considerations only tells us about the surface temperature. The internal temperature in the disk is quite different, and depends on the mechanism transporting energy to the surface. Because it is the internal temperature that determines the disk thickness $`H`$ (and probably also the viscosity), this transport needs to be considered in some detail for realistic disk models. This involves a calculation of the vertical structure of the disk. Because of the local (in $`r`$) nature of the balance between dissipation and energy loss, such calculations can be done as a grid of models in $`r`$, without having to deal with exchange of energy between neighboring models. Schemes borrowed from stellar structure computations are used (e.g. Meyer and Meyer-Hofmeister 1982, Pringle et al. 1986, Cannizzo et al. 1988).
An approximation to the temperature in the disk can be found when a number of additional assumptions is made. As in stellar interiors, the energy transport is radiative rather than convective at high temperatures. Assuming local thermodynamic equilibrium (LTE, e.g. Rybicki and Lightman 1979), the temperature structure of a radiative atmosphere is given, in the Eddington approximation by:
$$\frac{\mathrm{d}}{\mathrm{d}\tau }\sigma _\mathrm{r}T^4=\frac{3}{4}F.$$
(27)
The boundary condition that there is no incident flux from outside the atmosphere yields the approximate condition
$$\sigma _\mathrm{r}T^4(\tau =2/3)=F,$$
(28)
where $`\tau =_z^{\mathrm{}}\kappa \rho dz`$ is the optical depth at geometrical depth $`z`$, and $`F`$ the energy flux through the atmosphere. Assuming that most of heat is generated near the midplane (which is the case if $`\nu `$ is constant with height), $`F`$ is approximately constant with height and equal to $`\sigma _\mathrm{r}T_\mathrm{s}^4`$, given by (25). Eq (27) then yields
$$\sigma _\mathrm{r}T^4=\frac{3}{4}(\tau +\frac{2}{3})F.$$
(29)
Approximating the opacity $`\kappa `$ as constant with $`z`$, the optical depth at the midplane is $`\tau =\kappa \mathrm{\Sigma }/2`$. If $`\tau 1`$, the temperature at the midplane is then:
$$T^4=\frac{27}{64}\sigma _\mathrm{r}^1\mathrm{\Omega }^2\nu \mathrm{\Sigma }^2\kappa .$$
(30)
With the equation of state (1), valid when radiation pressure is small, we find for the disk thickness, using (20):
$`{\displaystyle \frac{H}{r}}=`$ $`(/\mu )^{2/5}\left({\displaystyle \frac{3}{64\pi ^2\sigma _\mathrm{r}}}\right)^{1/10}(\kappa /\alpha )^{1/10}(GM)^{7/20}r^{1/20}(f\dot{M})^{1/5}`$ (31)
$`=`$ $`\mathrm{5\hspace{0.17em}10}^3\alpha ^{1/10}r_6^{1/20}\left(M/M_{}\right)^{7/20}(f\dot{M}_{16})^{1/5},(P_rP)`$ (32)
where $`r_6=r/(10^6`$ cm), $`\dot{M}_{16}=\dot{M}/(10^{16}`$g/s), and
$$f=1\left(r_i/r\right)^{1/2}.$$
From this we conclude that: i) the disk is thin in X-ray binaries, $`H/r<0.01`$, ii) the disk thickness is relatively insensitive to the parameters, especially $`\alpha `$, $`\kappa `$ and $`r`$. It must be stressed, however, that this depends fairly strongly on the assumption that the energy is dissipated in the disk interior. If the dissipation takes place close to the surface \[such as in some magnetic reconnection models (Haardt et al. 1994, Di Matteo et al. 1999 and references therein)\], the internal disk temperature will be much closer to the surface temperature. The midplane temperature and $`H`$ are even smaller in such disks than calculated from (32).
The viscous dissipation rate per unit area of the disk, $`W_\mathrm{v}=(9/4)\mathrm{\Omega }^2\nu \mathrm{\Sigma }`$ \[cf. eq. 25)\] can be compared with the local rate $`W_\mathrm{G}`$ at which gravitational energy is liberated in the accretion flow. Since half the gravitational energy stays in the flow as orbital motion, we have
$$W_\mathrm{G}=\frac{1}{2\pi r}\frac{GM\dot{M}}{2r^2},$$
(33)
so that
$$W_\mathrm{v}/W_\mathrm{G}=3f=3[1(r_\mathrm{i}/r)^{1/2}].$$
(34)
At large distances from the inner edge, the dissipation rate is 3 times larger than the rate of gravitational energy release. This may seem odd, but becomes understandable when it is realized that there is a significant flux of energy through the disk associated with the viscous stress<sup>3</sup><sup>3</sup>3See LL section 16. Integrating the viscous energy dissipation over the whole disk, one finds
$$_{r_i}^{\mathrm{}}2\pi rW_\mathrm{v}dr=\frac{GM\dot{M}}{2r_\mathrm{i}},$$
(35)
as expected. That is, globally, but not locally, half of the gravitational energy is radiated from the disk while the other half remains in the orbital kinetic energy of the accreted material.
What happens to this remaining orbital energy depends on the nature of the accreting object. If the object is a slowly rotating black hole, the orbital energy is just swallowed by the hole. If it has a solid surface, the orbiting gas slows down until it corotates with the surface, dissipating the orbital energy into heat in a boundary layer. Unless the surface rotates close to the orbital rate (‘breakup’), the energy released in this way is of the same order as the total energy released in the accretion disk. The properties of this boundary layer are therefore crucial for accretion onto neutron stars and white dwarfs. See also section 9.1 and Inogamov and Sunyaev (1999, and elsewhere in this volume).
### 5.3 Radiation pressure dominated disks
In the inner regions of disks in XRB, the radiation pressure can dominate over the gas pressure, which results in a different expression for the disk thickness. The total pressure $`P`$ is
$$P=P_r+P_g=\frac{1}{3}aT^4+P_g.$$
(36)
Defining a ‘total sound speed’ by $`c_\mathrm{t}^2=P/\rho `$ the relation $`c_\mathrm{t}=\mathrm{\Omega }H`$ still holds. For $`P_rP_g`$ we get from (30), with (25) and $`\tau 1`$:
$$cH=\frac{3}{8\pi }\kappa f\dot{M},$$
(where the rather approximate relation $`\mathrm{\Sigma }=2H\rho _0`$ has been used). Thus,
$$\frac{H}{R}\frac{3}{8\pi }\frac{\kappa }{cR}f\dot{M}=\frac{3}{2}f\frac{\dot{M}}{\dot{M}_\mathrm{E}},$$
(37)
where $`R`$ is the stellar radius and $`\dot{M}_\mathrm{E}`$ the Eddington rate for this radius. It follows that the disk becomes thick near the star, if the accretion rate is near Eddington (though this is mitigated somewhat by the decrease of the factor $`f`$). Accretion near the Eddington limit is evidently not geometrically thin any more. In addition, other processes such as angular momentum loss by ‘photon drag’ have to be taken into account.
### 5.4 Time scales in a disk
Three locally defined time scales play a role in thin disks. The dynamical time scale $`t_\mathrm{d}`$ is the orbital time scale:
$$t_\mathrm{d}=\mathrm{\Omega }^1=(GM/r^3)^{1/2}.$$
(38)
The time scale for radial drift through the disk over a distance of order $`r`$ is the viscous time scale:
$$t_\mathrm{v}=r/(v_r)=\frac{2}{3}\frac{rf}{\nu }=\frac{2f}{3\alpha \mathrm{\Omega }}(\frac{r}{H})^2,$$
(39)
(using (19 and (20), valid for steady accretion). Finally, there are thermal time scales. If $`E_\mathrm{t}`$ is the thermal energy content (enthalpy) of the disk per unit of surface area, and $`W_\mathrm{v}=(9/4)\mathrm{\Omega }^2\nu \mathrm{\Sigma }`$ the heating rate by viscous dissipation, we can define a heating time scale:
$$t_\mathrm{h}=E_\mathrm{t}/W_\mathrm{v}.$$
(40)
In the same way, a cooling time scale is defined by the energy content and the radiative loss rate:
$$t_\mathrm{c}=E_\mathrm{t}/(2\sigma _\mathrm{r}T_\mathrm{s}^4).$$
(41)
For a thin disk, the two are equal since the viscous energy dissipation is locally balanced by radiation from the two disk surfaces. \[In thick disks (ADAFs), this balance does not hold, since the advection of heat with the accretion flow is not negligible. In ADAFs, $`t_\mathrm{c}>t_\mathrm{h}`$ (see elsewhere in this volume)\]. Thus, we can replace both time scales by a single thermal time scale $`t_\mathrm{t}`$, and find, with (24):
$$t_\mathrm{t}=\frac{1}{W_\mathrm{v}}_{\mathrm{}}^{\mathrm{}}\frac{\gamma P}{\gamma 1}dz,$$
(42)
where the enthalpy of an ideal gas of constant ratio of specific heats $`\gamma `$ has been used. Leaving out numerical factors of order unity, this yields
$$t_\mathrm{t}\frac{1}{\alpha \mathrm{\Omega }}.$$
(43)
That is, the thermal time scale of the disk is independent of most of the disk properties and of the order $`1/\alpha `$ times longer than the dynamical time scale. This independence is a consequence of the $`\alpha `$-parametrization used. If $`\alpha `$ is not a constant, but dependent on disk temperature for example, the dependence of the thermal time scale on disk properties will become apparent again.
If, as seems likely from observations, $`\alpha `$ is generally $`<1`$, we have in thin disks the ordering of time scales:
$$t_\mathrm{v}t_\mathrm{t}>t_\mathrm{d}.$$
(44)
## 6 Comparison with CV observations
The number of meaningful quantitative tests between the theory of disks and observations is somewhat limited since in the absence of a theory for $`\nu `$, it is a bit meagre on predictive power. The most detailed information perhaps comes from modeling of CV outbursts.
Two simple tests are possible (nearly) independently of $`\nu `$. These are the prediction that the disk is geometrically quite thin (eq. 32) and the prediction that the surface temperature $`T_\mathrm{s}r^{3/4}`$ in a steady disk. The latter can be tested in a subclass of the CV’s that do not show outbursts, the nova-like systems, which are believed to be approximately steady accreters. If the system is also eclipsing, eclipse mapping techniques can be used to derive the brightness distribution with $`r`$ in the disk (Horne, 1985, 1993). If this is done in a number of colors so that bolometric corrections can be made, the results (e.g. Rutten et al. 1992) show in general a fair agreement with the $`r^{3/4}`$ prediction. Two deviations occur: i) a few systems show significantly flatter distributions than predicted, and ii) most systems show a ‘hump’ near the outer edge of the disk. The latter deviation is easily explained, since we have not taken into account that the impact of the stream heats the outer edge of the disk. Though not important for the total light from the disk, it is an important local contribution near the edge.
Eclipse mapping of Dwarf Novae in quiescence gives a quite different picture. Here, the inferred surface temperature profile is often nearly flat (e.g. Wood et al. 1989a, 1992). This is understandable however since in quiescence the mass flux depends strongly on $`r`$. In the inner parts of the disk it is small, near the outer edge it is close to its average value. With eq. (25), this yields a flatter $`T_\mathrm{s}(r)`$. The lack of light from the inner disk is compensated during the outburst, when the accretion rate in the inner disk is higher than average (see Mineshige and Wood 1989 for a more detailed comparison). The effect is also seen in the 2-dimensional hydrodynamic simulations of accretion in a binary by Różyczka and Spruit (1993). These simulations show an outburst during which the accretion in the inner disk is enhanced, between two episodes in which mass accumulates in the outer disk.
## 7 Comparison with LMXB observations: irradiated disks
In low mass X-ray binaries a complication arises because of the much higher luminosity of the accreting object. Since a neutron star is roughly 1000 times smaller than a white dwarf, it produces 1000 times more luminosity for a given accretion rate.
Irradiation of the disk by the central source leads to a different surface temperature than predicted by (25). The central source (star plus inner disk) radiates the total accretion luminosity $`GM\dot{M}/R`$ (assuming sub-Eddington accretion, see section 2). If the disk is concave, it will intercept some of this luminosity. If the central source is approximated as a point source the irradiating flux on the disk surface is
$$F_{\mathrm{irr}}=ϵ\frac{GM\dot{M}}{4\pi Rr^2},$$
(45)
where $`ϵ`$ is the angle between the disk surface and the direction from a point on the disk surface to the central source:
$$ϵ=\mathrm{d}H/\mathrm{d}rH/r.$$
(46)
The disk is concave if $`ϵ`$ is positive. We have
$$\frac{F_{\mathrm{irr}}}{F}=\frac{2}{3}\frac{ϵ}{f}\frac{r}{R},$$
where $`F`$ is the flux generated internally in the disk, given by (25). On average, the angle $`ϵ`$ is of the order of the aspect ratio $`\delta =H/r`$. With $`f1`$, and our fiducial value $`\delta \mathrm{5\hspace{0.17em}10}^3`$, we find that irradiation in LMXB dominates for $`r>10^9`$cm. This is compatible with observations (for reviews see van Paradijs and McClintock 1993), which show that the optical and UV are dominated by reprocessed radiation.
When irradiation by an external source is included in the thin disk model, the surface boundary condition of the radiative transfer problem, equation (28) becomes
$$\sigma _\mathrm{r}T_\mathrm{s}^4=F+(1a)F_{\mathrm{irr}},$$
(47)
where $`a`$ is the X-ray albedo of the surface, i.e. $`1a`$ is the fraction of the incident flux that is absorbed in the optically thick layers of the disk (photons absorbed higher up only serve to heat up the corona of the disk). The surface temperature $`T_\mathrm{s}`$ increases in order to compensate for the additional incident heat flux. The magnitude of the incident flux is sensitive to the assumed disk shape $`H(r)`$, as well as on the assumed shape (plane or spherical, for example) of the central X-ray emitting region. The disk thickness depends on temperature, and thereby also on the irradiation. It turns out, however, that this dependence on the irradiating flux is small, if the disk is optically thick, and the energy transport is by radiation (Lyutyi and Sunyaev 1976). To see this, integrate (27) with the modified boundary condition (47). This yields
$$\sigma _\mathrm{r}T^4=\frac{3}{4}F(\tau +\frac{2}{3})+\frac{(1a)F_{\mathrm{irr}}}{F}.$$
(48)
The irradiation adds an additive constant to $`T^4(z)`$. At the midplane, this constant has much less effect than at the surface. For the midplane temperature and the disk thickness to be affected significantly, it is necessary that
$$F_{\mathrm{irr}}/F>\tau .$$
(49)
The reason for this weak dependence of the midplane conditions on irradiation is the same as in radiative envelopes of stars, which are also insensitive to the surface boundary condition. The situation is very different for convective disks. As in fully convective stars, the adiabatic stratification then causes the conditions at the midplane to depend much more directly on the surface temparture. The outer parts of the disks in LMXB with wide orbits may be convective, and their thickness affected by irradiation.
In the reprocessing region of the disks of LMXB, the conditions are such that $`F<<F_{\mathrm{irr}}\tau F`$, hence we must use eq. (32) for $`H`$. This yields $`ϵ=(21/20)H/r\mathrm{5\hspace{0.17em}10}^3`$, and $`T_\mathrm{s}r^{0.5}`$, and we still expect the disk to remain thin.
From the paucity of sources in which the central source is eclipsed by the companion one deduces that the companion is barely or not at all visible from the inner disk, presumably because the outer parts of the disk are much thicker than expected from the above arguments. This is consistent with the observation that the characteristic modulation of the optical light curve due to irradiation of the secondary’s surface by the X-rays is not very strong in LMXB (with the exception of Her X-1, which has a large companion). The place of the eclipsing systems is taken by the so-called ‘Accretion Disk Corona’ (ADC) systems, where shallow eclipses of a rather extended X-ray source are seen instead of the expected sharp eclipses of the inner disk (for reviews of the observations, see Lewin et al. 1995). The conclusion is that there is an extended X-ray scattering ‘corona’ above the disk. It scatters a few per cent of the X-ray luminosity.
What causes this corona and the large inferred thickness of the disk ? The thickness expected from disk theory is a rather stable small number. To ‘suspend’ matter at the inferred height of the disk forces are needed that are much larger than the pressure forces available in an optically thick disk. A thermally driven wind, produced by X-ray heating of the disk surface, has been invoked (Begelman et al. 1983, Schandl and Meyer 1994). For other explanations, see van Paradijs and McClintock (1995). Perhaps a magnetically driven wind from the disk, such as seen in protostellar objects (e.g. Königl and Ruden 1993) can explain both the shielding of the companion and the scattering. Such a model would resemble magnetically driven wind models for the broad-line region in AGN (e.g. Emmering et al., 1992, Königl and Kartje 1994). A promising possibility is that the reprocessing region at the disk edge consists of matter ‘kicked up’ at the impact of the mass transfering stream (Meyer-Hofmeister et al.1997, Armitage and Livio 1998, Spruit et al. 1998). This produces qualitatively the right dependence of X-ray absorption on orbital phase in ADC sources, and the light curves of the so-called supersoft sources.
### 7.1 Transients
Soft X-ray transients (also called X-ray Novae) are believed to be binaries similar to the other LMXB, but somehow the accretion is episodic, with very large outbursts recurring on time scales of decades (sometimes years). There are many black hole candidates among these transients (see Lewin et al. 1995 for a review). As with the Dwarf Novae, the time dependence of the accretion in transients can in principle be exploited to derive information on the disk viscosity, assuming that the outburst is caused by an instability in the disk. The closest relatives of soft transients among the White Dwarf plus main sequence star systems are probably the WZ Sge stars (van Paradijs and Verbunt 1984, Kuulkers et al. 1996), which show (in the optical) similar outbursts with similar recurrence times (cf. Warner 1987, O’Donoghue et al. 1991). Like the soft transients, they have low mass ratios ($`q<0.1`$). For a given angular momentum loss, systems with low mass ratios have low mass transfer rates, so the speculation is that the peculiar behavior of these systems is somehow connected with a low mean accretion rate.
### 7.2 Disk Instability
The most developed model for outbursts is the disk instability model of Osaki (1974), Hōshi (1979), Smak (1971, 1984), Meyer and Meyer-Hofmeister (1981), see also King (1995), Osaki (1993). In this model the instability that gives rise to cyclic accretion is due to a temperature dependence of the viscous stress. In any local process that causes an effective viscosity, the resulting $`\alpha `$\- parameter will be a function of the main dimensionless parameter of the disk, the aspect ratio $`H/r`$. If this is a sufficiently rapidly increasing function, such that $`\alpha `$ is large in hot disks and low in cool disks, an instability results by the following mechanism. Suppose we start the disk in a stationary state at the mean accretion rate. If this state is perturbed by a small temperature increase, $`\alpha `$ goes up, and by the increased viscous stress the mass flux $`\dot{M}`$ increases. By (25) this increases the disk temperature further, resulting in a runaway to a hot state. Since $`\dot{M}`$ is larger than the average, the disk empties partly, reducing the surface density and the central temperature (eq. 30). A cooling front then transforms the disk to a cool state with an accretion rate below the mean. The disk in this model switches back and forth between hot and cool states. By adjusting $`\alpha `$ in the hot and cool states, or by adjusting the functional dependence of $`\alpha `$ on $`H/r`$, outbursts are obtained that agree reasonably with the observations of soft transients (Lin and Taam 1984, Mineshige and Wheeler, 1989). A rather strong increase of $`\alpha `$ with $`H/r`$ is needed to get the observed long recurrence times.
Another possible mechanism for instability has been found in 2-D numerical simulations of accretion disks (Blaes and Hawley 1988, Różyczka and Spruit 1993). The outer edge of a disk is found, in these simulations, to become dynamically unstable to a oscillation which grows into a strong excentric perturbation (a crescent shaped density enhancement which rotates at the local orbital period). Shock waves generated by this perturbation spread mass over most of the Roche lobe, at the same time the accretion rate onto the central object is strongly enhanced. This process is different from the Smak-Osaki-Hōshi mechanism, since it requires 2 dimensions, and does not depend on the viscosity (instead, the internal dynamics in this instability generates the effective viscosity that causes a burst of accretion).
### 7.3 Other Instabilities
Instability to heating/cooling of the disk can be the due to several effects. The cooling rate of the disk, if it depends on temperature in an appropriate way, can cause a thermal instability like that in the interstellar medium. Other instabilities may result from the dependence of viscosity on conditions in the disk. For a general treatment see Piran (1978), for a shorter discussion see Treves et al., 1988.
## 8 Sources of Viscosity
The high Reynolds number of the flow in accretion disks (of the order $`10^{11}`$ in the outer parts of a CV disk) would, to most fluid dynamicists, seem an amply sufficient condition for the occurrence of hydrodynamic turbulence. A theoretical argument against such turbulence often used in astrophysics (Kippenhahn and Thomas 1981, Pringle 1981) is that in cool disks the gas moves almost on Kepler orbits, which are quite stable (except for the orbits that get close to the companion or near a black hole). This stability is related to the known stabilizing effect that rotation has on hydrodynamical turbulence (Bradshaw 1969, for a discussion see Tritton 1992). Kippenhahn and Thomas also point out that the one laboratory experiment that comes close to the situation in accretion disks, namely the rotating Couette flow, does not become unstable for parameters like in disks (for the rather limited range in Reynolds numbers available). A (not very strong) observational argument is that hydrodynamical turbulence as described above would produce an $`\alpha `$ that does not depend on the nature of the disk, so that all objects should have the same value. This is unlikely to be the case. From the modeling of CV outbursts one knows, for example, that $`\alpha `$ probably increases with temperature (more accurately, with $`H/r`$, see previous section). Also, there are indications from the inferred life times and sizes of protostellar disks (Strom et al. 1993) that $`\alpha `$ may be rather small there, $`10^3`$, whereas in outbursts of CV’s one infers values of the order $`0.11`$.
The indeterminate status of the hydrodynamic turbulence issue is an annoying problem in disk theory. Direct 3-D numerical simulation of the hydrodynamics in accretion disks is possible, and so far has not shown the expected turbulence. In fact, Balbus and Hawley (1996), and Hawley et al (1999) argue, on the basis of such simulations and a physical argument, that disks are actually quite stable against hydrodynamic turbulence, as long as the specific angular momentum increases outward. \[Such heresy would not pass a referee in a fluid mechanics journal.\] If it is true that disks are stable to hydrodynamic turbulence it will be an uphill struggle to convince the fluid mechanics community, since it can always be argued that one should go to even higher Reynolds numbers to see the expected turbulence in the simulations or experiments.
The astrophysical approach has been to circumvent the problem by finding plausible alternative mechanisms that might work just as well. Among the processes that have been proposed repeatedly as sources of viscosity is convection due to a vertical entropy gradient (e.g. Kley et al. 1993), which may have some limited effect in convective parts of disks. Another class are waves of various kinds. Their effect can be global, that is, not reducible to a local viscous term because by traveling across the disk they can communicate torques over large distances. For example, waves set up at the outer edge of the disk by tidal forces can travel inward and by dissipating there can effectively transport angular momentum outward (e.g. Narayan et al. 1987, Spruit et al. 1987). A nonlinear version of this idea are selfsimilar spiral shocks, observed in numerical simulations (Sawada et al. 1987) and studied analytically (Spruit 1987). Such shocks can produce accretion at an effective $`\alpha `$ of $`0.01`$ in hot disks, but are probably not very effective in disks as cool as those in CV’s and XRB. A second non-local mechanism is provided by a magnetically accelerated wind originating from the disk surface (Blandford 1976, Bisnovatyi-Kogan and Ruzmaikin 1976, Lovelace 1976, Blandford and Payne 1982, for reviews see Blandford 1989, Blandford and Rees 1992, for an introduction see Spruit 1996). In principle, such winds can take care of all the angular momentum loss needed to make accretion possible in the absence of a viscosity (Blandford 1976, Königl 1989). The attraction of this idea is that magnetic winds are a strong contender for explaining the strong outflows and jets seen in protostellar objects and AGN. It is not yet clear however if, even in these objects, the wind is actually the main source of angular momentum loss.
In sufficiently cool or massive disks, selfgravitating instabilities of the disk matter can produce internal friction. Paczýnski (1978) has proposed that the resulting heating would limit the instability and keep the disk in a well defined moderately unstable state. The angular momentum transport in such a disk has been modeled by several authors (e.g. Ostriker et al. 1999). Disks in XRB are too hot for selfgravity to play a role.
### 8.1 magnetic viscosity
Magnetic forces can be very effective at transporting angular momentum. If it can be shown that the shear flow in the disk produces some kind of small scale fast dynamo process, that is, some form of magnetic turbulence, an effective $`\alpha O(1)`$ expected (Shakura and Sunyaev 1973, Eardley and Lightman 1975, Pudritz 1981, Meyer and Meyer-Hofmeister 1982). Numerical simulations of initially weak magnetic fields in accretion disks have now shown that this does indeed happen in sufficiently ionized disks (Hawley et al. 1995, Brandenburg et al. 1995, Armitage 1998). These show a small scale magnetic field with azimuthal component dominating (due to stretching by differential rotation). The effective $`\alpha `$’s are of the order 0.05. The angular momentum transport is due to magnetic stresses. The fluid motions induced by the magnetic forces contribute only little to the angular momentum transport. In a perfectly conducting plasma this turbulence can develop from an arbitrarily small initial field through magnetic shear instability (also called magetorotational instability, Velikhov 1959, Chandrasekhar 1961, Balbus and Hawley 1991, 1992). The significance of this instability is that it shows that at large conductivity accretion disks must be magnetic. The actual form of the highly time dependent small scale magnetic field which develops can only be found from numerical simulations.
### 8.2 viscosity in radiatively supported disks
A disk in which the radiation pressure $`P_\mathrm{r}`$ dominates must be optically thick (otherwise the radiation would escape). The radiation pressure then adds to the total pressure is larger than it would be, for a given temperature, if only the gas pressure were effective. If the viscosity is then parametrized by (10), it turns out (Lightman and Eardley, 1974) that the disk is locally unstable. An increase in temperature increases the radiation pressure, which increases the viscous dissipation and the temperature, leading to a runaway. This has raised the question whether the radiation pressure should be included in the sound speed that enters expression (10). If it is left out, a lower viscosity results, and there is no thermal-viscous runaway. Without knowledge of the process causing the effective viscous stress, this question can not be answered. Sakimoto and Coroniti (1989) have shown, however, that if the stress is due to some form of magnetic turbulence, it most likely scales with the gas pressure alone, rather than the total pressure. Now that it seems likely, from the numerical simulations, that the stress is indeed magnetic, there is reason to believe that in the radiation pressure-dominated case the effective viscosity will scale as $`\nu \alpha P_\mathrm{g}/(\rho \mathrm{\Omega })`$ (this case has not been studied with simulations yet). Nayakshin and Rappaport (1999) show that, depending on how the viscosity scales in the intermediate regime $`P_\mathrm{g}P_{\mathrm{rad}}`$, interesting cyclic behavior can occur akin to the ‘S-curve’ instability in CV disks (section 7.2).
## 9 Beyond thin disks
Ultimately, much of the progress in developing useful models of accretion disks will depend on detailed numerical simulations in 2 or 3 dimensions. In the disks one is interested in, there is usually a large range in length scales (in LMXB disks, from less than the $`10`$km neutron star radius to the more than $`10^5`$km orbital scale). Correspondingly, there is a large range in time scales that have to be followed. This not technically possible at present and in the foreseeable future. In numerical simulations one is therefore limited to studying in an approximate way aspects that are either local or of limited dynamic range in $`r,t`$ (for examples, see Hawley 1991, Różyczka and Spruit 1993, Armitage 1998). For this reason, there is still a need for approaches that relax the strict thin disk framework somewhat without resorting to full simulations. Due to the thin disk assumptions, the pressure gradient does not contribute to the support in the radial direction and the transport of heat in the radial direction is negligible. Some of the physics of thick disks can be included in a fairly consistent way in the ‘slim disk’ approximation (Abramowicz et al., 1988). The so-called Advection Dominated Accretion Flows (ADAFs) are related to this approach (for a review see Yi 1998, for an introduction Spruit, elsewhere in this volume).
### 9.1 Boundary layers
In order to accrete onto a star rotating at the rate $`\mathrm{\Omega }_{}`$, the disk matter must dissipate an amount of energy given by
$$\frac{GM\dot{M}}{2R}\left[1\mathrm{\Omega }_{}/\mathrm{\Omega }_k(R)\right]^2.$$
(50)
The factor in brackets measures the kinetic energy of the matter at the inner edge of the disk ($`r=R`$), in the frame of the stellar surface. Due to this dissipation the disk inflates into a ‘belt’ at the equator of the star, of thickness $`H`$ and radial extent of the same order. Equating the radiation emitted from the surface of this belt to (50) one gets for the surface temperature $`T_{\mathrm{sb}}`$ of the belt, assuming optically thick conditions and a slowly rotating star ($`\mathrm{\Omega }_{}/\mathrm{\Omega }_k1`$):
$$\frac{GM\dot{M}}{8\pi R^2H}=\sigma _\mathrm{r}T_{\mathrm{sb}}^4$$
(51)
To find the temperature inside the belt and its thickness, use eq. (29). The value of the surface temperature is higher, by a factor of the order $`(R/H)^{1/4}`$, than the simplest thin disk estimate (25, ignoring the $`(r/r_i)^{1/2}`$ factor). In practice, this works out to a factor of a few. The surface of the belt is therefore not very hot. The situation is quite different if the boundary layer is not optically thick (Pringle and Savonije 1979). It then heats up to much higher temperatures. Analytical methods to obtain the boundary layer structure have been used by Regev and Hougerat (1988), numerical solutions of the slim disk type by Narayan and Popham (1993), Popham (1997), 2-D numerical simulations by Kley (1991). These considerations are primarily relevant for CV disks; in accreting neutron stars, the dominant effects of radiation pressure have to be included. More analytic progress on the structure of the boundary layer between a disk and a neutron star and the way in which it spreads over the surface of the star is reported by Inogamov and Sunyaev (1999, see also elsewhere in this volume).
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# On the controversy over the stochastic density functional equations
## 1 Introduction
In supercooled liquids, due to the dense packing and strong correlation of the constituent particles, the nonvibrational diffusive motions take much more time than collisions. In other words, the momentum and the energy flow much more quickly via collisions through the system than the slowly decaying number density. Recently there has been a considerable effort to describe such slow dynamics in liquids . We shall in particular discuss one of the approaches, the following stochastic equations in terms of the number density field $`\rho (x,t)`$ \[2-5\]:
$`{\displaystyle \frac{\rho (x,t)}{t}}=L[\rho (x,t)]{\displaystyle \frac{\delta H(\{\rho \})}{\delta \rho }}|_{\rho (x,t)}+\xi (x,t),`$ (1)
or its equivalent, i.e. the Fokker-Planck equation for the probability distribution functional $`P(\{\rho \},t)`$:
$$\frac{P(\{\rho \},t)}{t}=𝑑x\frac{\delta }{\delta \rho (x)}L[\rho (x)]\left[T\frac{\delta }{\delta \rho (x)}+\frac{\delta H(\{\rho \})}{\delta \rho (x)}\right]P(\{\rho \},t).$$
(2)
In equations (1) and (2), the Hamiltonian $`H`$ is of the free energy functional form as
$$H(\{\rho \})=\frac{1}{2}𝑑x𝑑y\rho (x)V(xy)\rho (y)+T𝑑x\rho (x)\mathrm{log}\rho (x),$$
(3)
$`L[\rho (x)]`$ is the kinetic coefficient written as $`L[\rho (x)]=\rho (x,t)\mathrm{\Gamma }`$ with $`\mathrm{\Gamma }`$ being the mobility of particles, and $`\xi `$ is the divergence of a random force and its correlation function is given by
$`\xi (x,t)\xi (y,t^{})=2T_xL[\rho (x)]_y\delta (xy)\delta (tt^{}).`$ (4)
For the explicit representation of this averaging, see equation (9) below.
These equations have attractive features for studying the slow dynamics: One is the physically clear incorporation of the thermally activated hopping processes via the last random term $`\xi `$ on the right hand side (rhs) of equation (1). Also of interest is the density dependence of the kinetic coefficient, $`L\rho `$, that produces the nonlinear term of dynamic origin and implies the relevance of these equations to describing dynamical heterogeneity .
Nevertheless the stochasticity for the number density variable is still controversial. To see this, let us first mention one of some attempts \[2-4\], presented by Dean , to justify the above stochastic equations (1) and (2): Consider here a canonical system of $`N`$ particles interacting via a pairwise potential $`V(x)`$ and surrounded by a thermal white noise heat bath. The $`i`$th particle then obeys the Langevin equation,
$$\frac{dX_i(t)}{dt}=\mathrm{\Gamma }\underset{j=1}{\overset{N}{}}_iV[X_i(t)X_j(t)]+\eta _i(t),$$
(5)
where the components of the noise $`\eta _i^\alpha `$ are taken to be uncorrelated as $`\eta _i^\alpha (t)\eta _i^\beta (t)=2T\mathrm{\Gamma }\delta _{ij}\delta ^{\alpha \beta }\delta (tt^{})`$. Dean shows, using the Ito prescription for the change of variables and summing over the $`i`$, that equation (5) is transformed to the equation of the density operator, $`\widehat{\rho }(x,t)=_i\widehat{\rho }_i(x,t)=_i\delta [xX_i(t)]`$:
$`\begin{array}{cc}\hfill {\displaystyle \frac{\widehat{\rho }(x,t)}{t}}=\widehat{\rho }(x,t){\displaystyle 𝑑y\widehat{\rho }(y,t)V(xy)}+T^2\widehat{\rho }(x,t)+\widehat{\xi }(x,t)& \\ \hfill \widehat{\xi }(x,t)={\displaystyle \underset{i}{}}[\widehat{\rho }_i(x,t)\eta _i(t)].& \end{array}`$ (8)
Since one finds
$`\widehat{\xi }(x,t)\widehat{\xi }(y,t)`$ $`{\displaystyle 𝑑\eta _i\widehat{\xi }(x,t)\widehat{\xi }(y,t)\mathrm{exp}\left(𝑑t\frac{\eta _i^2}{4T\mathrm{\Gamma }}\right)}`$ (9)
$`=2T_xL[\widehat{\rho }(x)]_x\delta (tt^{}),`$ (10)
equation (1) is verified so long as the operator $`\widehat{\rho }`$ reads $`\rho `$ of c-number.
As expected, though, Marconi and Tarazona (MT) subsequently object to the last supposition: they claim that $`\rho `$ is to be defined by averaging $`\widehat{\rho }`$ over the noise as $`\rho _{av}=\widehat{\rho }`$, where the subscript $`av`$ is appended for emphasizing the procedure. As a consequence dynamical density functional equation becomes deterministic:
$$\frac{\rho _{av}(x,t)}{t}=𝑑y\widehat{\rho }(x,t)\widehat{\rho }(y,t)V(xy)+T^2\rho _{av}(x,t),$$
(11)
whereby the Boltzmann distribution of number density is assured as time-independent solution .
To settle such controversy over the stochastic density functional equations (1) and (2), this letter aims at complementing Dean’s argument from (5) to (10) so that the above criticism by MT may become invalid. Our strategy is twofold: First, in the next section, we demonstrate that standard manipulations transform equation (8) of the density operator to the Fokker-Planck equation (2). Moreover we verify in section 3, with the help of the conditional grand canonical partition function, the static solution $`P_0(\{\rho \})`$ of (2):
$$P_0(\{\rho \})\mathrm{exp}(\beta H),$$
(12)
where $`\beta =T^1`$. In the final section, to clarify the connection between the stochastic and the deterministic equation, we confirm using the WKB-like approximation to the Fokker-Planck equation (2) that the noise-averaged deterministic equation (11) corresponds to that for the saddle-point path of $`P(\{\rho \},t)`$; this reveals that MT’s argument produces only the mean-field equation and not the first member of the BBGKY hierarchy including two-point equal-time correlation function.
## 2 From equation (8) to the Fokker-Planck equation (2)
Turning our attention to the functional space, we immediately find that the density operator $`\widehat{\rho }`$ may be directly mapped to the distribution functional $`P(\{\rho \},t)`$ as
$$P(\{\rho \},t)=\underset{x}{}\delta [\widehat{\rho }(x,t)\rho (x)],$$
(13)
not via the averaging $`\rho _{av}=\widehat{\rho }`$; essentially only this definition has dissolved the MT’s critique.
Let us then exhibit below that $`P(\{\rho \},t)`$ with equation (8) satisfies the Fokker-Planck equation (2). We first differentiate (13) with respect to time, and obtain
$`{\displaystyle \frac{P(\{\rho \},t)}{t}}={\displaystyle 𝑑x\frac{\widehat{\rho }(x,t)}{t}\frac{\delta }{\delta \widehat{\rho }(x,t)}\delta [\widehat{\rho }(x,t)\rho (x)]\underset{yx}{}\delta [\widehat{\rho }(y,t)\rho (y)]}`$
$`\begin{array}{cc}\hfill ={\displaystyle }dx[\widehat{\rho }(x,t){\displaystyle }dy\widehat{\rho }(y,t)V(xy)+T^2\widehat{\rho }(x,t)+\widehat{\xi }(x,t)]& \\ \hfill \times {\displaystyle \frac{\delta }{\delta \widehat{\rho }(x,t)}}\delta [\widehat{\rho }(x,t)\rho (x)]{\displaystyle \underset{yx}{}}\delta [\widehat{\rho }(y,t)\rho (y)].& \end{array}`$ (16)
We may replace $`\widehat{\rho }`$ by $`\rho `$ using the $`\delta `$–function, and hence equation (16) reads
$`\begin{array}{cc}\hfill {\displaystyle \frac{P(\{\rho \},t)}{t}}={\displaystyle 𝑑x\frac{\delta }{\delta \rho (x)}\left[\rho (x)𝑑y\rho (y)V(xy)+T^2\rho (x)\right]P(\{\rho \},t)}& \\ \hfill {\displaystyle 𝑑x\frac{\delta }{\delta \rho (x)}\widehat{\xi }(x,t)\underset{x}{}\delta [\widehat{\rho }(x,t)\rho (x)]}.& \end{array}`$ (19)
With use of the identity for an arbitrary function $`F(\{\eta _i^\alpha \})`$, $`F(\{\eta _i^\alpha \})\eta _i^\alpha (t)=2T\mathrm{\Gamma }\delta F(\{\eta _i^\alpha \})/\delta \eta _i^\alpha (t)`$, the last bracket term on the rhs of (19) is further transformed to
$`\widehat{\xi }(x,t){\displaystyle \underset{x}{}}\delta [\widehat{\rho }(x,t)\rho (x)]=2T\mathrm{\Gamma }{\displaystyle \underset{i,\alpha }{}}{\displaystyle \frac{\widehat{\rho }_i(x,t)}{x^\alpha }}{\displaystyle \frac{\delta \widehat{\rho }(y,t)}{\delta \eta _i^\alpha }}{\displaystyle \frac{\delta }{\delta \widehat{\rho }(y,t)}}{\displaystyle \underset{y}{}}\delta [\widehat{\rho }\rho ]`$
$`=T\mathrm{\Gamma }{\displaystyle \underset{i}{}}_x_y[\widehat{\rho }_i(x,t)\widehat{\rho }_i(y,t)]{\displaystyle \frac{\delta }{\delta \widehat{\rho }(y,t)}}{\displaystyle \underset{y}{}}\delta [\widehat{\rho }\rho ],`$ (20)
where the superscript $`\alpha `$ denotes the component of $`x`$ and $`\eta _i`$, and use has been made of
$$\frac{\delta \widehat{\rho }_i(y,t)}{\delta \eta _i^\alpha }\frac{1}{2}\frac{\widehat{\rho }_i(y,t)}{y_i^\alpha }$$
(21)
that is obtained from standard mathematical manipulation of the discretized Langevin equation . Also, noting that relation $`\widehat{\rho }_i(x,t)\widehat{\rho }_i(y,t)=\delta (xy)\rho _i(x,t)`$ gives
$$_x_y[\widehat{\rho }_i(x,t)\widehat{\rho }_i(y,t)]=_x\widehat{\rho }_i(x,t)_x\delta (xy)$$
(22)
and replacing $`\widehat{\rho }`$ by $`\rho `$ as before, the bracket term finally reads
$$\widehat{\xi }(x,t)\underset{x}{}\delta [\widehat{\rho }(x,t)\rho (x)]=T\mathrm{\Gamma }\rho (x,t)\frac{\delta P(\{\rho \},t)}{\delta \rho (x,t)}.$$
(23)
Equation (19) with this is none other than the Fokker-Planck equation (2).
Thus it has been demonstrated that the stochastic equation of density operator (8) leads to the Fokker-Planck equation (2) for the (c-number) density distribution functional \[or equation (1)\].
## 3 Verification of the equilibrium distribution functional (12)
We arrive at a time-independent solution of the Fokker-Planck equation in the large time limit: $`P_0(\{\rho \})=lim_t\mathrm{}P(\{\rho \},t)`$. Therefore it is plausible to suppose that the noise averaging in calculating $`P_0`$ becomes equivalent to the configurational one in equilibrium:
$$P_0(\{\rho \})\frac{1}{N!}\underset{i=1}{\overset{N}{}}𝑑X_i\underset{x}{}\delta [\widehat{\rho }(x,t)\rho (x)]\mathrm{exp}[\beta \underset{i,j}{}V(X_iX_j)].$$
(24)
The problem is then how to derive expression (12) from the above conditional partition function.
Let us once move to the grand canonical system where we are to consider
$$P_0^\mathrm{\Xi }=\underset{N=0}{\overset{\mathrm{}}{}}P_0\lambda ^N,$$
(25)
with $`\beta \mu =\mathrm{ln}\lambda `$ being chemical potential. Introducing the auxiliary field $`\psi `$ as $`\delta [\widehat{\rho }(x)\rho ]=𝑑\psi \mathrm{exp}[i\psi (\widehat{\rho }\rho )]`$, the configurational representation of $`P_0^\mathrm{\Xi }`$ given by (24) and (25) reads
$$P_0^\mathrm{\Xi }D\psi \mathrm{exp}\left[\beta 𝑑x𝑑y\frac{1}{2}\rho (x)V(xy)\rho (y)+i\rho (x)\psi (x)e^{i\psi (x)+\mu }\right],$$
(26)
where $`D\psi `$ is formally defined as $`_xd\psi (x)`$. Since there is no contribution to $`P_0^\mathrm{\Xi }`$ of principal quadratic fluctuation of the auxiliary field $`\psi `$ around the saddle point path $`\psi _{sp}`$ as shown elsewhere , Gaussian approximation for $`\psi `$ reduces the functional integral form (26) to
$$P_0^\mathrm{\Xi }\mathrm{exp}\left[\beta H+\beta 𝑑x\rho (x)+\mu \rho (x)\right],$$
(27)
as found from substituting $`\rho =e^{i\psi _{sp}+\mu }`$ into (26).
For returning to the canonical system, we have only to perform the following contour integral,
$$P_0=\frac{1}{2\pi i}𝑑\lambda \frac{P_0^\mathrm{\Xi }}{\lambda ^{N+1}},$$
(28)
where $`\lambda `$ is now a complex variable. This relation with use of the Cauchy’s integral theorem gives back the canonical form:
$`P_0\mathrm{exp}\left[\beta H+\beta {\displaystyle 𝑑x\rho (x)}\right]{\displaystyle \frac{1}{2\pi i}}{\displaystyle 𝑑\lambda \frac{1}{\lambda ^{1+[N{\scriptscriptstyle 𝑑x\rho (x)}]}}}`$
$`=\{\begin{array}{cc}e^{\beta (HN)}\hfill & \text{if }𝑑x\rho (x)=N\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}`$ (31)
The static solution (12) has been thus verified, and supplementing Dean’s discussion has been completed.
## 4 Discussion: connection with the deterministic equation (11)
Here we would like to first confirm the saddle-point path of the Fokker-Planck (2) from exploiting the WKB-like approach made in Ref. . Setting, similarly to the WKB approximation, that
$$P(\{\rho \},t)\mathrm{exp}\left[\beta \mathrm{\Phi }(\{\rho \},t)\right],$$
(32)
we obtain the Hamilton-Jacobi like equation:
$$\frac{\mathrm{\Phi }(\{\rho \},t)}{t}=𝑑x\frac{\delta \mathrm{\Phi }(\{\rho \},t)}{\delta \rho (x)}L[\rho (x)]\left[\frac{\delta \mathrm{\Phi }(\{\rho \},t)}{\delta \rho (x)}\frac{\delta H(\{\rho \})}{\delta \rho (x)}\right].$$
(33)
A short-cut way of deriving from this the most probable (or saddle-point) path $`\{\overline{\rho }\}`$ is to expand $`\mathrm{\Phi }`$ and $`H`$ around $`\{\overline{\rho }\}`$ as
$`\begin{array}{cc}\hfill \mathrm{\Phi }(\{\rho \},t)=\mathrm{\Phi }(\{\overline{\rho }\},t)+{\displaystyle \frac{1}{2}}{\displaystyle 𝑑x𝑑y[\rho (x)\overline{\rho }(x)]\mathrm{\Phi }\mathrm{"}(xy)[\rho (y)\overline{\rho }(y)]}+\mathrm{}& \\ \hfill J(\{\rho \})=J(\{\overline{\rho }\})+{\displaystyle \frac{\delta J(\{\rho \})}{\delta \rho }}|_{\{\rho \}=\{\overline{\rho }\}}[\rho (x)\overline{\rho }(x)]+\mathrm{},& \end{array}`$ (36)
with $`J(\{\rho \})L[\rho (x)]\delta H(\{\rho \})/\delta \rho (x)`$. Substitution of these into equation (33) yields in $`𝒪[\rho (x)\overline{\rho }(x)]`$
$$\frac{\overline{\rho }(x,t)}{t}=L[\overline{\rho }(x,t)]\frac{\delta H(\{\rho \})}{\delta \rho }|_{\{\rho \}=\{\overline{\rho }\}}.$$
(37)
Since $`\{\overline{\rho }\}`$ is to be in accord with the noise averaged density $`\{\rho _{av}\}`$, equation (37) implies that $`\widehat{\rho }(x,t)\widehat{\rho }(y,t)=\rho _{av}(x,t)\rho _{av}(y,t)`$ for the first term on the rhs of (11), i.e., no spatial correlation of noise-averaging. In other words, the above derivation reveals that the noise-averaged equation (11) is not the first member of the dynamical BBGKY hierarchy unlike the proposal by MT, but is only the mean-field equation for the saddle-point path of $`P(\{\rho \},t)`$.
We have thus validated the stochastic density functional equations, which must be a powerful tool for the understanding of supercooled fluids and glasses, via proving the irrelevance of MT’s objection to Dean’s argument in three ways: (i) demonstrating that standard manipulations enable to replace with the c-number density field $`\rho `$ the corresponding operator variable $`\widehat{\rho }`$ in the stochastic equation (8) derived by Dean, (ii) verifying the static solution (12) of the Fokker-Planck equation for the density distribution functional with the help of the conditional grand canonical partition function, and (iii) pointing out that the noise averaged path satisfying the deterministic equation (11) merely corresponds to the saddle-point one.
The next problem is how to solve these dynamical density functional equations, stochastic or deterministic. In previous works , the static density functional theory has been exploited as input, and some justifications have been also described by Kawasaki and MT . However, the present discussion does not support these; from our point of view, what to suppose for incorporating the static theory remains an open problem.
We acknowledge the financial support from the Ministry of Education, Science, Culture, and Sports of Japan.
## References
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# Parametric Amplification of Metric Fluctuations During Reheating in Two Field Models
## I Introduction
It was recently suggested that parametric resonance during the reheating phase of an inflationary Universe may lead to an exponential amplification of super-Hubble scale gravitational fluctuations. If true, this would affect the usual predictions of inflationary models for observables such as the matter power spectrum and the spectrum of cosmic microwave anisotropies. In particular, it would require the coupling constants in the particle physics model of inflation to be exponentially smaller than previously thought in order that the theory does not generate a too large amplitude for the fluctuations.
In Ref. it was shown that, although there are no causality constraints which prohibit the amplification of super-Hubble (but sub-horizon) modes during reheating, the effect does not occur in a simple massive scalar field model of chaotic inflation based on the potential $`V(\varphi )=m^2\varphi ^2/2`$ (Here, $`\varphi `$ is the inflaton field). This is true even beyond the linear analysis . Similarly, there is no effect for a quartic potential $`V(\varphi )=\lambda \varphi ^4/4`$, nor for a potential containing both quadratic and quartic terms . These results agree with the earlier analyses in Refs. and .
It was then suggested that the amplification of super-Hubble-scale modes would occur for two field models of inflation, e.g. for a model with potential
$$V(\varphi ,\chi )=\frac{1}{2}m^2\varphi ^2+\frac{1}{2}g^2\varphi ^2\chi ^2,$$
(1)
where as before $`\varphi `$ is the inflaton field and $`\chi `$ is a second scalar matter field. This model had earlier been analyzed by Taruya and Nambu who claimed that the isocurvature mode of the fluctuations will be parametrically amplified during reheating. However, as was shown in Refs. and (and more recently in ), the fluctuations in the $`\chi `$ field are exponentially suppressed during inflation for values of the coupling constant for which the equation of motion of the metric fluctuations corresponds to broad resonance, thus rendering the effect studied in Ref. completely inefficient.
The suppression of fluctuations in the $`\chi `$ field which renders the parametric amplification of gravitational fluctuations ineffective in the model given by (1) occurs since during inflation the induced mass $`m_\chi `$ of the $`\chi `$ field which is given by $`m_\chi =g|\varphi |`$ is larger than the Hubble expansion parameter $`H`$, and hence, as can be easily seen by considering the equation of motion
$$\ddot{\delta }\chi +3H\dot{\delta }\chi +\left(\frac{k^2}{a^2}+g^2\varphi ^2\right)\delta \chi =\mathrm{\hspace{0.17em}0}$$
(2)
for the linearized fluctuation of the $`\chi `$ field with comoving wave number $`k`$ (the scale factor is denoted by $`a(t)`$), $`\delta \chi `$ undergoes damped oscillatory motion.
A model in which $`m_\chi <H`$ during the stage of inflation when scales of cosmological interest today exit the Hubble radius was recently studied by Bassett and Viniegra . It is a two field model given by the potential
$$V(\varphi ,\chi )=\frac{1}{4}\lambda \varphi ^4+\frac{1}{2}g^2\varphi ^2\chi ^2.$$
(3)
In the absence of metric fluctuations, this model was studied in detail in Ref. (see also ), where it was shown that for values of the coupling constants satisfying
$$g^2\mathrm{\hspace{0.17em}2}\lambda $$
(4)
long wavelength modes ($`k0`$) are in the first broad instability band of the Floquet-type equation of motion derived from (2) after field rescaling which describes the parametric resonance of matter fluctuations in an unperturbed expanding space-time. Bassett and Viniegra showed that in this model the quantity $`\zeta `$ increases exponentially during the initial stages of reheating. Note that $`\zeta `$ is a measure of the curvature fluctuations and is believed to be conserved on super-Hubble scales in the absence of isocurvature fluctuations (see e.g. for a review of the theory of cosmological fluctuations). However, since the model given by (3) admits isocurvature fluctuations, a growth of $`\zeta `$ on super-Hubble modes is expected also in the “usual” analysis of the evolution of fluctuations in inflationary cosmology.
In this paper we take a closer look at the theory given by the potential (3). Subject to certain assumptions on initial conditions of the background field dynamics we recover results similar to Bassett and Viniegra: exponential growth of $`\zeta `$ during the initial stages of reheating. Furthermore, we demonstrate that this effect is not taken into account by the conventional theory of isocurvature perturbations. We then discuss some criteria which an inflationary Universe model must satisfy in order to have substantial parametric growth of $`\zeta `$ during reheating via scalar field interactions. We argue that these conditions are naturally satisfied in some models of hybrid inflation, and we study a couple of concrete examples in which super-Hubble-scale gravitational fluctuations grow exponentially during reheating (some other examples where exponential growth of super-Hubble-scale modes could occur are given in ).
## II Massless Two Field Model Reconsidered
In this section we will consider the two-field model with potential (3) in which Bassett and Viniegra showed the parametric resonance of super-Hubble scale gravitational modes. This model has been studied in detail in in the absence of gravitational perturbations (see also ).
As in our previous paper , we shall work in longitudinal gauge in which the metric including linearized scalar metric fluctuations takes on the form
$$ds^2=(12\mathrm{\Phi })dt^2a^2(t)(1+2\mathrm{\Psi })dx^idx_i$$
(5)
where the $`x^i`$ are the comoving spatial coordinates and $`t`$ is physical time. Since for the matter model considered the off-diagonal components of the spatial part of the energy-momentum tensor vanish, the corresponding components of the Einstein equations imply $`\mathrm{\Psi }=\mathrm{\Phi }`$. For the moment we will write down the equations for a general system with multiple scalar fields ($`i=1,\mathrm{},n`$), and only at a later stage will we specialize to the specific model considered in this section.
The remaining independent equations of motion for linearized perturbations in this Einstein-Higgs system are the perturbed energy constraint and momentum constraint equations as well as the equations of motion for the Higgs field perturbations $`\delta \varphi _i`$:
$`3H\dot{\mathrm{\Phi }}`$ $``$ $`\left({\displaystyle \frac{k^2}{a^2}}+3H^2\right)\mathrm{\Phi }`$ (6)
$`=`$ $`4\pi G{\displaystyle \underset{i=1}{\overset{n}{}}}[\dot{\varphi }_i\dot{\delta \varphi _i}\mathrm{\Phi }\dot{\varphi _i}^2+V,_i\delta \varphi _i],`$ (7)
$`\dot{\mathrm{\Phi }}`$ $`+`$ $`H\mathrm{\Phi }=4\pi G{\displaystyle \underset{i=1}{\overset{n}{}}}\dot{\varphi }_i\delta \varphi _i,`$ (8)
$`\ddot{\delta \varphi _i}`$ $`+`$ $`3H\dot{\delta \varphi _i}+[{\displaystyle \frac{k^2}{a^2}}\delta \varphi _i+{\displaystyle \underset{j=1}{\overset{n}{}}}V,_{ij}\delta \varphi _j]`$ (9)
$`=`$ $`4\dot{\mathrm{\Phi }}\dot{\varphi }_i2V,_i\mathrm{\Phi },`$ (10)
where $`k`$ denotes the comoving wavenumber, $`V,_i`$ indicates the derivative of $`V`$ with respect to $`\varphi _i`$, and $`G`$ is Newton’s constant.
The Sasaki-Mukhanov variables for the $`n`$ matter fields are
$$Q_i=\delta \varphi _i+\frac{\dot{\varphi }_i}{H}\mathrm{\Phi }$$
(11)
and satisfy the following system of equations:
$$\ddot{Q}_i+3H\dot{Q}_i+\frac{k^2}{a^2}Q_i+\underset{j=1}{\overset{n}{}}[V_{,ij}\frac{8\pi G}{a^3}(\frac{a^3}{H}\dot{\varphi }_i\dot{\varphi }_j)^.]Q_j=0.$$
(12)
We will now specialize to our two field model. If the homogeneous part of the second scalar field $`\chi `$ vanishes (which will not be true if parametric resonance is to excite gravitational fluctuations - see later), then the inflaton $`\varphi `$ during the initial stages of reheating (when back-reaction effects are negligible) oscillates as follows
$$\varphi (\eta )=a^1\varphi _0cn(xx_0,\frac{1}{\sqrt{2}}),$$
(13)
where $`\varphi _0`$ is the amplitude of $`\varphi `$ at the end of the slow-rolling period, $`cn`$ is the Jacobi elliptic cosine function, $`\eta `$ is conformal time and $`x=\sqrt{\lambda }\varphi _0`$ is a rescaled dimensionless conformal time coordinate. Following , it is insightful to rescale all the fields $`f`$ by the scale factor and use new fields $`\stackrel{~}{f}=af`$. The equations of motion (12) then become
$`\stackrel{~}{Q_\varphi }^{^{\prime \prime }}`$ $`+`$ $`\left[\kappa ^2{\displaystyle \frac{a^{^{\prime \prime }}}{a}}+3cn^2(xx_0,{\displaystyle \frac{1}{\sqrt{2}}})\right]\stackrel{~}{Q_\varphi }`$ (14)
$`=`$ $`2{\displaystyle \frac{g^2}{\lambda \varphi _o}}cn(xx_0,{\displaystyle \frac{1}{\sqrt{2}}})\stackrel{~}{\chi }\stackrel{~}{Q_\chi }+M_{\varphi \varphi }\stackrel{~}{Q_\varphi }+M_{\varphi \chi }\stackrel{~}{Q_\chi }`$ (15)
$`\stackrel{~}{Q_\chi }^{^{\prime \prime }}`$ $`+`$ $`\left[\kappa ^2{\displaystyle \frac{a^{^{\prime \prime }}}{a}}+{\displaystyle \frac{g^2}{\lambda }}cn^2(xx_0,{\displaystyle \frac{1}{\sqrt{2}}})\right]\stackrel{~}{Q_\chi }`$ (16)
$`=`$ $`2{\displaystyle \frac{g^2}{\lambda \varphi _o}}cn(xx_0,{\displaystyle \frac{1}{\sqrt{2}}})\stackrel{~}{\chi }\stackrel{~}{Q_\varphi }+M_{\varphi \chi }\stackrel{~}{Q_\varphi }+M_{\chi \chi }\stackrel{~}{Q_\chi }`$ (17)
where, following the notation of (up to a prefactor containing $`a`$), we have have used the abbreviation
$$M_{\varphi _1\varphi _2}=\frac{8\pi G}{a^2}\left[\frac{1}{aH}(\stackrel{~}{\varphi _1}^{^{}}\frac{a^{^{}}}{a}\stackrel{~}{\varphi _1})(\stackrel{~}{\varphi _2}^{^{}}\frac{a^{^{}}}{a}\stackrel{~}{\varphi _2})\right]^{},$$
(18)
and where
$$\kappa ^2=\frac{k^2}{\lambda \varphi _0^2}.$$
(19)
In the model considered, the time-averaged equation of state is that of radiation and hence the scale factor is linear in $`\eta `$ and thus $`a^{^{\prime \prime }}=0`$.
Neglecting for a moment the terms on the right hand side of the equations, both equations (14) and (16) are of the form of Lamé equations. Lamé equations in the context of reheating were first noticed in and then studied in detail in , and . The coefficient in front of the $`cn^2`$ term (which in the case of (16) is $`g^2/\lambda `$) is crucial to the resonance structure of the equation. Since we are interested in super-Hubble modes, we need to know for which values of the coefficient the mode $`\kappa =0`$ lies in the first instability band of the equation. This is the case for values <sup>*</sup><sup>*</sup>*Note that this condition and the corresponding conditions in all the other examples discussed in this paper are stable against perturbative coupling constant renormalizations, the reason being that we consider $`g^2\lambda `$ and $`\lambda `$ should be constrained by the CMB anisotropy results to be a small number $`<<1`$. Thus, the perturbative correction terms which are of order $`g^4`$ are much smaller than either $`g^2`$ or $`\lambda `$.
$$1<\frac{g^2}{\lambda }<\mathrm{\hspace{0.17em}3}.$$
(20)
This implies that resonance occurs only in the equation (16) for $`\stackrel{~}{Q_\chi }`$, and then only for values of $`g^2/\lambda `$ in the above range (or in the range corresponding to other instability bands). However, if the condition (20) is satisfied, and if the $`\chi `$ field is indeed excited, then parametric amplification of $`\stackrel{~}{Q_\chi }`$ is expected, and via the terms on the right hand side of (14), induced exponential growth of $`\stackrel{~}{Q_\varphi }`$ should occur.
Note, however, that it is not sufficient to show that an amplification of $`\stackrel{~}{Q_\varphi }`$ or $`\stackrel{~}{Q_\chi }`$ occurs in order to demonstrate that parametric resonance during reheating will have a crucial effect on the amplitude of gravitational fluctuations. Note that also in the conventional treatment of fluctuations the amplitude of the $`\stackrel{~}{Q}`$ variable grows during the interval in which the equation of state of the background changes. As emphasized in , a straight forward way to check if the effect discussed here is a new effect is to consider the time evolution of the “traditional conserved quantity” $`\zeta `$ () which gives a measure of the adiabatic component of the metric fluctuations. In the multi-field case, $`\zeta `$ is given by
$$\zeta =\frac{H}{_j\dot{\varphi }_{j}^{}{}_{}{}^{2}}\underset{i}{}\dot{\varphi }_iQ_i.$$
(21)
In the two field case, the evolution equation for $`\zeta `$ is
$`\dot{\zeta }`$ $`=`$ $`{\displaystyle \frac{H}{\dot{H}}}{\displaystyle \frac{^2}{a^2}}\mathrm{\Phi }`$ (22)
$`+`$ $`{\displaystyle \frac{H}{2}}\left[{\displaystyle \frac{Q_\varphi }{\dot{\varphi }}}{\displaystyle \frac{Q_\chi }{\dot{\chi }}}\right]{\displaystyle \frac{d}{dt}}\left({\displaystyle \frac{\dot{\varphi }^2\dot{\chi }^2}{\dot{\varphi }^2+\dot{\chi }^2}}\right).`$ (23)
In the case of a single scalar field, the second term vanishes, but the equation cannot be applied during reheating since $`\dot{H}=0`$ at times when $`\dot{\varphi }=0`$.
In the single field case, the evolution equation for $`\zeta `$ which applies also during reheating is (for long-wavelength perturbations for which the spatial gradient term can be neglected)
$$(1+w)\dot{\zeta }=0.$$
(24)
This shows that $`\zeta `$ is conserved unless $`w=1`$. Unless matter is given by an oscillating scalar field (in which case $`w=1`$ will occur at the turnaround points when $`\dot{\varphi }=0`$), Eq. (24) implies that the variable $`\zeta `$ is conserved on scales outside the Hubble radius . However, reheating corresponds to an oscillating inflaton field, in which case the conclusion that $`\zeta `$ is constant may break down, as discussed in . Nevertheless, in the specific single field models which have been analyzed it was found that no net increase of $`\zeta `$ occurs during the initial stages of reheating, and a general proof of the constancy of $`\zeta `$ in single field models was suggested in .
If matter is described, more realistically, in terms of multiple scalar fields (each of which is given by a conventional action), then it appears extremely unlikely that $`w=1`$ will occur at all, since at the points in time when $`\dot{\varphi }=0`$, the other fields will not all also be at rest, and thus the net value of $`w`$ will be greater than $`1`$. Therefore, the only realistic possibility for growth of $`\zeta `$ is as a consequence of the second term on the right hand side of (22), a term which corresponds to an isocurvature perturbation. Inspection of (22), however, immediately shows that during reheating there is the chance of having a very large increase in $`\zeta `$ as a consequence of the zeros in $`\dot{\varphi }`$ which arise periodically in time. This effect is missed if the scalar fields are treated in the slow-roll approximation, or if the change in the equation of state during reheating is modelled as a monotonic change from a nearly de-Sitter equation of state during inflation to a radiative equation of state after reheating.
However, to see if there is indeed an exponential growth of isocurvature perturbations, it is important to take a closer look at the equations. In the specific two field model of (3), symmetric initial conditions for the homogeneous part of $`\chi `$ give $`\chi =0`$ and $`\dot{\chi }=0`$. In this case, it follows from (21) that $`\zeta `$ only depends on $`Q_\varphi `$, and, since by (14) there is no parametric amplification of super-Hubble modes of $`Q_\varphi `$ given that the coupling to $`Q_\chi `$ vanishes, that there will therefore be no parametric amplification of $`\zeta `$. The same result can also be seen from (22) since for symmetric initial conditions, the time derivative on the right hand side of the equation acts on a constant. The result is confirmed by our numerical analysis (see Figure 1).
However, due to quantum fluctuations we expect that the average of $`\chi `$ over a volume corresponding to a particular super-Hubble (but sub-horizon) mode will not vanish. It is reasonable to use for the homogeneous value of $`\chi `$ the r.m.s. value of the renormalized quantum fluctuations. It follows from (16) that $`Q_\chi `$ will experience parametric amplification during the initial stages of reheating. It will grow as $`exp[\mu _0\eta ]`$, where $`\mu _0`$ is the Floquet index of $`k=0`$ (which for continuity reasons cannot be too different from the Floquet exponent for long wavelengths). Via the non-vanishing source terms in (14), this will induce parametric growth of $`Q_\varphi `$, and this quantity mainly will contribute to the parametric growth of $`\zeta `$. Our numerical analysis confirms the above considerations. In Figure 2 we depict the growth of $`\zeta `$ during the initial stages of reheating for two different values of the homogeneous component of $`\chi `$, which shows how the onset of the parametric growth of $`\zeta `$ is dependent on the value of $`\chi `$. In Figure 3 we show how the growth of $`\zeta `$ is similar to the growth of $`Q_\varphi `$.
## III Isocurvature Perturbations
Having determined in the previous section that in the model (3) there is indeed amplification of $`\zeta `$ during reheating, we must now show that this effect is indeed a consequence of parametric resonance, and not just an effect due to the change in the equation of state between the inflationary era and the post-inflationary era, an effect which is already taken care of in the “usual” theory of isocurvature perturbations in inflationary cosmology, which we define as the results obtained when the transition in the equation of state between the inflationary slow-rolling phase ($`p\rho `$) and the post-inflationary radiation-dominated phase ($`p=\frac{1}{3}\rho `$) is taken to be monotonic. We will show that in the “usual” analysis there can be no exponential increase in the isocurvature perturbation, and that therefore the exponential increase we obtain here is a result of parametric resonance.
The fact that isocurvature perturbations can induce an adiabatic component on super-Hubble scales has been known for a long time . Entropy perturbations act as a source for $`\zeta `$ even on scales larger than the Hubble radius. This is true even in the case when matter is given by a single scalar field. In this case, the evolution equation for $`\zeta `$ becomes
$$\dot{\zeta }=\mathrm{\hspace{0.17em}3}H\left(\frac{\dot{p}}{\dot{\rho }}\frac{\delta p}{\delta \rho }\right).$$
(25)
These perturbations, however, are suppressed on scales larger than the Hubble radius.
In models with two or more scalar fields, the equation for $`\zeta `$ is given by (22), and it is thus clear that even on super-Hubble scales one should expect $`\dot{\zeta }0`$. In the approximation in which both fields are slowly rolling, the time evolution of $`\zeta `$ on scales larger than the Hubble radius was studied in detail in , with particular emphasis on calculating the deviations from scale-invariance of the resulting power spectrum of density fluctuations. However, since the analyses made use of the slow-rolling approximation, no effects of the dynamics of reheating were considered.
More recently, Taruya and Nambu and Bassett et al. considered the effect of reheating on the spectrum of density fluctuations and discovered a large growth of $`\zeta `$ due to the initial isocurvature perturbations, however in a model in which the necessary initial $`\chi `$ field fluctuations are exponentially suppressed during inflation. Bassett and Viniegra then pointed out that the suppression would be absent in the model (3).
It has been known for a long time that isocurvature perturbations can be produced in inflationary models with more than one scalar field. This issue was initially considered in the context of axion perturbations in , extended to more general two field models in and studied in detail in . It was discovered that initial super-Hubble-scale isocurvature perturbations induce an adiabatic component by the time that the scales re-enter the Hubble radius.
The gauge invariant expression for the total isocurvature perturbation in a multi-fluid system is
$$p\mathrm{\Gamma }\underset{i}{}(\delta p_i^{\mathrm{gi}}c_s^2\delta \rho _i^{\mathrm{gi}}),$$
(26)
where $`\delta p_i^{\mathrm{gi}}`$ and $`\delta \rho _i^{\mathrm{gi}}`$ are the gauge invariant pressure and density perturbations with respect to the total matter rest frame and the total speed of sound $`c_s^2`$ is defined as the weighted sum of the i-th speed of sound :
$$c_s^2\frac{\dot{p}}{\dot{\rho }}=\frac{1}{_j\dot{\varphi }_j^2}\underset{i}{}c_{si}^2\dot{\varphi }_i^2$$
(27)
with
$$c_{si}^2=1+2\frac{V,_i}{3H\dot{\varphi }_i}$$
(28)
The total isocurvature perturbation can be written as the sum of the non-adiabatic pressure component of the single component and of the relative isocurvature perturbation $`S_{ij}`$ as:
$$p\mathrm{\Gamma }=\underset{i}{}(\delta p_ic_{si}^2\delta \rho _i)+\frac{1}{_i\dot{\varphi }_i^2}\underset{i,j}{}\frac{\dot{\varphi }_i^2\dot{\varphi }_j^2}{2}S_{ij}(c_{si}^2c_{sj}^2)$$
(29)
where
$$S_{ij}\frac{\delta \rho _i^{\mathrm{gi}}}{\dot{\varphi }_i^2}\frac{\delta \rho _j^{\mathrm{gi}}}{\dot{\varphi }_j^2}.$$
(30)
The relative isocurvature perturbation $`S_{ij}`$ with respect to the total matter frame can be written for our two field model as
$`S_{\varphi \chi }`$ $`=`$ $`{\displaystyle \frac{\delta \rho _\varphi }{\dot{\varphi }^2}}{\displaystyle \frac{\delta \rho _\chi }{\dot{\chi }^2}}{\displaystyle \frac{av}{k}}Q_\varphi \left({\displaystyle \frac{1}{\dot{\varphi }^2}}+{\displaystyle \frac{1}{\dot{\chi }^2}}\right)`$ (31)
$`=`$ $`{\displaystyle \frac{\delta \rho _\varphi }{\dot{\varphi }^2}}{\displaystyle \frac{\delta \rho _\chi }{\dot{\chi }^2}}{\displaystyle \frac{3g^2}{8\pi G}}{\displaystyle \frac{\varphi \chi }{\dot{\varphi }^2\dot{\chi }^2}}(\dot{\varphi }\chi \dot{\chi }\varphi ){\displaystyle \frac{\dot{H}}{H}}(\zeta \mathrm{\Phi }),`$ (32)
where $`v`$ is the total perturbed velocity for matter, $`Q_\varphi `$ is the homogeneous energy transfer to the $`\varphi `$ component ($`Q_\varphi +Q_\chi =0`$ because of total energy conservation), and now $`\delta \rho _i`$ are the density perturbations in the longitudinal gauge:
$$\delta \rho _i=\dot{\varphi }_i\delta \dot{\varphi }_i\mathrm{\Phi }\dot{\varphi }_i^2+V,_i\delta \varphi _i.$$
(33)
From Equation (32) it follows that the parametric resonance from the matter sector of the theory (the $`Q`$ variables to be specific) induces exponential growth of the relative isocurvature perturbation, and hence also of the total isocurvature perturbation (see Figure 4).
In turn, isocurvature perturbations determine the change in $`\zeta `$ via Eq. (22). This shows that in the presence of scalar field interaction terms, there is a correlated exponential growth of $`\zeta `$ and of the relative isocurvature perturbation $`S_{\varphi \chi }`$. This exponential growth is a consequence of parametric resonance and is absent if the phase transition is modelled with a monotonically increasing value of $`w`$.
Analogously to Eq. (21) for the Bardeen parameter, the total non adiabatic pressure $`p\mathrm{\Gamma }`$ can be expressed in terms of the Sasaki-Mukhanov variables in the following way:
$$p\mathrm{\Gamma }=\underset{i}{}[\frac{\dot{V}}{\rho }\dot{\varphi }_iQ_i2V,_iQ_i2\frac{\dot{V}}{3H_i\dot{\varphi }_i^2}(\dot{\varphi }_i\dot{Q}_i+V,_iQ_i)].$$
(34)
As mentioned above, the exponential growth of $`S_{\varphi \chi }`$ and $`\zeta `$ during reheating is a new effect due entirely to parametric resonance. The growth of fluctuations in inflationary models with two uncoupled fields were studied in in an approximation in which the oscillations of the inflaton field were neglected. In this case there is no growth of $`S_{\varphi \chi }`$. An initial isocurvature perturbation does induce the growth of an adiabatic component on super-Hubble scales, but the final amplitude of the adiabatic mode is not much larger than the initial amplitude of the isocurvature perturbation, in agreement with the earlier analysis in .
## IV Three Golden Rules
Based on the analysis of Section 2, it appears that several conditions are required in order to have efficient parametric resonance of super-Hubble-scale metric fluctuations.
1. In the absence of gravitational perturbations there must be broad-band parametric resonance in the matter sector of the theory corresponding to isocurvature fluctuations, and $`k=0`$ must be part of the resonance band.
2. The fluctuations in the matter field which undergoes parametric resonance must be effectively massless during inflation. More precisely, there should be no large net suppression of these fluctuations before the phase of parametric resonance.
3. The homogeneous value of the matter field which undergoes resonance must be non-vanishing. This is the weakest of the three conditions since it is only required if we work strictly to first order in perturbations.
We show how these three rules are satisfied in another model with massless fields, but now based on negative coupling instability . The potential is the following:
$$V(\varphi ,\chi )=\frac{1}{4}\lambda \varphi ^4\frac{1}{2}g^2\varphi ^2\chi ^2+\frac{1}{4}\lambda _\chi \chi ^4$$
(35)
with the parameter $`r\lambda \lambda _\chi /g^4>1`$ in order to have a potential bounded from below . This model has an attractor for $`\chi `$ in the point which minimize the potential for $`\chi `$
$$\overline{\chi }(t)\frac{g}{\sqrt{\lambda _\chi }}\varphi (t).$$
(36)
In this way even the third and weakest of the above conditions is satisfied. The second one is easily satisfied because of the negative effective mass for the $`\chi `$ fluctuations when the background $`\chi `$ is small. In order to find the unstable bands one can use the attractor solution and estimate the frequency of the inflaton $`\varphi `$ and of the fluctuations $`\delta \chi `$ during the period of coherent oscillations:
$`\omega _\varphi ^2`$ $`=`$ $`\lambda \varphi ^2g^2\chi ^2\lambda \varphi ^2(1{\displaystyle \frac{1}{r}})\stackrel{~}{\lambda }\varphi ^2`$ (37)
$`\omega _{\delta \chi }^2`$ $`=`$ $`{\displaystyle \frac{k^2}{a^2}}+3\lambda _\chi \chi ^2g^2\varphi ^2{\displaystyle \frac{k^2}{a^2}}+2g^2\varphi ^2.`$ (38)
If $`\chi `$ is small compared to the inflaton $`\varphi `$, then an unstable band for $`k=0`$ should be located at $`2g^2=2\stackrel{~}{\lambda }`$. This gives a second order equation for $`g^2`$ whose positive root is:
$$g^2=\lambda _\chi \frac{1+\sqrt{1+4\lambda /\lambda _\chi }}{2}$$
(39)
We confirm numerically this analytical estimate in Figure 5 for two allowed values of $`g^2`$: $`g^2=\lambda `$ with $`\lambda \lambda _\chi `$ and $`g^2=\lambda /2(\sqrt{5}1)`$ with $`\lambda =\lambda _\chi `$. The reason for the growth of $`\zeta `$ is similar to the previous case: $`Q_\chi `$ is parametrically amplified, feeds the growth of $`Q_\varphi `$ and in this case both contribute to the growth of $`\zeta `$.
## V A Model Motivated by Hybrid Inflation
Another natural scenario in which the above conditions can all be satisfied is hybrid inflation . Hybrid inflation is also an attractive framework for implementing inflation in the context of supergravity models . Since (at least) two fields are involved in the dynamics of hybrid inflation, the generation of isocurvature perturbations is rather natural. In hybrid models, the phase of inflation during which the inflaton field $`\varphi `$ is slowly rolling towards $`\varphi =0`$ is terminated by a phase transition in the second scalar field $`\chi `$, a field with the double-well potential. This implies that during the oscillations of $`\varphi `$, the background value of $`\chi `$ is non-vanishing, leading to an obvious realization of condition (3) above.
Parametric resonance in the matter sector of hybrid inflation models was studied in detail by Garcia-Bellido and Linde , and, in supersymmetric hybrid inflation, by Bastero-Gil et al. . The resonance of the fluctuations of the two fields $`\varphi `$ and $`\chi `$ is inefficient for a large set of the parameter space since both the $`\delta \varphi `$ and $`\delta \chi `$ fields are effectively massive during the regime of coherent oscillations ($`\delta \chi `$ becomes massive through the Higgs mechanism). Quite generically, parametric resonance could be much more efficient if there is a third field $`\psi `$ which couples to both $`\varphi `$ and $`\chi `$. Another interesting possibility is to consider a doublet for the field $`\chi `$. Then, even in the phase of coherent oscillations, there is a massless degree of freedom, namely the “Goldstone” mode. Such a situation arises naturally in supergravity models .
Therefore, as a toy model we will consider the following potential for the inflaton field $`\varphi `$ and the doublet $`\chi =\frac{1}{\sqrt{2}}(\chi _1,\chi _2)`$:
$`V`$ $`=`$ $`\lambda ({\displaystyle \frac{M^2}{2\lambda }}|\chi |^2)^2+{\displaystyle \frac{1}{2}}m^2\varphi ^2+g^2\varphi ^2|\chi |^2`$ (40)
$`=`$ $`{\displaystyle \frac{1}{2}}m^2\varphi ^2+{\displaystyle \frac{1}{2}}g^2\varphi ^2(\chi _1^2+\chi _2^2)+{\displaystyle \frac{M^4}{4\lambda }}`$ (41)
$``$ $`{\displaystyle \frac{1}{2}}M^2(\chi _1^2+\chi _2^2)+{\displaystyle \frac{\lambda }{2}}\chi _1^2\chi _2^2`$ (42)
$`+`$ $`{\displaystyle \frac{\lambda }{4}}(\chi _1^4+\chi _2^4).`$ (43)
For supersymmetric hybrid inflation, there is only one independent coupling constant since
$$\lambda =\frac{g^2}{2}.$$
(44)
The values of the masses and coupling constants are constrained by the amplitude of density fluctuations at Hubble radius crossing, which is given by
$$\frac{\delta \rho }{\rho }\frac{g}{\lambda ^{3/2}}\left(\frac{M}{M_{pl}}\right)^3\left(\frac{M}{m}\right)^2$$
(45)
which should be about $`10^5`$. In our numerical simulations, we choose $`g^2=10^3`$, $`\lambda =g^2/2`$ (as in a supersymmetric model), $`M^2/M_{pl}^2=10^{12}`$ and $`M^2/m^2=10^{10}`$. With these values, the Hubble parameter during the late stages of inflation is much larger than $`m`$ which ensures slow rolling of $`\varphi `$.
In this model, slow-roll inflation takes place while the value of the inflaton $`\varphi `$ is larger than $`\varphi _c=M/g`$. For these values, the effective square mass $`m_\chi ^2`$ of $`\chi `$ (evaluated at $`\chi =0`$) is positive. Once $`\varphi `$ drops below $`\varphi _c`$, the $`m_\chi ^2`$ turns negative and inflation ends (at time $`t_c`$) via a symmetry breaking transition in the matter sector. We choose the basis of the $`\chi `$ fields such that the order parameter of the transition is $`\chi _1`$. However, since in general the initial ratio of $`\chi `$ fields and $`\chi `$ velocities at $`t_c`$ is not the same
$$\frac{\dot{\chi }_2}{\dot{\chi }_1}(t_c)\frac{\chi _2}{\chi _1}(t_c),$$
(46)
we can with no loss of generality assume that the value of $`\chi _2`$ at the time $`t_r`$, when the $`\chi `$ transition is complete and the $`\varphi `$ oscillations start, does not vanish (as in the previous case a reasonable value for $`\chi _2`$ is the r.m..s. of the renormalized quantum fluctuations). Thus, we have argued that the third of the conditions mentioned at the beginning of this section (non-vanishing background matter fields) is naturally satisfied in this model (in contrast to the model considered in Section 2). The initial values of the matter fields at the beginning of the period of oscillation will be
$`\chi _1(t_r)`$ $``$ $`M/\sqrt{\lambda }`$ (47)
$`0\chi _2(t_r)`$ $``$ $`M/\sqrt{\lambda }.`$ (48)
In contrast to the model (1) in which the matter fluctuations are exponentially suppressed during inflation, the effective negative coupling instability in the matter fields in the time interval between $`t_c`$ and $`t_r`$ leads to the conclusion that the $`\chi `$ fluctuations are not suppressed. In fact, they are suppressed during slow-roll inflation, but build up again exponentially fast in the time interval $`t_c<t<t_r`$. To see this, we focus on the evolution of the fluctuations in of the $`\chi `$ field since, as we shall see below, these are essential for the effectiveness of parametric resonance. We will consider field fluctuations neglecting metric perturbations and the mixing terms deriving from particle interactions (even if in the hybrid models these mixing terms are not perturbatively small and their importance has been emphasized in ). Under these approximations the evolution equation for $`\delta \chi _2`$ is
$$\ddot{\delta }\chi _2+3H\dot{\delta }\chi _2=\left(\frac{k^2}{a^2}+g^2\varphi ^2+\lambda \chi _1^2+3\lambda \chi _2^2M^2\right)\delta \chi _2$$
(49)
For $`\varphi >\varphi _c`$, the effective square mass is larger than $`H^2`$ and positive, thus leading to damped oscillatory solutions
$$\delta \chi _2a^{3/2}(t)exp(i\omega t),$$
(50)
with $`\omega =g\varphi `$ (in the adiabatic limit). During this time interval, however, the homogeneous components of the matter fields are also damped. The evolution equation for the order parameter $`\chi _1`$ is
$$\ddot{\chi }_1+3H\dot{\chi }_1=\left(g^2\varphi ^2+\lambda \chi _1^2+\lambda \chi _2^2M^2\right)\chi _1$$
(51)
Since (up to the contributions from $`\chi _2^2`$ which are negligible during inflation) the effective masses in (51) and (49) are the same, the damping rates of $`\chi _1`$ and $`\delta \chi _2`$ are also the same for $`t<t_c`$. In the time interval between $`t_c`$ and $`t_r`$, the signs of the effective square masses in both Equation (51) for $`\chi _1`$ and Equation (49) for $`\delta \chi _2`$ are reversed. In both cases, the effective $`m^2`$ is now $`M^2`$, leading to exponential increase in both $`\chi _1`$ and $`\delta \chi _2`$. This period ends when $`\chi _1`$ reaches the minimum of the potential at time $`t_r`$. To summarize the above discussion, the evolution of $`\chi _1`$ follows
$`{\displaystyle \frac{M}{\sqrt{\lambda }}}\chi _1(t_r)`$ $``$ $`e^{M(t_rt_c)}\chi _1(t_c)`$ (52)
$``$ $`e^{M(t_rt_c)}e^{\frac{3}{2}H(t_ct_i)}\chi _1(t_i)`$ (53)
where $`t_i`$ is the time at the beginning of inflation and $`H`$ is the Hubble constant during inflation, assumed to be constant to make the equation simple (this assumption does not affect the basic point we are making). In comparison, the evolution of $`\delta \chi _2`$ obeys
$`\delta \chi _2(t_r)`$ $``$ $`e^{M(t_rt_c)}\delta \chi _2(t_c)`$ (54)
$``$ $`e^{M(t_rt_c)}e^{\frac{3}{2}H(t_ct_i)}\delta \chi _2(t_i)`$ (55)
which shows that the exponential growth of $`\delta \chi _2`$ for $`t_c<t<t_r`$ precisely makes up for the exponential decay during the period $`t_1<t<t_c`$, exactly as it does for the evolution of $`\chi _1`$. Equations (52) and (54) can be combined to give
$$\delta \chi _2(t_r)\frac{M/\sqrt{\lambda }}{\chi _1(t_i)}\delta \chi _2(t_i).$$
(56)
This demonstrates that there is no overall suppression of the fluctuations in $`\delta \chi _2`$ before the onset of parametric resonance, showing that the second condition for the effectiveness of parametric amplification of super-Hubble gravitational modes mentioned at the beginning of this section is satisfied.
The final conditions to discuss are the criteria for parametric resonance of the $`k=0`$ modes of the matter perturbations. To do this, we consider the mode equation for $`\delta \chi _2`$ during the period in which the inflaton $`\varphi `$ oscillates. For general hybrid models, the background dynamics is chaotic since both $`\varphi `$ and $`\chi _1`$ oscillate with different frequencies. However, in the supersymmetric case , the frequencies coincide and the background dynamics becomes non-chaotic. Both $`\varphi `$ and $`\chi _1`$ oscillate with the frequency $`\sqrt{2}M`$. To simplify the analysis, we shall neglect the back-reaction of particle production and expansion on the inflaton, and neglect the Hubble damping term in the equation of motion (this is a good approximation since we are considering a case in which $`H<<\sqrt{2}M`$ and the fields oscillation are not damped by the expansion of the universe). Therefore, we take the inflaton to oscillate with amplitude $`\varphi _a<\varphi _c`$, and $`\chi _1`$ will oscillate about its ground state as
$$\chi _1(z)=\frac{M}{\sqrt{\lambda }}(1+f(z)),$$
(57)
where $`f(z)`$ is periodic with period $`2\pi `$. It is convenient to introduce the dimensionless time $`z=\sqrt{2}Mt`$. Denoting the derivative with respect to $`z`$ by a prime, the equation for the Fourier mode $`\chi _{2k}`$ of $`\delta \chi _2`$ becomes
$$\chi _{2k}^{^{\prime \prime }}+\chi _{2k}\left(\frac{k^2}{2a^2M^2}+\frac{g^2\varphi _a^2}{4M^2}+\frac{g^2\varphi _a^2}{4M^2}cos(2z)+f+\frac{f^2}{2}\right)=\mathrm{\hspace{0.17em}0},$$
(58)
where we have neglected the terms in $`\chi _2`$. In the absence of the final term (the term containing $`f(z)`$), this has the form of the Mathieu equation
$$\chi _{2k}^{^{\prime \prime }}+\chi _{2k}\left(A(k)2qcos(2z)\right)=\mathrm{\hspace{0.17em}0}.$$
(59)
The value of $`q`$ is $`q1/8`$, the maximal value being taken on if $`\varphi _a=\varphi _c`$, and for long wavelengths $`A(k)2q`$. As can be seen from the Floquet instability charts (see e.g. Fig. 1 in ), these values do not correspond to efficient resonance. From the evolution of the background fields obtained from the full numerical solution of the background field equations (see Figure 6) it follows that the amplitude of oscillation $`\varphi _a`$ is in fact substantially smaller than $`\varphi _c`$. In contrast, $`\chi _1`$ oscillates with a large amplitude. Hence, the term containing $`f(z)`$ in Equation (58) is more important. This term leads to negative coupling instability (see for a general discussion of resonant particle production by negative coupling instability) for small values of $`k`$. Hence, we expect parametric amplification of long wavelength gravitational fluctuations in our model.
The above considerations are supported by our numerical results. In Figure 6 we show the evolution of the background fields $`\varphi `$, $`\chi _1`$, $`\chi _2`$ and $`H`$ as a function of time in a simulation with parameters mentioned after Equation (45), and with initial conditions $`\varphi _0=3M_{\mathrm{pl}}`$, $`\chi _1=.01M_{\mathrm{pl}}`$, $`\chi _2=.0001M_{\mathrm{pl}}`$.
As is evident, following an initial transient period the three scalar fields oscillate with the same frequency. The results for the fluctuation variables $`Q_\varphi `$, $`Q_{\chi _1}`$, $`Q_{\chi _2}`$ and $`\zeta `$ are shown in Figure 7. The initial perturbation amplitudes were chosen to be $`Q_\varphi (t_0)=1`$, $`Q_{\chi _1}(t_0)=Q_{\chi _1}(t_0)=10^4`$, and all their derivatives set to zero for a wavelength outside the Hubble radius ($`k=0`$). As is evident, after the initial transient period, all four quantities grow almost with the same Floquet exponent, as expected from our analytical analysis.
At this point, an obvious question is whether the field $`\chi _2`$ is essential in order to obtain parametric resonance of super-Hubble-scale cosmological fluctuations. In fact, the field fluctuations $`\delta \chi _1`$ also will experience an effective negative coupling instability , and therefore the presence of $`\chi _2`$ is not essential for this supersymmetric choice of the parameters. The equation of motion for $`\delta \chi _1`$ for a hybrid model with two field is
$$\ddot{\delta }\chi _1+3H\dot{\delta }\chi _1=\left(\frac{k^2}{a^2}+g^2\varphi ^2+3\lambda \chi _1^2M^2\right)\delta \chi _1.$$
(60)
The effective squared mass is large and positive during slow-rolling. At the beginning of the transient period (when $`\chi _1`$ starts rolling down its potential but is not yet close to the minimum of the potential) the effective squared mass turns negative. Once $`\chi _1`$ gets close to its equilibrium position, the effective squared mass will again be large and positive (the factor $`3`$ in the third term on the r.h.s. of (60) is crucial). However, since $`\chi _1`$ is oscillating with a large amplitude, the effect of the large mass will be periodically cancelled out by these oscillations. Neglecting the expansion of the background, Eq. (60) can be written as
$$\chi _{1k}^{^{\prime \prime }}+\chi _{1k}\left(\frac{k^2}{2a^2M^2}+\frac{g^2\varphi ^2}{2M^2}+1+3f(z)+\frac{3}{2}f^2\right)=\mathrm{\hspace{0.17em}0}.$$
(61)
From our numerical results (Figure 6) we expect that the amplitude of $`f(z)`$ will be only slightly smaller than 1. Hence, we expect negative coupling instability for long wavelength metric perturbations also in the two field case . However, as demonstrated below for values of the coupling constants which do not correspond to the supersymmetric point, the Floquet index in the two field case will be smaller than in the three field model.
Other interesting effects happen if we go away from the supersymmetric point $`g^2=2\lambda `$. In analogy with the results of which show that random noise in the inflaton leads to an increase in the strength of the parametric instability, we expect that the chaotic background dynamics will not eliminate but rather strengthen the resonance. Chaotic background dynamics is expected for $`g^2\lambda `$ in hybrid models . This effect is shown by our numerical simulations (Figure 8 and Figure 9) which show that the parametric resonance of super-Hubble-scale gravitational fluctuations for the choice $`g^2=\lambda `$ is larger than in the supersymmetric case, where no chaoticity is present . Figure 10 shows how the presence of the “Goldstone” mode $`\chi _2`$ changes the development of the resonance in this chaotic case. In the two field case the Floquet index with which $`\zeta `$ grows is smaller than the corresponding index in the three field case.
## VI Discussion
We have studied the parametric amplification of long wavelength gravitational fluctuations during reheating in two field inflationary Universe models. We have partially confirmed the results of Bassett and Viniegra and shown that this effect is possible for certain models. We have established criteria under which an exponential increase in the amplitude of cosmological perturbations during the period when the inflaton field oscillates should be expected. It is crucial that there must be either broad-band parametric instability or negative coupling instability in the matter sector of the theory (i.e., in the absence of gravitational perturbations). This will excite isocurvature fluctuations during reheating. It is important that there be no net exponential damping in the amplitude of the isocurvature fluctuations before reheating. The resonance in the matter sector then induces a resonance in the gravitational sector provided that the background values of the matter fields do not vanish. Since large coupling constants are not necessary in order to have efficient resonance, the effect is stable against perturbative coupling constant renormalizations. We have shown that in this case the resulting increase in the amplitude of the adiabatic mode, conveniently tracked in terms of the variable $`\zeta `$, and of isocurvature fluctuations, tracked in terms of the non adiabatic pressure $`p\mathrm{\Gamma }`$, is exponential and is due to the oscillations in the inflaton field. This means that the effect is absent if the phase transition is modelled by a monotonic change in $`w=p/\rho `$.
We then argue that the conditions under which parametric amplification of long wavelength gravitational fluctuations occurs are naturally satisfied in a class of models of hybrid inflation. The presence of a complex matter scalar field enhances the resonance, since it ensures the existence of a field which is massless in the true vacuum of the theory, but it is not crucial if there is negative coupling instability. However, note that the existence of massless modes is helpful for the effect to occur. Such massless modes arise quite generically in string theory (see e.g. for a recent review). Thus, the parametric amplification of long wavelength fluctuations may be also present in models of inflation based on string theory.
Acknowledgements
We are grateful to Bruce Bassett, Serguei Khlebnikov, Lev Kofman and Bill Unruh for stimulating discussions, and Jim Zibin for comments on the draft. R.B. wishes to thank Bill Unruh for hospitality at the University of British Columbia where this work was completed. F. F. wishes to thank Brown University for hospitality. The research was supported in part (at Purdue) by the U.S. Department of Energy under Contract DE-FG02-91ER40681, TASK B, and (at Brown) by DE-FG02-91ER40688, TASK A.
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# A note on asymptotically isometric copies of 𝑙¹ and 𝑐₀
## Abstract.
Nonreflexive Banach spaces that are complemented in their bidual by an L-projection - like preduals of von Neumann algebras or the Hardy space $`H^1`$ \- contain, roughly speaking, many copies of $`l^1`$ which are very close to isometric copies. Such $`l^1`$-copies are known to fail the fixed point property. Similar dual results hold for $`c_0`$.
###### Key words and phrases:
asymptotically isometric copies of $`l_1`$, James’ distortion, L-summands, L-embedded, M-ideals, fixed point property
In it is shown that an isomorphic $`l^1`$-copy does not necessarily contain asymptotically isometric $`l^1`$-copies although by James’ classical distortion theorem it always contains almost isomorphic $`l^1`$-copies. (For definitions see below.) Within this context and the context of the fixed point property Dowling and Lennard show that the presence of an asymptotic $`l^1`$-copy makes a Banach space fail the fixed point property. Then they prove that every nonreflexive subspace of $`L^1[0,1]`$ fails the fixed point property by observing that the proof of a theorem of Kadec and Pełczyński \[10, Th. 6\] yields an asymptotic $`l^1`$-copy inside such supspaces of $`L^1[0,1]`$. Alspach’s example may be considered as an early forerunner of these results.
In the present note we modify a construction of Godefroy in order to show that every nonreflexive subspace of any L-embedded Banach space contains an asymptotic $`l^1`$-copy and thus, in particular, fails the fixed point property. Analogous results hold for $`c_0`$ and M-embedded spaces.
Let $`(x_n)`$ be a sequence of nonzero elements in a Banach space $`X`$.
We say that $`(x_n)`$ spans $`l^1`$ $`r`$-isomorphically or just isomorphically if there exists $`r>0`$ (trivially $`r1`$) such that $`r\left(_{n=1}^{\mathrm{}}|\alpha _n|\right)_{n=1}^{\mathrm{}}\alpha _nx_n_{n=1}^{\mathrm{}}|\alpha _n|`$ for all scalars $`\alpha _n`$. We say that $`(x_n)`$ spans $`l^1`$ almost isometrically if it is such that there exists a sequence $`(\delta _m)`$ in $`[0,1[`$ tending to $`0`$ such that $`(1\delta _m)\left(_{n=m}^{\mathrm{}}|\alpha _n|\right)_{n=m}^{\mathrm{}}\alpha _nx_n\left(_{n=m}^{\mathrm{}}|\alpha _n|\right)`$ for all $`m\mathrm{I}N`$. Trivially the property of spanning $`l^1`$ almost isometrically passes to subsequences. Analogously, we say that $`(x_n)`$ spans $`c_0`$ almost isometrically if there exists a sequence $`(\delta _m)`$ as above such that $`(1\delta _m)sup_{mnm^{}}|\alpha _n|_{n=m}^m^{}\alpha _nx_n(1+\delta _m)sup_{mnm^{}}|\alpha _n|`$ for all $`mm^{}`$.
Recall that James distortion theorem (or , ) for $`l^1`$ and $`c_0`$ says that every isomorphic copy of $`l^1`$ (of $`c_0`$) contains an almost isometric copy of $`l^1`$ (of $`c_0`$).
Following we say that $`(x_n)`$ spans $`l^1`$ asymptotically isometrically (or just that $`(x_n)`$ spans $`l^1`$ asymptotically) if there is a sequence $`(\delta _n)`$ in $`[0,1[`$ tending to $`0`$ such that $`_{n=1}^{\mathrm{}}(1\delta _n)|\alpha _n|_{n=1}^{\mathrm{}}\alpha _nx_n_{n=1}^{\mathrm{}}|\alpha _n|`$ for all scalars $`\alpha _n`$. A sequence $`(x_n)`$ is said to span $`c_0`$ asymptotically isometrically (or just to span $`c_0`$ asymptotically) if there is a sequence $`(\delta _n)`$ as above such that $`sup_{nm}(1\delta _n)|\alpha _n|_{n=1}^m\alpha _nx_nsup_{nm}(1+\delta _n)|\alpha _n|`$ for all $`m\mathrm{I}N`$. (Note that in the definition of asymptotic $`c_0`$’s and $`l^1`$’s in the sequence $`(\delta _n)`$ is supposed to be decreasing but that this difference is not essential.) Finally we say that a Banach space is isomorphic (respectively almost isometric, respectively asymptotically isometric) to $`l^1`$ (to $`c_0`$) if it has a basis with the corresponding property. Clearly a sequence spanning $`l^1`$ asymptotically spans $`l^1`$ almost isometrically. The main result of states that the converse does not hold because there are almost isometric copies of $`l^1`$ which do not contain $`l^1`$ asymptotically. Analogously, it was proved in \[4, Th. 3\] that there exist almost isometric copies of $`c_0`$ which do not contain asymptotic $`c_0`$-copies.
Some notation: The results are stated for complex scalars but hold also for real Banach spaces. Operator means linear bounded map. As usual, we consider a Banach space as a subspace of its bidual omitting the canonical embedding. $`[x_n]`$ denotes the complete linear hull of $`(x_n)`$, $`(e_n)`$ is the standard basis of $`l^1`$. Basic properties and definitions of Banach space theory can be found in or in .
A Banach space $`X`$ is said to have the fixed point property if any contractive (not necessarily linear) map $`f:CC`$ on any non-empty closed bounded convex subset $`CX`$ has a fixed point where contractive means that $`f(x)f(y)<xy`$ for all $`x,yX`$.
Let $`Y`$ be a subspace of a Banach space $`X`$ and $`P`$ be a projection on $`X`$. $`P`$ is called an L-projection provided $`x=Px+(\mathrm{id}_XP)x`$ for all $`xX`$. A subspace $`YX`$ is called an M-ideal in $`X`$ if its annihilator $`Y^{}`$ in $`X^{}`$ is the range of an L-projection on $`X^{}`$. $`Y`$ is called an L-summand in $`X`$ if it is the range of an L-projection on $`X`$. In the special case where $`X=Y^{\prime \prime }`$ and where $`Y`$ is an M-summand (respectively an L-summand) in $`Y^{\prime \prime }`$ we say that $`Y`$ is M-embedded (respectively L-embedded). As examples we only mention that preduals of von Neumann algebras, in particular $`l^1`$ and $`L^1`$-spaces, furthermore the Hardy space $`H_0^1`$ and the dual of the disc algebra are L-embedded. The sequence space $`c_0`$, the space of compact operators on a Hilbert space, and the quotient $`C/A`$ of the continuous functions on the unit circle by the disc algebra $`A`$ are examples among M-embedded spaces. The dual of an M-embedded space is L-embedded; the converse is false \[8, III.1.3\]. Throughout this note, if $`X`$ is an L-embedded Banach space we will write $`X_s`$ for the complement of (the canonical embedding of) $`X`$ in $`X^{\prime \prime }`$ that is $`X^{\prime \prime }=X_1X_s`$. In this case $`P`$ (or $`P_X`$ to avoid confusion) will denote the L-projection from $`X^{\prime \prime }`$ onto $`X`$. There is a useful criterion for L-embeddedness of subsapces of L-embedded spaces due to Li ( or \[8, Th. IV.1.2\]): A closed subspace $`Y`$ of an L-embedded Banach space $`X`$ is L-embedded if and only if $`PY^{}=Y`$. In this case if $`Y`$ is L-embedded and if one identifies $`Y^{\prime \prime }=Y_1Y_s`$ and $`Y^{}X^{\prime \prime }`$ then $`Y_s=Y^{}X_s`$. Since biduals of L-embedded spaces are quite ”big” and therefore difficult to handle we mention in passing that a theorem of Buhvalov-Lozanovskii (, \[8, IV.3.4\]) provides a characterisation of L-embeddedness of subspaces of $`L^1[0,1]`$ only in terms of the spaces themselves: Subspaces of $`L^1[0,1]`$ are L-embedded if and only if their unit balls are closed with respect to the measure topology \[8, IV.3.5\].
The standard reference for M- and L-embedded spaces is the monograph .
There are few stability results for L-embeddedness: Neither subspaces nor quotients inherit this property \[8, IV.1\]. Therefore the following lemma reveals a nice exception. In particular it underlines the idea that $`l^1`$-copies are the building blocks of L-embedded spaces.
###### Lemma 1.
Almost isometric copies of $`l^1`$ which are subspaces of L-embedded Banach spaces are L-embedded.
Proof: Consider first the L-embedded subspaces $`U_m=[e_n]_{nm}`$ of $`l^1`$, $`m\mathrm{I}N`$. We have $`U_m^{}=U_m_1(c_0U_m^{})`$. An element $`\mu (l^1)_s=c_0^{}`$ belongs to each $`U_m^{}`$. \[Denote by $`\rho _m`$ the projection $`(\alpha _n)(0,\mathrm{},0,\alpha _{m+1},\alpha _{m+2})`$ on $`l^1`$. Then $`\mu `$ annihilates $`\text{ker}\rho _m^{}`$ because $`\text{ker}\rho _m^{}c_0`$. Thus $`\mu \text{ran}(\rho ^{\prime \prime })`$ and $`\mu \overline{[e_n]}_{nm}^w^{}c_0^{}`$.\]
Let $`X`$ be an L-embedded Banach space with L-decomposition $`X^{\prime \prime }=X_1X_s`$ and with L-projection $`P`$ from $`X^{\prime \prime }`$ onto $`X`$. Let $`(x_n)`$ be a sequence spanning an almost isometric copy $`Y`$ of $`l^1`$ in $`X`$, put $`Y_m=[x_n]_{nm}`$, $`m\mathrm{I}N`$. Via the isomorphism between $`Y`$ and $`l^1`$ induced by $`x_ne_n`$ the situation described for $`l^1`$ carries over to $`Y`$. That is, there is $`ZY^{}X^{\prime \prime }`$ such that $`Y^{}=YZ`$, $`Y_m^{}=Y_m(ZY_m^{})`$ and $`zZY_m^{}`$ for all $`m\mathrm{I}N`$, $`zZ`$.
Since $`(x_n)`$ is supposed to span $`Y`$ almost isometrically, there are numbers $`\eta _m0`$ tending to $`0`$ such that $`y+z(1\eta _m)(y+z)`$ for all $`yY_m`$, $`zZY_m^{}`$.
Let $`zZ`$. In order to show that $`Y`$ is L-embedded it is enough to show that $`Pz=0`$ because then Li’s criterion $`PY^{}=Y`$ is fulfilled. But $`zZY_m^{}`$ and by a quantitative version of Li’s result (see \[8, IV.1.4\] or \[14, Lem. 2\]) applied to $`X`$ and $`Y_m`$, we have $`Pz3\eta _m^{1/2}z`$ for all $`m\mathrm{I}N`$. Hence $`Pz=0`$.
In passing we note (for Corollary 3 below) that a $`w^{}`$-accumulation point of $`\{e_n|n\mathrm{I}N\}`$ belongs to $`c_0^{}U_m^{}`$ and has norm one. Accordingly, if $`z`$ is a $`w^{}`$-accumulation point of $`\{x_n|n\mathrm{I}N\}`$, then $`zX_sY_m^{}`$ and $`z=1`$, the latter because each $`Y_m^{}`$ is $`(1\delta _m)`$-isomorphic to $`(l^1)^{\prime \prime }`$ with $`\delta _m0`$.
In addition to Lemma 1 we note that if $`Y`$ is an almost isometric $`l^1`$-copy in an L-embedded space $`X`$ then $`X/Y`$ is L-embedded, too, by \[8, IV.1.3\].
Here is a way of constructing asymptotically isometric $`l^1`$-copies in L-embedded Banach spaces. It is essentially due to Godefroy \[8, IV.2.5\]:
###### Theorem 2.
Let $`X`$ be an L-embedded Banach space with L-decomposition $`X^{\prime \prime }=X_1X_s`$, $`(\mathrm{\Gamma },)`$ a directed set and $`(x_\gamma )_{\gamma \mathrm{\Gamma }}`$ a net in the unit ball of $`X`$. If $`x_\gamma \stackrel{w^{}}{}x_sX_s`$ in the $`w^{}`$-topology of $`X^{\prime \prime }`$ and $`x_s=1`$ then there is a sequence $`(x_{\gamma _n})_{nN}`$ which spans $`l^1`$ asymptotically.
Proof: Let $`(\delta _n)`$ be a sequence of strictly positive numbers converging to $`0`$. Set $`\eta _1=\frac{1}{3}\delta _1`$ and $`\eta _{n+1}=\frac{1}{3}\mathrm{min}(\eta _n,\delta _{n+1})`$ for $`n\mathrm{I}N`$. By induction over $`n\mathrm{I}N`$ we will construct $`\gamma _n\mathrm{\Gamma }`$ such that
(1) $`\left({\displaystyle \underset{i=1}{\overset{n}{}}}(1\delta _i)|\alpha _i|\right)+\eta _n{\displaystyle \underset{i=1}{\overset{n}{}}}|\alpha _i|{\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _ix_{\gamma _i}{\displaystyle \underset{i=1}{\overset{n}{}}}|\alpha _i|n\mathrm{I}N,\alpha _i\mathrm{C}.`$
(The last inequality is trivial because $`x_\gamma 1`$.) We recall that the norm is $`w^{}`$-lower semicontinuous. Thus $`lim\; infx_\gamma x_s=1`$ and for the first induction step we choose $`x_{\gamma _1}`$ such that $`x_{\gamma _1}>1\delta _1+\eta _1`$.
For the induction step $`nn+1`$ we suppose $`x_{\gamma _1},\mathrm{},x_{\gamma _n}`$ to be constructed such that (1) holds. Fix an element $`\alpha =(\alpha _i)_{i=1}^{n+1}`$ in the unit sphere of $`l_{n+1}^1`$ such that $`\alpha _{n+1}0`$. The $`w^{}`$-convergence (along $`\gamma `$) of $`(_{i=1}^n\alpha _ix_{\gamma _i})+\alpha _{n+1}x_\gamma `$ to $`(_{i=1}^n\alpha _ix_{\gamma _i})+\alpha _{n+1}x_s`$ yields
$`\underset{\gamma }{lim\; inf}\left({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _ix_{\gamma _i}\right)+\alpha _{n+1}x_\gamma `$ $``$ $`\left({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _ix_{\gamma _i}\right)+\alpha _{n+1}x_s`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _ix_{\gamma _i}+\alpha _{n+1}x_s`$
$`\stackrel{(\text{1})}{}`$ $`\left({\displaystyle \underset{i=1}{\overset{n}{}}}(1\delta _i)|\alpha _i|\right)+\eta _n\left({\displaystyle \underset{i=1}{\overset{n}{}}}|\alpha _i|\right)+|\alpha _{n+1}|`$
$`=`$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+\eta _n\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}|\alpha _i|\right)`$
$`(\eta _n\delta _{n+1})|\alpha _{n+1}|`$
$``$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+\mathrm{min}(\eta _n,\delta _{n+1})`$
because $`\alpha =1`$ and $`|\alpha _{n+1}|1`$. Thus there is an index $`\gamma _0\mathrm{\Gamma }`$ such that $`inf_{\gamma \gamma _0}\left(_{i=1}^n\alpha _ix_{\gamma _i}\right)+\alpha _{n+1}x_\gamma >\left(_{i=1}^{n+1}(1\delta _i)|\alpha _i|\right)+2\eta _{n+1}`$; note that the subnet $`(x_\gamma )_{\gamma \gamma _0}`$ still $`w^{}`$-converges to $`x_s`$.
Choose a finite $`\eta _{n+1}`$-net $`(\alpha ^l)_{l=1}^{L_{n+1}}`$ in the unit sphere of $`l_{n+1}^1`$ in the sense that for each $`\alpha `$ in the unit sphere of $`l_{n+1}^1`$ there is $`lL_{n+1}`$ such that $`\alpha \alpha ^l=_{i=1}^{n+1}|\alpha _i\alpha _i^l|<\eta _{n+1}`$. Then we may repeat the reasoning above finitely many times for $`l=1,\mathrm{},L_{n+1}`$ in order to get $`x_{\gamma _{n+1}}`$ such that
$`{\displaystyle \underset{i=1}{\overset{n+1}{}}}\alpha _i^lx_{\gamma _i}>\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i^l|\right)+2\eta _{n+1}lL_{n+1}.`$
For an arbitrary $`\alpha `$ in the unit sphere of $`l_{n+1}^1`$ choose $`lL_{n+1}`$ such that $`\alpha \alpha ^l<\eta _{n+1}`$. Then
(2) $`{\displaystyle \underset{i=1}{\overset{n+1}{}}}\alpha _ix_{\gamma _i}`$ $``$ $`{\displaystyle \underset{i=1}{\overset{n+1}{}}}\alpha _i^lx_{\gamma _i}{\displaystyle \underset{i=1}{\overset{n+1}{}}}(\alpha _i\alpha _i^l)x_{\gamma _i}`$
$``$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+2\eta _{n+1}\alpha \alpha ^l`$
$``$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+\eta _{n+1}`$
$`=`$ $`\left({\displaystyle \underset{i=1}{\overset{n+1}{}}}(1\delta _i)|\alpha _i|\right)+\eta _{n+1}{\displaystyle \underset{i=1}{\overset{n+1}{}}}|\alpha _i|.`$
This extends to all $`\alpha l_{n+1}^1`$ and thus ends the induction and the proof.
Within L-embedded spaces James’ distortion theorem is more efficient:
###### Corollary 3.
A sequence spanning $`l^1`$ almost isometrically in an L-embedded Banach space admits of a subsequence which spans $`l^1`$ asymptotically.
Proof: We use the notation of the proof of Lemma 1 and the remark in the end of it. Hence there is a net $`(x_{n_\gamma })`$ which $`w^{}`$-converges to a normalized $`w^{}`$-accumulation point $`x_s`$ of $`\{x_n|n\mathrm{I}N\}`$. Proposition 2 applies.
###### Corollary 4.
Every nonreflexive subspace of an L-embedded Banach space contains an asymptotic copy of $`l^1`$.
In particular, every nonreflexive subspace of an L-embedded Banach space fails the fixed point property.
Proof: L-embedded spaces are $`w`$-sequentially complete (, or \[8, IV.2.2\]). Hence the first assertion follows from theorems of Eberlein-Šmuljan, Rosenthal, James and from Corollary 3. Combine this with \[5, Th. 1.2\] to get the second assertion.
By (or \[8, IV.2.7\]) the basis $`(x_{\gamma _n})`$ in Proposition 2 admits a subsequence whose closed linear span is complemented in $`X`$. In this way Corollary 4 recovers the theorem of Kadec and Pełczyński mentioned in the introduction.
From Lemma 1 we see that both almost isometric and asymptotically isometric $`l^1`$-copies are L-embedded provided they are contained in an L-embedded Banach space. But at the time of this writing it is not clear whether asymptotic $`l^1`$-copies are always L-embedded (or whether asymptotic $`c_0`$-copies are M-embedded). For almost isometric $`l^1`$-copies the situation is clearer:
###### Corollary 5.
There are almost isometric copies of $`l^1`$ which are not L-embedded.
Proof: Combine Corollary 3 or 4 and the counterexample of .
Now we briefly turn to $`c_0`$. Some straightforward modifications of the proof of (or \[8, IV.2.7\]) show that the dual of a nonreflexive L-embedded space contains asymptotic copies of $`c_0`$ although, of course, not every nonreflexive subspace in the dual of an L-embedded space contains $`c_0`$-copies. (Take $`l^{\mathrm{}}`$ for example.) M-embedded spaces provide a more natural frame:
###### Proposition 6.
Every nonreflexive subspace of an M-embedded space contains an asymptotic $`c_0`$-copy.
Proof: The result could be obtained by appropriate modifications of the proof of \[8, III.3.4, 3.7a\] but we suggest a more direct (and shorter) argument.
Since M-embeddedness passes to subspaces it is enough to prove that every nonreflexive M-embedded space $`Z`$ contains an asymptotic $`c_0`$-copy. Let $`(\delta _m)`$ be a sequence in $`]0,1[`$ converging to $`0`$. By induction over $`n\mathrm{I}N`$ we will construct a sequence $`(z_n)`$ in $`Z`$ such that
(3) $`\underset{in}{\mathrm{max}}(1(12^n)\delta _i)|\alpha _i|{\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _iz_i\underset{in}{\mathrm{max}}(1+(12^n)\delta _i)|\alpha _i|\alpha _i\mathrm{C}.`$
For the beginning of the induction we choose $`z_1`$ in the unit sphere of $`Z`$. Suppose that $`z_1,\mathrm{},z_n`$ have been constructed and satisfy (3). Let $`P`$ denote the L-projection on $`Z^{\prime \prime \prime }`$ with range $`Z^{}`$. We have $`P^{}|_Z=\mathrm{id}_Z`$ and $`Z^{\prime \prime \prime \prime }=Z^{}_{\mathrm{}}Z^{}^{}`$. There exists an element $`z^{}{}_{}{}^{}Z^{}^{}`$ with $`z^{}{}_{}{}^{}=1`$ because $`Z`$ is not reflexive. Put $`E=\text{lin}(\{z_i|in\}\{z^{}{}_{}{}^{}\})`$ and choose $`\eta >0`$ such that
$`(1\eta )(1(12^n)\delta _i)`$ $`>`$ $`(1(12^{(n+1)})\delta _i)`$
$`(1+\eta )(1+(12^n)\delta _i)`$ $`<`$ $`(1+(12^{(n+1)})\delta _i)`$
for all $`in`$. The principle of local reflexivity provides an operator $`T_1:EZ^{\prime \prime }`$ and an operator $`T_2:T_1(E)Z`$ such that $`T=T_2T_1`$ fulfills
(4) $`T|_{EZ}=\mathrm{id}_{EZ}\text{and}(1\eta )eTe(1+\eta )eeE.`$
Put $`z_{n+1}=Tz^{}^{}`$. Then we get (3, $`n+1`$):
$`\underset{in+1}{\mathrm{max}}(1(12^{(n+1)})\delta _i)|\alpha _i|`$ $`<`$ $`(1\eta )\mathrm{max}(\underset{in}{\mathrm{max}}(1(12^n)\delta _i)|\alpha _i|,|\alpha _{n+1}|)`$
$`\stackrel{(\text{3})}{}`$ $`(1\eta )\mathrm{max}({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _iz_i,\alpha _{n+1}z^{}{}_{}{}^{})`$
$`=`$ $`(1\eta )\left({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _iz_i\right)+\alpha _{n+1}z^{}{}_{}{}^{}`$
$`\stackrel{(\text{4})}{}`$ $`{\displaystyle \underset{i=1}{\overset{n+1}{}}}\alpha _iz_i`$
$`\stackrel{(\text{4})}{}`$ $`(1+\eta )\left({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _iz_i\right)+\alpha _{n+1}z^{}{}_{}{}^{}`$
$`=`$ $`(1+\eta )\mathrm{max}({\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _iz_i,|\alpha _{n+1}|)`$
$`<`$ $`\underset{in+1}{\mathrm{max}}(1+(12^{(n+1)})\delta _i)|\alpha _i|.`$
This ends the induction and the proof.
It is now immediate that every nonreflexive subspace of an M-embedded space fails the fixed point property (see \[6, Prop. 7\]).
The analogue of Lemma 1 holds trivially, simply because M-embeddedness passes to subspaces \[8, III.1.6\]; thus every $`c_0`$-copy inside an M-embedded space is M-embedded. And similarly as for Corollary 5, the counterexample of combined with Proposition 6 provides an almost isometric $`c_0`$-copy which is not M-embedded.
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# Critical state in thin anisotropic superconductors of arbitrary shape
## I Introduction
Numerous publications deal with the critical state problem in isotropic thin superconductors placed in a perpendicular magnetic field. By thin we mean that the sample thickness $`d`$ is considerably less than its lateral extension. Analytic solutions of the problem were obtained in the cases of a circular disk , a strip , and elliptic-shaped films. For thin superconductors of an arbitrary shape, numerical methods for solving this problem were elaborated in Ref. , and the critical state of a rectangular platelet was investigated in detail in Ref. . All the above-mentioned results can be derived considering the appropriate flat superconductors to be infinitely thin. In the framework of this approach a possible in-plane anisotropy of pinning can be taken into account, yielding good agreement with appropriate magnetooptic experiments.
It is important to note that in the highly anisotropic high-$`T_c`$ superconductors the flux-line pinning in general depends on the angle $`\theta `$ between the local direction of the magnetic induction $`𝐁`$ and the $`c`$ axis which in typical experiments is normal to the plane of the samples. For example, this type of anisotropy occurs when one takes into account the intrinsic pinning exerted by the CuO planes. Besides this, twin boundaries, columnar defects, and other extended defects give rise to such an out-of-plane anisotropy. In all these cases, the flux-line curvature that always occurs in thin superconductors placed in a perpendicular magnetic field, leads to a dependence of the critical current density $`j_c`$ on the coordinate $`z`$ across the thickness of the sample. Therefore, the critical state problem becomes three-dimensional (3D), and the feasibility of treating such superconductors as infinitely thin requires special consideration.
For a thin anisotropic disk with radius $`R`$ placed in a perpendicular magnetic field it was shown that the small ratio $`d/R`$ enables one to split the 2D critical state problem into two one-dimensional problems. The first one treats the critical state across the thickness of the disk, while the second one treats the disk as infinitely thin and gives the distributions of magnetic field and sheet current along its radius. Using this method, the magnetic moment of a thin anisotropic disk was theoretically studied in the case when the angular dependent $`j_c`$ has a peak at $`\theta =0`$. In this special geometry of the circular disk the directions of field and current are fixed and known in advance, and the critical state problem thus is essentially not 3D but only 2D. This simplification does not take place in thin samples of arbitrary shape, where the current stream lines, magnetic contour lines, and penetrating flux fronts in general do not coincide. Moreover, according to Ref. , in the general case the directions of the circulating currents may change across the thickness of the sample, and a rotation or twist of the flux line arrangement is possible.
In the present paper we show that even in the general 3D critical state problem of a thin superconductor with arbitrary shape and arbitrary anisotropy of $`j_c`$, the 3D pinning problem can be split into two simpler problems which can be solved: The 1D problem across the thickness of the sample is solved analytically, while the in-plane problem reduces to that of the infinitely thin superconductor which was treated before. Since the methods of solving the latter problem are well elaborated, our finding enables one to study the critical state in flat samples of arbitrary shape and with any type of anisotropy. As concrete examples, we shall investigate the critical states of long strips and rectangular platelets with various types of out-of-plane anisotropy.
## II Splitting the 3D problem
In what follows we assume that the pinning-caused characteristic magnetic field $`H_{cr}=j_cd/2`$ considerably exceeds the lower critical field $`H_{c1}`$, and the thickness $`d`$ of the sample is much greater than the London penetration depth $`\lambda `$. Under these conditions the magnetic induction is practically equal to the magnetic field in the superconductor, and the so-called geometrical barrier is negligible. We study the critical state problem in the framework of the macroscopic approach, in which all considered lengths are larger than the flux-line spacing, and treat the superconductor as a uniform anisotropic medium.
Let us place the coordinate system so that its $`xy`$ plane coincides with the middle plane of the sample and the $`z`$ axis is directed along the external magnetic field; the case of a rectangular platelet is shown in Fig. 1. Since the thickness $`d`$ of the sample is assumed to be much less than its smallest lateral dimension $`L`$, the strict equations $`𝐣=0`$ and $`𝐁=0`$ give that, in leading order in the parameter $`d/L`$, one has $`j_z=0`$, and the perpendicular field component $`B_z`$ is independent of $`z`$ inside the sample, i.e., $`B_z=B_z(x,y)`$. Accounting for the fact that the scale of changes with $`x`$ and $`y`$ of $`B_z`$ and of the in-plane field $`𝐁_t=(B_x,B_y)`$ is $`L`$, while the appropriate scale for changes of $`𝐁_t`$ across the thickness is $`d`$, we may neglect the derivatives of $`𝐁`$ with respect to $`x`$ and $`y`$. The critical state equation can thus be written in the form:
$`\widehat{𝐳}\times {\displaystyle \frac{𝐁_t}{z}}=\mu _0𝐣_c,`$ (1)
where $`\widehat{𝐳}`$ is the unit vector along $`z`$ and $`𝐣_c`$ is the critical current density. This equation means that the critical gradient develops predominantly across the thickness of the sample, and we can consider the fields and the currents to be independent of $`x,y`$ in sufficiently small parts of the sample with dimensions in the $`x,y`$ directions much greater than $`d`$ and much less than $`L`$. To proceed our analysis we write $`𝐁_t=B_t\widehat{𝐭}`$ and $`𝐣_c=j_c\widehat{𝐧}`$, where $`\widehat{𝐧}=(\mathrm{cos}\phi ,\mathrm{sin}\phi )`$ and $`\widehat{𝐭}=(\mathrm{cos}\psi ,\mathrm{sin}\psi )`$. The angles $`\phi (x,y,z)`$ and $`\psi (x,y,z)`$ define the directions of $`𝐣_c`$ and $`𝐁_t`$ in $`xy`$ planes. The above critical state equation can then be rewritten as
$`{\displaystyle \frac{B_t}{z}}=\mu _0j_c\mathrm{sin}(\phi \psi ),`$ (2)
$`B_t{\displaystyle \frac{\psi }{z}}=\mu _0j_c\mathrm{cos}(\phi \psi ).`$ (3)
We emphasize that in the general case the critical current density $`j_c`$ may depend on $`B_z`$, $`B_t`$, $`\phi \psi `$, $`\psi `$, i.e., $`j_c=j_c(B_z,B_t,\phi \psi ,\psi )`$. The quantities $`B_t`$ and $`B_z`$ can be also expressed in terms of the magnitude $`B|𝐁|`$ and the tilt angle $`\theta `$ between the local direction of $`𝐁`$ and the z axis:
$`B_z=B\mathrm{cos}\theta ,B_t=B\mathrm{sin}\theta .`$
The dependence on $`\psi `$ occurs only if there is an anisotropy of pinning in the $`xy`$ plane (e.g., when twin boundaries exist in the sample). In the isotropic case, when the flux-line pinning depends neither on $`\psi `$ nor on the tilt angle $`\theta `$, the dependence of $`j_c`$ on $`B_z`$ and $`B_t`$ occurs only through the magnitude $`B`$. Note also that, since $`B_z`$ is independent of $`z`$, it enters Eqs. (1,2) as a parameter, and thus we have only two equations for the three functions $`B_t`$, $`\phi `$, and $`\psi `$. One more equation can be obtained if one takes into account the prehistory of the critical state, telling how the flux lines were pushed into the sample.
In increasing moderate magnetic field there is a flux and current-free core in the superconductor, Fig. 1. The surface of the core, $`\gamma `$, is the penetrating flux front, which we describe by a function $`z=z_\gamma (x,y)`$. In the limiting case of an infinitely thin superconductor ($`d0`$) the flux front may be thought of as a flat curve $`\gamma _0`$, forming the outer rim (equator) of $`\gamma `$, since we have $`B_z=0`$ inside this contour. In the case of small but finite $`dL`$, the component $`B_z`$ is still practically zero in the region inside $`\gamma _0`$ since there the flux lines are almost parallel to the flat surfaces of the sample. Note that these flux lines may have different orientations at different values of $`z`$ if the sample is not a circular disk or a strip; thus, in principle, flux-cutting processes may occur in the superconductor.
We now consider the third required equation in the region $`z_\gamma (x,y)|z|d/2`$ with $`x,y`$ inside the curve $`\gamma _0`$, i.e. between the lens-shaped empty core and the flat surfaces. In increasing applied field the flux lines are pushed continuously from the flat surfaces into this region, such that the in-plane orientation of the flux lines is a unique function of $`B_t`$, $`x`$, and $`y`$:
$`\psi =\mathrm{\Psi }(B_t,x,y).`$ (4)
The function $`\mathrm{\Psi }`$ is found from the solution of the critical state problem for the infinitely thin sample (see below). Equation (3) is just the third required equation. Now we are able to solve Eqs. (1–3). Dividing Eq. (2) by Eq. (1) we obtain
$`B_t{\displaystyle \frac{\psi }{B_t}}=\mathrm{cot}(\phi \psi ).`$ (5)
Inserting here the known field orientation, Eq. (3), we find the difference $`\phi \psi `$ and the angle of the currents $`\phi `$ as functions of $`B_t`$. Finally, the dependence of $`B_t`$ on $`z`$ can be obtained from the implicit relation
$`zz_\gamma ={\displaystyle _0^{B_t}}{\displaystyle \frac{db_t}{\mu _0j_c(0,b_t,\pi /2,\mathrm{\Psi }(b_t,x,y))}},`$ (6)
which follows from Eq. (1), and from the formula
$`j_c(0,B_t,\phi \psi ,\psi )={\displaystyle \frac{j_c(0,B_t,\pi /2,\psi )}{\mathrm{sin}(\phi \psi )}},`$ (7)
reflecting the fact that only the component of $`j_c`$ normal to $`𝐁`$ is essential (in absence of flux cutting). The function $`z_\gamma (x,y)`$ describing the shape of the flux front is obtained from the equality:
$`{\displaystyle \frac{d}{2}}z_\gamma ={\displaystyle _0^{B_s}}{\displaystyle \frac{db_t}{\mu _0j_c(0,b_t,\pi /2,\mathrm{\Psi }(b_t,x,y))}},`$ (8)
where $`B_s(x,y)=B_t(x,y,d/2)`$ is the tangential component of the surface field. This $`B_s`$ may be obtained by solving the critical state equation for infinitely thin superconductors. Thus, Eqs. (3–6) give the solution of the critical state problem across the thickness of the sample in the region inside the curve $`\gamma _0`$.
It should be emphasized that the above formula for the dependence of $`j_c`$ on $`\phi \psi `$ and Eqs. (3,5,6) hold only if flux cutting does not occur in the superconductor when the critical state is established in it. Flux cutting will only happen if $`|\mathrm{cot}(\phi \psi )|`$, calculated from Eq. (4), turns out to be greater than some critical value $`\chi `$. In this case the so-called generalized critical-state model must be used to describe the critical state in the region of the sample inside $`\gamma _0`$. However, since at present we cannot point out the shape of the superconductor for which this flux cutting really occurs, we do not analyze this case here.
Let us now derive the third equation for the region where $`B_z0`$, i.e., outside the curve $`\gamma _0`$. In deriving it is necessary to take into account that in an anisotropic superconductor the direction of the flux-line velocity in general does not coincide with the direction of the Lorentz force causing this movement, and thus the appropriate electric field $`𝐄`$ is not always parallel to the current density (see Appendix A). With this in mind, the third equation results from the following two conditions: First, the shape of a flux line cannot change when an applied magnetic field $`H`$ increases slightly so that $`\gamma _0`$ and all the flux lines in this region of the sample are shifted by distances which are much less than the sample width $`L`$. Second, during the shift, each line element moves along the local direction of $`[𝐄\times 𝐁]`$. The geometrical analysis of these conditions gives
$`{\displaystyle \frac{B_z^2\mathrm{sin}\xi B_t^2\mathrm{cos}\psi \mathrm{sin}(\psi \xi )}{B_z^2\mathrm{cos}\xi +B_t^2\mathrm{sin}\psi \mathrm{sin}(\psi \xi )}}=\mathrm{const}`$
where the angle $`\xi `$ is determined by Eqs. (A2,A3,A5) in Appendix A. Since $`B_t(0)=0`$, the obtained relation can be rewritten as
$`\mathrm{tan}[\xi \xi (0)]={\displaystyle \frac{B_t^2\mathrm{sin}[\psi \xi (0)]\mathrm{cos}[\psi \xi (0)]}{B_t^2\mathrm{cos}^2[\psi \xi (0)]+B_z^2}},`$ (9)
with $`\xi (0)\xi (x,y,\mathrm{\hspace{0.17em}0})`$. This is the third required equation in the region where $`B_z0`$. When an in-plane anisotropy of the flux line pinning is absent, it is possible to find the solution of Eqs. (1,2,7) in the interval $`0zd/2`$. The initial conditions to these equations are
$`B_t(0)=0,\phi (0)\psi (0)=\pi /2,`$ (10)
where the last equality follows from Eq. (2) with $`B_t(0)=0`$. Since the solution of the equations is unique, it is sufficient to guess it. The solution has the form
$`\phi (x,y,z)=\phi _0(x,y),`$ (11)
$`\psi (x,y,z)=\phi _0(x,y)\pi /2\psi _0(x,y),`$ (12)
$`z={\displaystyle _0^{B_t}}{\displaystyle \frac{db_t}{\mu _0j_c(B_z,b_t,\pi /2,\psi _0)}}.`$ (13)
Here the angle $`\phi _0(x,y)`$ defines the direction of the sheet current in the critical state of the infinitely thin superconductor and is assumed to be known. In obtaining the solution we have taken into account that $`\xi =\phi `$ at $`\phi \psi =\pi /2`$ when the in-plane anisotropy is absent. Eqs. (9–11) mean that the flux lines are curved in the region where $`B_z0`$, but they are not twisted there. In other words, the distributions of the magnetic induction and the current across the thickness of the sample are the same as in a circular disk or a strip if one considers small regions in the $`xy`$ plane with the same values of $`B_z`$. Eq. (11) allows us to find an implicit relation between the sheet current $`J`$ and $`B_z`$:
$`{\displaystyle \frac{d}{2}}={\displaystyle _0^{J/2}}{\displaystyle \frac{dh_t}{j_c(B_z,\mu _0h_t,\pi /2,\psi _0)}}.`$ (14)
Here we have used that $`B_t(z=d/2)B_s=\mu _0J/2`$, where $`J|𝐉|`$ and $`𝐉`$ is the current density integrated over the film thickness. It should be noted that although we indicate explicitly the angle $`\psi _0`$ as an argument of $`j_c`$ in Eqs. (11,12), the critical current density is independent of this angle in cases without in-plane anisotropy.
We now discuss the case when some in-plane anisotropy of flux-line pinning exists in the superconductor together with the out-of-plane anisotropy, i.e. we allow for a dependence of $`j_c`$ on $`\psi `$. This case can occur, e.g. in twinned crystals of high-$`T_c`$ superconductors. When an in-plane anisotropy exists, it is impossible to find the solution of Eqs. (1,2,7) in the general form. We may only state that in this situation the in-plane angle $`\psi `$ generally depends on $`z`$, and thus flux lines are not only curved but also twisted. Then, according to Eq. (2), the current is not normal to the magnetic induction. The magnitude of the twist is determined by the relative strengths of the in-plane and out-of-plane anisotropies (and, of course, by the thickness of the sample). Interestingly, when $`j_c`$ does not depend on $`B_t`$, the solution described by Eqs. (9–11) becomes true again, and the twist disappears. (This situation occurs, e.g., in sufficiently large magnetic fields, $`HH_{cr}`$, or if there is no out-of-plane anisotropy in the superconductor and $`j_c`$ is practically independent of the magnitude of the magnetic induction). However, if $`H`$ changes, in this case a flux line is shifted in a direction which does not lie in the plane containing the line. According to the results of Appendix A, the direction makes an angle $`\alpha `$ with the above-mentioned plane where
$`\mathrm{tan}\alpha ={\displaystyle \frac{[\mathrm{ln}j_c(B_z,\phi \psi _0,\psi _0)]}{\phi }}|_{\phi =\psi _0+{\displaystyle \frac{\pi }{2}}}.`$
To sum up, it is important to emphasize that at a given dependence of $`j_c`$ on $`B_z`$, $`B_t`$, $`\phi \psi `$ and $`\psi `$, the one-dimensional equations (1), (2), (7) with initial conditions (8) can be solved (at least numerically) in the case of an arbitrary anisotropy of flux-line pinning, and a relation $`J(B_z,\phi _0)`$ generalizing Eq. (12) can be obtained. Here $`\phi _0`$ defines the direction of the sheet current in the $`xy`$ plane.
Thus, we have reduced the 3D critical state problem to a 2D problem in which the $`B_z`$-dependence of the sheet current $`J(B_z,\phi _0)`$ is determined not only by the $`B_z`$-dependence of the critical current density $`j_c`$ but also by its dependence on the tilt angle $`\theta `$ of the flux lines away from the film normal, see Eq. (12). It should be noted, however, that to obtain this reduction we have assumed that $`d`$ is much less than the characteristic scale of changes of $`𝐁`$ in the $`xy`$ plane. Numerical solutions of the critical state equation for infinitely thin superconductors of various shapes show that this assumption indeed is true everywhere; the scale of variation is of the order of $`L`$ except for a small region near the flux front $`\gamma _0`$ where the derivatives $`B_x/y`$ and $`B_y/x`$ are not small compared with $`𝐁_t/z`$. However, the width of this region is determined by the thickness $`d`$ and is thus small.
Knowing the function $`J(B_z,\phi _0)`$ obtained from Eq. (12) or from the above-mentioned generalization of this equation, it is possible to solve the critical state equations for an infinitely thin superconductor of any shape. The static equations to be solved are
$$𝐉=0$$
(15)
and the Biot-Savart law connecting $`B_z`$ with the sheet current $`𝐉(x,y)=_{d/2}^{d/2}𝐣(x,y,z)𝑑z`$,
$$\mu _0^1B_z(𝐫)=H+_S\frac{[𝐑\times 𝐉]}{4\pi R^3}d^2r^{},$$
(16)
where $`𝐑=𝐫𝐫^{}`$; $`𝐫`$ and $`𝐫^{}`$ are 2D vectors in the $`xy`$ plane, and the integration is carried out over the specimen area $`S`$. In addition, one has to take into account that $`B_z=0`$ inside the 2D flux front $`\gamma _0`$, while outside $`\gamma _0`$ the sheet current is equal to $`J(B_z,\phi _0)`$. At given applied field $`H`$, a solution of this static problem exists only for a certain curve $`\gamma _0`$, which determines the shape and position of the penetrating flux front. In Refs. , in fact, the more general dynamic equations were considered, which allow one to find the solution of the 2D critical state equations very effectively for any prehistory $`H(t)`$. The appropriate formulas are presented in Appendix B.
Solving these equations for the infinitely thin superconductor of arbitrary shape in a perpendicular applied field $`H`$, one obtains $`B_z(x,y,H)`$, $`J(x,y,H)`$, $`\phi _0(x,y,H)`$, $`B_s(x,y,H)=\mu _0J/2`$, the direction of $`𝐁_s`$, i.e., the function $`\psi _0(x,y,H)=\phi _0(x,y,H)\pi /2`$, and also the position of $`\gamma _0`$. Eliminating $`H`$ from $`\psi _0(x,y,H)`$ and $`B_s(x,y,H)`$ in the region inside $`\gamma _0`$, and replacing $`B_s`$ by $`B_t`$ we arrive at the function $`\mathrm{\Psi }(B_t,x,y)`$ that enters Eq. (3). Thus, the solution of the 2D critical state problem for the infinitely thin superconductor enables one to obtain all the necessary information for the calculations of the core shape $`z_\gamma (x,y)`$ and distributions of $`B_t`$, $`\psi `$, $`\phi `$ across $`z`$ at any $`x`$ and $`y`$. The knowledge of all these functions permits to describe the 3D critical state of thin anisotropic superconductors in detail.
## III Out-of-plane Anisotropy
We now use the results obtained above to investigate the critical states of thin superconductors with out-of-plane anisotropy alone. That is, we consider the case when the critical current density $`j_c(B_z,B_t,\pi /2,\psi )`$ does not depend on the in-plane angle $`\psi `$. In addition we here neglect its dependence on $`B`$ such that $`j_c`$ depends only on the flux-line tilt angle $`\theta `$, $`j_c=j_c(\theta )`$. This situation occurs when the flux-line pinning is isotropic in the $`xy`$ plane, and the scale $`B_0`$ characterizing the dependence of $`j_c`$ on $`B`$ considerably exceeds the self-fields of the critical currents. As shown above, the anisotropy of $`j_c(\theta )`$ can be accounted for by introducing an appropriate $`B_z`$ dependence of the sheet current $`J=J(B_z)`$, implicitly given by
Eq. (12). In the considered case, a more explicit form of this dependence may be obtained in the following way. Differentiation of Eq. (12) with respect to $`B_z`$ yields the two equations:
$`j_c(\theta )d=J(B_z)B_z{\displaystyle \frac{dJ(B_z)}{dB_z}},`$ (17)
$`\mathrm{tan}\theta ={\displaystyle \frac{\mu _0J(B_z)}{2B_z}}.`$ (18)
In fact, Eqs. (15,16) are the parametric form of that function $`j_c(\theta )`$ which results in a given dependence $`J(B_z)`$ via Eq. (12).
To fix ideas, we consider the following rather general model dependence for $`J(B_z)`$:
$`J(B_z)=j_c(0)d\left[1+p\mathrm{exp}\left({\displaystyle \frac{qB_z}{B_{cr}}}\right)\right],`$ (19)
where $`B_{cr}=\mu _0j_c(0)d/2`$, while $`j_c(0)`$, $`p`$ and $`q`$ are the parameters of this example ($`p`$ and $`q`$ are dimensionless). One then easily verifies that the corresponding angular dependence of the critical current density takes the form:
$`j_c(\theta )=j_c(0)\left[1+p(1+qt)\mathrm{exp}(qt)\right],`$ (20)
$`\mathrm{tan}\theta =t^1\left[1+p\mathrm{exp}(qt)\right],`$ (21)
where $`t`$ is a curve parameter with range $`0t\mathrm{}`$ equivalent to $`\pi /2\theta 0`$. This model dependence $`j_c(\theta )`$ is presented in Fig. 2 together with the corresponding $`J(B_z)`$, Eq. (17). For appropriate choices of the parameters $`p`$ and $`q`$ this model describes the intrinsic pinning by the CuO planes in high-$`T_c`$ superconductors ($`p>0`$, maximum $`j_c`$ at $`\theta =\pi /2`$) and pinning by columnar defects perpendicular to the film ($`p<0`$, maximum $`j_c`$ at $`\theta =0`$). It may thus be used to simulate these important cases. We remark that the critical state of a circular disk with some model for intrinsic pinning was studied numerically in Ref. , while the magnetization of the disk in the case when $`j_c`$ has a peak at $`\theta =0`$ was analyzed in Ref. . However, in those papers only the fully penetrated critical state was considered.
Below we calculate the partially and fully penetrated critical states in strips and rectangular platelets with anisotropic pinning using the angular dependence $`j_c(\theta )`$ described by Eqs. (18). The corresponding results for circular disks are very similar to the results for the strip. In fact the algorithm used below to time-integrate the 1D equation of motion for the sheet current $`J(x,t)`$ in thin strips is almost identical to the algorithm used for thin disks.
## IV Thin Strip
We begin with the calculation of the critical state in an infinitely thin strip (see Appendix B). In this calculation the function $`J(B_z)`$ described by Eq. (17) is considered as the critical value of the sheet current. The obtained results are presented in Figs. 3–7.
Figure 3 shows the induction profiles $`B_z(x)`$ and sheet current profiles $`J(x)`$ in increasing $`H`$ for thin strips with various anisotropies of $`j_c(\theta )`$ depending on the parameters $`p`$ and $`q`$. The profiles are presented for several values of $`x_0`$, the position of the 2D flux front $`\gamma _0`$. The reference case of isotropic pinning ($`p=0`$) is also shown here. In Fig. 3a one can see that for $`p>0`$ (pinning by the CuO layers) the penetrating flux-density profile is very steep at $`x=x_0`$: $`B_z(x)`$ jumps almost abruptly and $`J(x)`$ has a very sharp peak which ideally should reach the maximum value $`J(x_0)=(1+p)j_c(0)d`$, cf. Eq. (17). In our computation this ideal maximum height is not reached due to our finite creep exponent $`n=101`$ and finite number $`N=300`$ of equidistant grid points (to illustrate the influence of $`n`$ the case of a small value $`n=11`$ is shown in Fig. 3d). As opposed to this, in Figs. 3b and 3c the induction profiles for $`p<0`$ (pinning by linear defects perpendicular to the strip plane) are less sharp: at $`x=x_0`$, $`B_z(x)`$ vanishes less steeply (almost linearly for Fig. 3b) and $`J(x)`$ decreases monotonically and exhibits an inflection point with infinite slope, $`dJ/dx|_{x=x_0}=\mathrm{}`$, $`J(x_0)=(1+p)j_c(0)d`$. These general features of the flux front of superconductors with
(a)
(b)
(c)
(d)
anisotropic pinning are derived analytically in a forthcoming paper. It turns out that in the vicinity of $`x_0`$ the profiles are $`B_z(x)(xx_0)^\beta `$ and $`J(x)J(x_0)|xx_0|^\beta `$ where $`\beta 0.5\pi ^1\mathrm{arctan}(pq)`$.
In the isotropic case ($`p=0`$) our numerical results practically completely coincide with the exact solution
$`{\displaystyle \frac{J(x)}{j_cd}}={\displaystyle \frac{2}{\pi }}\mathrm{arccot}\left\{\left[\mathrm{max}(0,{\displaystyle \frac{(x_0^2x^2)a^2}{(a^2x_0^2)x^2}})\right]^{1/2}\right\},`$ (22)
$`x_0={\displaystyle \frac{a}{\mathrm{cosh}(\pi H/j_cd)}}.`$ (23)
where $`0xa`$ and $`2a`$ is the width of the strip. It should be also mentioned that in the case of monotonically decreasing $`J_c(B_z)`$ the critical states of a thin strip and a thin disk were studied in Refs. , respectively. In Ref. the Kim model $`J_c(B_z)=J_c(0)B_0/(B_0+|B_z|)`$ was considered, while in Ref. the Kim model and the exponential dependence $`J_c(B_z)=J_c(0)\mathrm{exp}(|B_z|/B_0)`$ were analyzed where $`J_c(0)`$ and $`B_0`$ are some constant parameters. Our results in the case $`p>0`$ are similar to the results in these papers. Equations (15,16) show that both models may be interpreted as angular dependences of $`j_c(\theta )`$ with maxima at $`\theta =\pi /2`$ if $`B_0`$ is of the order of $`J_c(0)`$.
The position $`x=x_0`$ of the flux front $`\gamma _0`$ as a function of the increasing applied field $`H`$ is shown in Fig. 4 for the same anisotropic strips. One realizes that for $`p>0`$ the penetration depth $`ax_0`$ is smaller, and for $`p<0`$ larger, than for the isotropic case $`p=0`$. The isotropic case is shown twice in Fig. 4: The numerical result for $`x_0`$ obtained with a finite creep exponent $`n=101`$ is only slightly smaller than the ideal static result, Eq. (20). This demonstrates the accuracy of our computations.
Using the above results for the infinitely thin strip and equations of Sec. II, we can now describe the two dimensional critical state of the anisotropic strip with small but finite thickness.
The general 3D flux front $`\gamma `$ (Fig. 1) is determined by Eq. (6) which is rewritten for our model Eqs. (17,18) as
$`z_\gamma (x,y)={\displaystyle \frac{d}{2}}\left[1{\displaystyle \frac{J(x,y)}{J_{\gamma _0}}}\right],`$ (24)
where $`J_{\gamma _0}=(1+p)j_c(0)d`$ and $`J(x,y)`$ is the value of the sheet current at a given point inside $`\gamma _0`$. Fig. 5 shows $`z_\gamma (x)`$, the spatial profile or cross section of flux fronts $`\gamma `$, in thin strips with anisotropic pinning. Note that for $`p>0`$ the fronts have a wide flat part near the central plane $`z=0`$, while for $`p<0`$ the fronts have a more rounded shape near $`x=x_0`$ as compared to the isotropic case $`p=0`$.
Physically, the wide flat flux front for $`p>0`$ results from the fact that in this case the flux lines oriented perpendicularly to the film (and to the CuO planes) are less pinned and can thus penetrate more easily than the flux lines oriented parallel to the film plane. In the opposite case $`p<0`$, flux lines in the film plane move more easily and cause a more pointed wedge-like front.
The critical states of thick isotropic strips and disks were calculated in Refs. using the appropriate 2D equations. The isotropic fronts ($`p=0`$) obtained from Eq. (21) and depicted in Fig. 5 coincide with the flux fronts computed for such strips and disks if one takes $`da`$ and $`j_c=j_c(0)`$. These flux front profiles are equal for strips and disks since the appropriate static sheet currents $`J(x)`$ and $`J(r)`$ have identical forms described by Eqs. (19,20); for the disk $`x`$ should be replaced by the radial coordinate $`r`$ and $`a`$ means the disk radius. Interestingly, the thin film solution can well describe features of thick strips and disks up to quite large aspect ratios $`d/2a0.2`$ as long as the flux front is not too close to the center ($`x_0>d`$).
In a strip the sheet current has a fixed direction, and thus there is no rotation of the flux lines in the region inside $`\gamma _0`$, i.e., $`\psi =`$const in Eq. (3). On the other hand, since in the considered case any in-plane anisotropy of pinning is absent, the flux lines are not twisted in the region $`xx_0`$, and the solution (9–11) is valid. The distribution of $`B_t=B_x(x,y,z)`$ in the strip can be found using Eqs. (5,11). For our model (17,18), these equations are solved analytically, yielding for $`xx_0`$
$`B_x(x,z)=0,|z|z_\gamma (x),`$ (25)
$`|B_x(x,z)|=j_c(0)(1+p)[|z|z_\gamma (x)],`$ (26)
$`z_\gamma (x)|z|d/2,`$ (27)
while in the region $`xx_0`$ we arrive at
$`|B_x(x,z)|=j_c(0)|z|\left[1+p\mathrm{exp}\left({\displaystyle \frac{qB_z(x)d}{2|z|B_{cr}}}\right)\right],`$ (28)
$`|z|d/2.`$ (29)
Note that in the anisotropic case ($`p0`$) $`B_x`$ is not linear in $`z`$ in the region $`xx_0`$. Taking into account these equations and the fact that $`B_z`$ is independent of $`z`$, one finds the 2D distributions of $`𝐁`$ inside the thin anisotropic strips shown in Fig. 6. As can be seen, in the isotropic case our approach gives the expected magnetic field lines except in a narrow region near the flux front, where this approximation breaks down (see Sec. II). The width of this region is less than, or of the order of $`d`$.
Knowing the direction of $`𝐁(x,z)`$ we can further calculate the distribution of $`j=j_c`$ in the region outside the 3D front $`\gamma `$ (inside $`\gamma `$ one has $`j=0`$). When $`x<x_0`$, the flux lines are practically parallel to the $`xy`$ plane, and thus $`j_c=j_c(\pi /2)=j_c(0)(1+p)`$. If $`x>x_0`$, $`j_c`$ is given by the first equation (18) with $`t`$ replaced by the ratio $`B_z(x)/[\mu _0j_c(0)z]`$. Note that for $`p>0`$ and $`x>x_0`$, $`j_c(x,z)`$ is greater near the surfaces ($`|z|d/2`$) than in the central plane of the strip ($`z=0`$). The opposite situation occurs when $`p<0`$.
In Fig. 7 the virgin magnetization curves and magnetic hysteresis loops $`M(H)`$ are presented for strips with various anisotropies of pinning. For $`ab`$ plane pinning ($`p>0`$) a so-called central peak appears. Thus, in high-$`T_c`$ superconductors this often observed peak in principle may be caused by intrinsic pinning by the CuO planes. Note that in this case the width of the peak is proportional to the thickness $`d`$ of the sample (our field units are $`j_c(0)d`$). From this one may estimate the contribution of the intrinsic pinning when the peak is observed. In the opposite case $`p<0`$, i.e. when $`j_c(\theta )`$ has a minimum at $`\theta =\pi /2`$, a dip in $`|M(H)|`$ occurs in the region of weak fields. An additional dependence of $`j_c`$ on $`|𝐁|`$ may then lead to a so-called fishtail (or peak) effect with a maximum of $`|M|`$ at some intermediate field. In Fig. 7 the zig-zag sweep of the applied field $`H(t)`$ has constant ramp rate, $`|dH/dt|=1`$ in units where $`j_c(0)d=a=\mu _0=E_0=1`$ in the used current–voltage law $`E(J)=E_0(J/J_c)^n`$ with $`n=51`$, cf. Appendix B. When the creep exponent $`n`$ is lowered to a more realistic value $`n=11`$ measured e.g. in Ref. , the amplitudes of the peak and dip slightly decrease, and the sharp edges of $`M(H)`$ at the largest and lowest $`H_a`$ are rounded, see also the figures in Refs. d, .
## V Thin Rectangular Plates
The critical state of rectangular platelets is qualitatively similar to that of the strip. In particular, for aspect ratios $`b/a>1.4`$ (see Figs. 1, 8, 9) the profiles $`J_y(x,0)`$ and $`B_z(x,0)`$ along the shorter axis are practically identical to the profiles in long strips with $`b/a\mathrm{}`$. There are, however, several qualitative differences as compared to the magnetic behavior of long strips or circular disks:
a) In platelets with shape different from strips or circular disks, the stream lines of the sheet current $`𝐉(x,y)=(J_x,J_y)`$ do no longer coincide with the contour lines of $`B_z(x,y)`$ and the penetrating flux fronts $`\gamma _0`$ (Fig. 1). This new feature is described in detail in Ref. 4 for films with elliptical shape.
b) In increasing applied field $`H`$, the direction of the sheet current $`𝐉(x,y)`$ changes with $`H`$ at any given point $`(x,y)`$ positioned inside the front $`\gamma _0`$ and away from the symmetry axes $`x=0`$, $`y=0`$, or edges $`x=\pm a`$, $`y=\pm b`$. In the isotropic case this rotation of $`𝐉(x,y)`$ comes to a halt when the flux front passes through the point.
Both features a) and b) occur only in perpendicular geometry. In long rods with the same rectangular cross section in longitudinal field, the $`B_z`$ contours, $`𝐉`$ stream lines, and flux fronts are concentric rectangles, and inside the flux front both $`B_z(x,y)`$ and $`𝐉(x,y)`$ vanish. From the feature b) follows a further interesting effect:
In thin rectangles a rotation of the penetrating flux lines may occur inside the front $`\gamma _0`$. In this region the flux lines surround a thin flat field-free core and are nearly parallel to the $`x,y`$ plane. However, at a given point $`(x,y)`$ the flux lines lying at different $`z`$ can have different in-plane orientations. This rotation of flux lines is due to the rotation of the direction of the sheet current and occurs in any isotropic or anisotropic thin superconductor with shape different from a strip or circular disk. An example of this rotation is shown in Fig. 8.
The stream lines of the sheet current and the contours of $`B_z(x,y)`$ in a thin rectangular plate with strong out-of-plane anisotropy of pinning are depicted in Fig. 9 for increasing applied field $`H`$. The stream lines are contours of the function $`g(x,y)`$ introduced in Appendix B. In all contour plots in Fig. 9 the same constant level spacing is used: $`\mathrm{\Delta }g=1/20`$ and $`\mathrm{\Delta }B_z=1/10`$ in units where $`j_c(0)d=a=\mu _0=1`$. In our computation based on Ref. we used $`|dH/dt|=1`$ and $`E_0=1`$ and $`n=51`$ in $`E=E_0(J/J_c)^n`$. The shape of the depicted loops is very similar to their shape in the isotropic plate, i.e. different $`J_c(B_z)`$ dependences lead to almost the same penetration pattern. This is so since the shape of these loops is determined mainly by the specimen shape: At small $`HJ_c(0)`$, one has almost complete field expulsion and the stream lines are thus independent of any material property, being rectangular near the edges and circular near the center of the rectangle. After the flux front has passed through a given point, $`J`$ abruptly tends to nearly a constant, see Fig. 3 and Eq. (17). The situation then becomes similar to the isotropic case, and the stream lines practically coincide with concentric rectangles with constant spacing proportional to $`1/J`$.
Since the current stream lines in the rectangular platelet do not coincide with the lines $`B_z=`$const, and $`J`$ depends on $`B_z`$, the penetration pattern in the anisotropic case could differ from that obtained in the isotropic superconductor. However, even a close look can hardly discover a qualitative difference in the shapes of the contours in isotropic and strongly out-of-plane anisotropic (Fig. 9) rectangular plates, in stark contrast to the situation with in-plane-anisotropy. A minor difference is that in the partly penetrated state ($`H=0.4`$ and $`H=0.6`$ in Fig. 9) the stream lines outside the flux front and near the axes $`x=0`$ and $`y=0`$, are slightly convex towards the center, while for isotropic Bean pinning the stream lines are straight there (or are slightly concave due to the finite creep exponent $`n`$). As expected from the strip results, Fig. 3, with out-of-plane pinning the density of the sheet-current stream-lines near the flux front is larger (when $`p>1`$) or smaller (when $`p<1`$) as compared to the case of isotropic pinning ($`p=0`$), cf. Eq. (17). The same statement is true for the contours of $`B_z`$ near the front.
Figure 10 shows magnetization curves of thin rectangular plates with side ratio 1.4 and various out-of-plane anisotropies. These magnetization loops exhibit central peaks or dips caused by the anisotropy and are very similar to those shown in Fig. 7 for strips. Thus, one can conclude that the features of the loops at $`Hj_cd`$ are specified mainly by the anisotropy of flux-line pinning (i.e., by $`p`$ and $`q`$) rather than by the shape of the sample. This property should be useful for determining these pinning parameters from experiments.
## VI Conclusions
In many papers thin superconductors are considered as 2D samples with a constant critical sheet current $`J_c`$ or with $`J_c(B_z)=j_c(B_z)d`$ if $`j_c`$ depends on the magnitude of the magnetic induction. However, when an out-of-plane anisotropy of pinning exists, the current density also depends on the angle $`\theta `$ between $`𝐁`$ (the flux-line direction) and the normal to the sample plane (our $`z`$ axis), and thus the critical state problem becomes three-dimensional. We have shown here that the general 3D critical state problem for thin superconductors of arbitrary shape and with arbitrary 3D anisotropy of pinning can be separated into a 1D problem, which treats the current density and induction across the thickness of the sample, and the 2D problem of thin superconductors with a new induction dependence of the critical sheet current $`J_c(B_z)`$ flowing in the plane of the superconductor. The new $`J_c(B_z)`$ dependence is determined by the out-of-plane anisotropy of the critical current density $`j_c(\theta )`$ and modifies the original $`B_z`$ dependence and in-plane anisotropy of $`J_c`$, if they exist. The resulting 2D thin film problem with given $`J_c(B_z)`$ can be solved by standard numerical and analytical methods. Our theory generalizes the approach of Ref. for a circular disk.
Knowing the solutions of these 1D and 2D problems, the critical state in three-dimensional thin flat superconductors (Fig. 1) can be understood. Two qualitatively new features of the critical state occur in this general case as compared to the highly symmetric cases of isotropic disks or strips, which are commonly considered. First, when the shape of the sample differs from a strip or a disk then a rotation of flux lines (differently oriented in different planes $`z=`$ const) may occur in the region where $`B_z=0`$, i.e. where the flux lines surround a non-penetrated flat core and are nearly parallel to the sample plane since the sample is thin. Second, if both out-of-plane and in-plane anisotropies of pinning exist, the flux lines even in thin superconductors are not only curved but also twisted in the region where $`B_z0`$, i.e. outside the front line $`\gamma _0`$ in Fig. 1. Twist means that the in-plane component of the direction of the flux lines at a given point $`x,y`$ changes with the depth $`z`$.
Using the above approach, the critical states of strips and rectangular platelets with only out-of-plane anisotropy are analyzed. We consider a rather general analytical model (18) for this anisotropy, which can simulate both the intrinsic pinning by CuO planes and pinning by extended defects aligned with the $`c`$ axis in high-$`T_c`$ superconductors. We obtain the shape of the flux front and the two-dimensional distributions of the magnetic induction and current density for the anisotropic strip. In the case of intrinsic pinning, the flux front has a wide flat part near the central plane, Fig. 5; the current density peaks near this flat front, and the induction profile is very steep there, see Fig. 3. In fact, a so-called current string occurs in this case, i.e., an additional current flows along this steep front as compared to the isotropic situation. In Refs. the string was considered in connection with a finite lower critical field $`H_{c1}`$. However, in the considered case the string is not due to a finite $`H_{c1}`$ but results from the anisotropy of pinning, and the jump of $`B_z`$ at the flux front considerably exceeds $`H_{c1}`$ since $`H_{cr}=j_cd/2H_{c1}`$ was assumed above. This type of out-of-plane anisotropy leads to a central peak in the magnetization curve, Fig. 7.
A different situation occurs when the $`\theta `$ dependence of $`j_c`$ has a minimum at $`\theta =\pi /2`$. In this case the flux front has the shape of a wedge with a rounded point that becomes sharper with increasing anisotropy. As can be seen in Figs. 3b,c the current density near the front line $`\gamma _0`$ now has no peak and no sharp step (as for isotropic pinning) but it decreases monotonically and has only a rounded step and an inflection point with vertical slope at this front line. The induction $`B_z`$ now vanishes less steeply than in the isotropic case, and in the limit of large such anisotropy the induction profile near the front decreases almost linearly. The magnetization loop for this type of anisotropy has a central dip, see Fig. 7.
All these features remain qualitatively the same when a smaller creep exponent $`n`$ is chosen in $`EJ^n`$, Eq. (B6), or when the superconductor has a rectangular shape. In particular, with smaller $`n`$ the cusp in $`J(x)`$ remains sharp and the profile of $`B_z(x)`$ remains steep, see Fig. 3d with $`n=11`$. As shown in Figs. 7 and 10, the magnetization loops of thin rectangular plates with out-of-plane anisotropy exhibit the same central peak or dip as thin strips.
For any given dependence $`j_c(𝐁)`$, the equations of this paper enable one to compute various characteristics of the critical state in thin superconducting samples of realistic shapes. With this theory at hand, experimental investigation of flux-density profiles, of the $`H`$ dependence of the penetrating flux front (e.g., by magneto-optics), and of magnetization loops, will yield information not only on the strengh but also anisotropy of flux-line pinning in superconductors.
###### Acknowledgements.
G.P.M. acknowledges the hospitality of the Max-Planck-Institute für Metallforschung, Stuttgart.
## A electric field in anisotropic superconductors
In a uniform anisotropic superconductor the direction of the electric field $`𝐄`$ generated by vortex motion in general will not coincide with the direction of the current. To study this effect, we start with the situation when a current with the density $`j`$ flows in the plane normal to the magnetic induction $`𝐁`$ (below we call this plane the $`N`$ plane). The creep activation barrier $`U[B,\phi ,j/j_0(B,\phi )]`$ is assumed to vanish at $`j=j_0(B,\phi )`$ \[this is the definition of the factor $`j_0(B,\phi )`$\]. Here the angle $`\phi `$ specifies the direction of the current in the $`N`$ plane, $`𝐣(\mathrm{cos}\phi ,\mathrm{sin}\phi )`$, and $`U[B,\phi ,j/j_0(B,\phi )]`$ is the value of the barrier which a flux bundle has to overcome to hop in the direction of the Lorentz force, i.e., along the vector $`[𝐣\times 𝐁]`$. At given temperature $`T`$, the expression $`\mathrm{exp}\{U[B,\zeta ,j\mathrm{cos}(\zeta \phi )/j_0(B,\zeta )]/T\}`$ gives the probability for a flux-bundle to jump in the direction normal to any given in-plane vector $`(\mathrm{cos}\zeta ,\mathrm{sin}\zeta )`$. Hence, the electric field $`𝐄`$ is obtained by averaging this expression over all angles $`\zeta `$ in the interval $`\phi \pi /2\zeta \phi +\pi /2`$ (we disregard flux-bundle jumps against the Lorentz force). If the effective depth of the flux-pinning well $`[U(B,\phi ,x)/x]_{x=1}`$ considerably exceeds $`T`$, the direction of $`𝐄`$ in the $`N`$ plane, $`(\mathrm{cos}\zeta _0,\mathrm{sin}\zeta _0)`$, can be found as the position of the minimum of $`U[B,\zeta ,j\mathrm{cos}(\zeta \phi )/j_0(B,\zeta )]`$ with respect to variations of $`\zeta `$. Since, according to the definition of $`j_0(B,\zeta )`$, one has $`U(B,\zeta ,1)/\zeta =0`$, one obtains the following formula for $`\zeta _0`$ in the critical state:
$`\mathrm{tan}(\zeta _0\phi )={\displaystyle \frac{j_0^{}(B,\zeta _0)}{j_0(B,\zeta _0)}},`$
where $`j_0^{}(B,\zeta )j_0(B,\zeta )/\zeta `$. As for the critical current density $`j_c(B,\phi )`$, we have
$`j_c(B,\phi )={\displaystyle \frac{j_0(B,\zeta _0)}{\mathrm{cos}(\zeta _0\phi )}}=[j_0(B,\zeta _0)^2+j_0^{}(B,\zeta _0)^2]^{1/2}.`$
It is also useful to express the angle $`\zeta _0`$ in terms of $`j_c(B,\phi )`$, i.e., the quantity determined in experiments,
$$\delta \zeta _0\phi =\mathrm{arctan}\left\{\frac{j_c^{}(B,\phi )}{j_c(B,\phi )}\right\},$$
(A1)
where $`j_c^{}(B,\phi )=j_c(B,\phi )/\phi `$. The above formulas have been derived under the assumption that $`\mathrm{cos}(\zeta \phi )/j_0(B_z,\zeta )`$ has only one maximum in the interval $`\phi \pi /2<\zeta <\phi +\pi /2`$. This situation does not always occur in twinned crystals. In such crystals the angle $`\zeta _0`$ may have a fixed value in some interval of angles $`\phi `$, and thus $`j_c(\phi )=j_c(\zeta _0)/\mathrm{cos}(\zeta _0\phi )`$ inside the interval. However, note that Eq. (A1) is still valid in this case.
We now present the formulas for the case considered in Sec. II: critical currents flow in the $`xy`$ plane with $`𝐣_c=(j_c\mathrm{cos}\phi ,j_c\mathrm{sin}\phi )`$, and a flux-line element has the direction $`(\mathrm{sin}\theta \mathrm{cos}\psi ,\mathrm{sin}\theta \mathrm{sin}\psi ,\mathrm{cos}\theta )`$. Since in the absence of flux cutting only the component of $`𝐣_c`$ normal to $`𝐁`$ is essential, the critical current density in the $`N`$ plane, $`j_c^{}`$, has the form
$$j_c^{}(B_z,B_t,\phi \psi ,\psi )=n_{}(\phi )j_c(B_z,B_t,\phi \psi ,\psi ),$$
(A2)
where $`n_{}(\phi )=[1\mathrm{cos}^2(\phi \psi )\mathrm{sin}^2\theta ]^{1/2}`$. It follows from Eq. (A1) that the angle $`\delta `$ between $`𝐄`$ and $`𝐣_c^{}`$ is given by the formula
$$\mathrm{tan}\delta =\frac{[\mathrm{ln}j_c^{}(B_z,B_t,\phi \psi ,\psi )]}{\phi }\frac{n_{}^2(\phi )}{\mathrm{cos}\theta }.$$
(A3)
The unit vector $`𝐞`$ directed along $`𝐄`$ has the components:
$`e_x={\displaystyle \frac{1}{n_{}(\xi )}}[\mathrm{cos}\xi \mathrm{sin}^2\theta \mathrm{cos}\psi \mathrm{cos}(\xi \psi )],`$ (A4)
$`e_y={\displaystyle \frac{1}{n_{}(\xi )}}[\mathrm{sin}\xi \mathrm{sin}^2\theta \mathrm{sin}\psi \mathrm{cos}(\xi \psi )],`$ (A5)
$`e_z={\displaystyle \frac{1}{n_{}(\xi )}}\mathrm{sin}\theta \mathrm{cos}\theta \mathrm{cos}(\xi \psi ),`$ (A6)
where the angle $`\xi `$ is defined by the expression
$$\mathrm{tan}(\xi \psi )=\mathrm{cos}\theta \frac{\mathrm{tan}(\phi \psi )+\mathrm{cos}\theta \mathrm{tan}\delta }{\mathrm{cos}\theta \mathrm{tan}\delta \mathrm{tan}(\phi \psi )}.$$
(A7)
In other words, e is the normalized projection of the vector $`(\mathrm{cos}\xi ,\mathrm{sin}\xi )`$ onto the $`N`$ plane. This explains the meaning of the angle $`\xi `$.
It is important to keep in mind one special case. If an in-plane anisotropy of flux-line pinning is absent, it follows from symmetry considerations that at fixed $`B_z`$, $`B_t`$, $`\psi `$ the derivatives $`j_c^{}/\phi `$ and $`j_c/\phi `$ are equal to zero when $`\phi \psi =\pm \pi /2`$. Thus, in this case one has $`\delta =0`$, $`\xi =\phi `$, and the directions of $`𝐄`$, $`𝐣_c`$ and $`𝐣_c^{}`$ coincide.
## B dynamic equations
It is convenient to express the sheet current in the film through a scalar function $`g(x,y)=g(𝐫)`$ as
$`𝐉(𝐫)=\times \widehat{𝐳}g(𝐫)=\widehat{𝐳}\times g(𝐫).`$
This substitution guarantees that $`𝐉=0`$ and that the current flows along the specimen boundary $`\mathrm{\Gamma }`$ if one puts $`g(𝐫)=\mathrm{const}=0`$ there (the lines $`g(𝐫)=\mathrm{const}`$ coincide with the current stream lines). Then, Eq. (14) is transformed as follows,
$$\frac{B_z(𝐫)}{\mu _0}=H+C(𝐫)g(𝐫)_S\frac{g(𝐫^{})g(𝐫)}{4\pi R^3}d^2r^{},$$
(B1)
with
$$4\pi C(𝐫)=_{\overline{S}}\frac{d^2r^{}}{R^3}=_0^{2\pi }\frac{d\varphi }{R(\varphi )}.$$
(B2)
The integral in Eq. (B1) is taken in the Cauchy sense. In formulas (B2) the integration area $`\overline{S}`$ is the entire $`x,y`$ plane with exclusion of the sample area $`S`$; $`R(\varphi )|𝐬𝐫|`$; $`𝐬`$ defines the position of a point on the boundary $`\mathrm{\Gamma }`$; and $`\varphi `$ is the angle of the vector $`𝐬𝐫`$ relative to an arbitrary fixed direction. Eq. (B1) admits inversion,
$$g(𝐫)=_SK(𝐫,𝐫^{})[\mu _0^1B_z(𝐫^{})H]d^2r^{},$$
(B3)
where the kernel $`K(𝐫,𝐫_1)`$ is found from the following equation,
$`K(𝐫,𝐫_1)={\displaystyle \frac{1}{C(𝐫)}}\{\delta (𝐫𝐫_1)`$ (B4)
$`+{\displaystyle _S}{\displaystyle \frac{K(𝐫^{},𝐫_1)K(𝐫,𝐫_1)}{4\pi R^3}}d^2r^{}\},`$ (B5)
which can be solved by iteration starting with $`K(𝐫,𝐫_1)=0`$.
We now differentiate Eq. (B3) with respect to time $`t`$, take into account the Maxwell equation, $`B_z/t=\widehat{𝐳}[\times 𝐄]`$, and the current-voltage law,
$`𝐄=𝐄(𝐉,B_z),`$
which expresses the electric field $`𝐄`$ through the sheet current $`𝐉`$. Then, we arrive at an equation for $`g(𝐫)`$,
$`{\displaystyle \frac{g(𝐫,t)}{t}}={\displaystyle _S}K(𝐫,𝐫^{})`$ (B6)
$`\left\{[\widehat{𝐳}\times ]𝐄[\widehat{𝐳}\times g(𝐫^{},t),B_z(𝐫^{})]{\displaystyle \frac{H(t)}{t}}\right\}d^2r^{}.`$ (B7)
Equations (B1) and (B5) allow one to find $`B_z(𝐫,t)`$, $`g(𝐫,t)`$ \[and thus $`𝐉(𝐫,t)`$\] if the prehistory $`H(t)`$ and the current-voltage law are given.
The static critical state can be calculated by a dynamic approach using any model law $`𝐄`$ which has a sharp bend at $`J=J_c`$. The specific form of this dependence is irrelevant; e.g. we may use the power law
$$𝐄=𝐄_0\left(\frac{|𝐉|}{J_c}\right)^n$$
(B8)
and take the limit $`n1`$. That the static critical state is indeed reached in the limit $`n\mathrm{}`$ was shown rigorously in Ref. . The function $`J(B_z,\phi _0)`$ resulting from Eq. (12) or from the appropriate generalization of this equation should be used as the critical sheet current $`J_c`$ in Eq. (B6) for currents with direction $`(\mathrm{cos}\phi _0,\mathrm{sin}\phi _0)`$. The direction of the constant vector $`𝐄_0`$ in general does not coincide with $`𝐉`$ if there is anisotropy of flux-line pinning in the $`xy`$ plane (see Appendix A). At given $`\phi _0`$, the angle $`\xi _0`$ defining the direction of $`𝐄(\mathrm{cos}\xi _0,\mathrm{sin}\xi _0)`$ is determined by the formula
$$\mathrm{tan}(\xi _0\phi _0)=\frac{[\mathrm{ln}J(B_z,\phi _0)]}{\phi _0}.$$
(B9)
Thus, to describe the critical state of an infinitely thin superconductor, it is sufficient to solve the dynamical Eqs. (B1,B5,B6) with the initial condition $`g(x,y)|_{t=0}=0`$. This dynamical approach appears to be more convenient than solving the static Eqs. (13,14) with an initially unknown flux-front $`\gamma _0`$.
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# Untitled Document
Rectangular Scott-type Permanents
Guo-Niu Han et Christian Krattenthaler
Abstract.— Let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of a polynomial $`P(x)`$ of degree $`n`$ and $`y_1,y_2,\mathrm{},y_m`$ be the zeroes of another polynomial $`Q(y)`$ of degree $`m`$. Our object of study is the permanent $`\text{per}(1/(x_iy_j))_{1in,\mathrm{\hspace{0.17em}1}jm}`$, here named “Scott-type” permanent, the case of $`P(x)=x^n1`$ and $`Q(y)=y^n+1`$ having been considered by R. F. Scott. We present an efficient approach to determining explicit evaluations of Scott-type permanents, based on generalizations of classical theorems by Cauchy and Borchardt, and of a recent theorem by Lascoux. This continues and extends the work initiated by the first author (“Généralisation de l’identité de Scott sur les permanents,” to appear in Linear Algebra Appl.). Our approach enables us to provide numerous closed form evaluations of Scott-type permanents for special choices of the polynomials $`P(x)`$ and $`Q(y)`$, including generalizations of all the results from the above mentioned paper and of Scott’s permanent itself. For example, we prove that if $`P(x)=x^n1`$ and $`Q(y)=y^{2n}+y^n+1`$ then the corresponding Scott-type permanent is equal to $`(1)^{n+1}n!`$.
Résumé.—Soient $`x_1,x_2,\mathrm{},x_n`$ les zéros d’un polynôme $`P(x)`$ de degré $`n`$ et $`y_1,y_2,\mathrm{},y_m`$ les zéros d’un autre polynôme $`Q(y)`$ de degré $`m`$. Notre objet d’étude est le permanent $`\text{per}(1/(x_iy_j))_{1in,\mathrm{\hspace{0.17em}1}jm}`$, appelé ici permanent de type Scott. Le cas de $`P(x)=x^n1`$ et $`Q(y)=y^n+1`$ a été considéré par R. F. Scott. Nous présentons une approche efficace pour déterminer les évaluations explicites des permanents de type Scott, basée sur des généralisations des théorèmes classiques de Cauchy et Borchardt, et d’un théorème récent de Lascoux. La présente étude prolonge le travail du premier auteur (“Généralisation de l’identité de Scott sur les permanents,” à apparaître dans Linear Algebra Appl.). Notre approche nous permet de fournir de nombreuses évaluations explicites des permanents de type Scott pour des choix spéciaux des polynômes $`P(x)`$ et $`Q(y)`$, y compris des généalisations de tous les résultats de l’article mentionné ci-dessus et du permanent de Scott lui-même. Par exemple, nous prouvons que si $`P(x)=x^n1`$ et $`Q(y)=y^{2n}+y^n+1`$ alors le permanent correspondant de type Scott est égal à $`(1)^{n+1}n!`$.
1. Introduction
In 1881, Scott stated, without proof, the following result:
Let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of $`x^n1`$ and $`y_1,y_2,\mathrm{},y_n`$ be the zeroes of $`y^n+1`$. Let $`A`$ be the $`n\times n`$ matrix $`(1/(x_iy_j))_{1i,jn}`$. Then
$$\text{per}(A)=\{\begin{array}{cc}(1)^{\frac{n1}{2}}\frac{n(135\mathrm{}(n2))^2}{2^n},\hfill & \text{if }n\text{ is odd,}\hfill \\ 0,\hfill & \text{if }n\text{ is even.}\hfill \end{array}$$
In 1978, in his monograph Permanents, Minc \[13, p. 155\] included this result in a list of conjectures on permanents. Since then, several proofs have been given , one of which by Minc himself. All of these proofs are heavily based on the fact that the zeroes of the polynomials $`x^n1`$ and $`y^n+1`$ can be written in simple explicit terms. Thus, neither of these proofs extends to the more general problem which is the subject of this paper:
Given a polynomial $`P(x)`$ of degree $`n`$ and a polynomial $`Q(y)`$ of degree $`m`$, let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of $`P(x)`$ and $`y_1,y_2,\mathrm{},y_m`$ be the zeroes of $`Q(y)`$. Evaluate the permanent of the $`n\times m`$ matrix $`(1/(x_iy_j))_{1in,\mathrm{\hspace{0.17em}1}jm}`$.
As usual, the permanent $`\text{per}(A)`$ of an $`n\times m`$ matrix $`A`$ is defined as the sum of all possible products of $`n`$ coefficients of $`A`$ chosen such that no two of the coefficients are taken from the same row nor from the same column (see \[13, Ch. 1, (1.1)\]). Given the assumptions of the above problem, we call the permanent of the matrix $`(1/(x_iy_j))_{1in,\mathrm{\hspace{0.17em}1}jm}`$ a Scott-type permanent, and denote it by $`\text{PER}(P(x),Q(y))`$.
In , the first author presented a new approach to this type of problem in the case $`n=m`$, i.e., in the case that both polynomials have the same degree. This approach does not rely at all on explicit analytic forms of zeroes of polynomials. Instead, it makes essential use of recent symmetric functions techniques, in particular of a theorem due to Lascoux , which the latter author established in his étude on the square ice model of statistical mechanics.
In the present paper, we are going to extend this approach to arbitrary $`n`$ and $`m`$. This requires extensions of classical theorems of Cauchy and Borchardt (see Theorems (Cauchy+) and (Borchardt+) in Section 2), and an extension of Lascoux’s theorem (see Theorem (Lascoux+)). As a result (see Theorem 1), we are able to express any Scott-type permanent as the quotient of a determinant which features complete homogeneous and elementary symmetric functions in the zeroes of the two polynomials, divided by the resultant of the two polynomials. In particular, it follows immediately that any Scott-type permanent is rational in the coefficients of the polynomials $`P(x)`$ and $`Q(y)`$.
In Section 5 we apply this result to obtain explicit evaluations of Scott-type permanents in numerous special cases. Amongst others, we provide generalizations of all the results from , thus also covering Scott’s permanent itself. For the proofs of the results in Section 5, we make use of two particular specializations of our main theorem, Theorem 1, which we derive in Section 3 (see Theorems 3 and 4), and of four determinant evaluations, which we state and establish separately in Section 4.
Finally, we also comment briefly on an alternative approach to the evaluation of Scott-type determinants, due to the first author . It allows to express Scott-type permanents in terms of weighted sums over involutions (see Theorem 2 in Section 2). By combining this result with some of the evaluations in Section 5, we obtain interesting summation theorems, which are presented in Section 6.
2. The general theory
In , the main ingredients are theorems by Cauchy, Borchardt, Lascoux, and a lemma on the resultant. Since we intend to extend the approach of to the case of rectangular Scott-type permanents (corresponding to polynomials of, possibly, different degrees), we have to first provide the appropriate extensions of these theorems.
Given positive integers $`m`$ and $`n`$ and two sets $`X=\{x_1,x_2,\mathrm{},x_n\}`$ and $`Y=\{y_1,y_2,\mathrm{},y_m\}`$ of variables, we use the following notations:
$$R(X,Y):=\underset{i=1}{\overset{n}{}}\underset{j=1}{\overset{m}{}}(x_iy_j)\text{ and }\mathrm{\Delta }(X):=\underset{i<j}{}(x_ix_j).$$
$$\left(\frac{1}{x_iy_j}\right):=\left(\begin{array}{cccc}\frac{1}{x_1y_1}& \frac{1}{x_1y_2}& \mathrm{}& \frac{1}{x_1y_m}\\ \frac{1}{x_2y_1}& \frac{1}{x_2y_2}& \mathrm{}& \frac{1}{x_2y_m}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \\ \frac{1}{x_ny_1}& \frac{1}{x_ny_2}& \mathrm{}& \frac{1}{x_ny_m}\end{array}\right).$$
We first state a variation on Cauchy’s evaluation of his double alternant (cf. \[4; 15, vol. 1, pp. 342–345\]).
Theorem (Cauchy+).— For $`mn`$, let
$$C(X,Y):=\left(\begin{array}{cccc}\frac{1}{x_1y_1}& \frac{1}{x_1y_2}& \mathrm{}& \frac{1}{x_1y_m}\\ \frac{1}{x_2y_1}& \frac{1}{x_2y_2}& \mathrm{}& \frac{1}{x_2y_m}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \\ \frac{1}{x_ny_1}& \frac{1}{x_ny_2}& \mathrm{}& \frac{1}{x_ny_m}\\ \\ 1& 1& \mathrm{}& 1\\ y_1& y_2& \mathrm{}& y_m\\ y_1^2& y_2^2& \mathrm{}& y_m^2\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ y_1^{mn1}& y_2^{mn1}& \mathrm{}& y_m^{mn1}\end{array}\right).$$
Then
$$det(C(X,Y))=(1)^{n(n1)/2}\frac{\mathrm{\Delta }(X)\mathrm{\Delta }(Y)}{R(X,Y)}.$$
Proof.—If $`n=m`$, this is exactly Cauchy’s theorem. The general case can be either established directly, or, it may be observed that the “general” case is in fact implied by Cauchy’s theorem. To see this, consider the above identity with $`n=m`$. Given $`k<m`$, expand both sides as power series in $`1/x_{k+1}`$, …, $`1/x_m`$, and compare coefficients of $`1/x_{k+1}x_{k+2}^2\mathrm{}x_m^{mk}`$ on both sides.
Next we state the required extension of Borchardt’s theorem \[2; 15, vol. 2, pp. 173–175\]. It can be established by reading through the proof of Borchardt’s theorem given in \[1, Proof of Cor. 5.1\], ignoring however the restriction $`m=n`$ (see also ).
Theorem (Borchardt+).— For $`mn`$, let
$$B(X,Y):=\left(\begin{array}{cccc}\frac{1}{(x_1y_1)^2}& \frac{1}{(x_1y_2)^2}& \mathrm{}& \frac{1}{(x_1y_m)^2}\\ \frac{1}{(x_2y_1)^2}& \frac{1}{(x_2y_2)^2}& \mathrm{}& \frac{1}{(x_2y_m)^2}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \\ \frac{1}{(x_ny_1)^2}& \frac{1}{(x_ny_2)^2}& \mathrm{}& \frac{1}{(x_ny_m)^2}\\ \\ 1& 1& \mathrm{}& 1\\ y_1& y_2& \mathrm{}& y_m\\ y_1^2& y_2^2& \mathrm{}& y_m^2\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ y_1^{mn1}& y_2^{mn1}& \mathrm{}& y_m^{mn1}\end{array}\right).$$
Then
$$det(B(X,Y))=det(C(X,Y))\times \text{per}\left(\frac{1}{x_iy_j}\right).$$
Our next goal is to derive the required extension of (a special case of) Lascoux’s theorem \[11, Theorem q\]. Let $`Z=\{z_1,z_2,\mathrm{},z_n\}`$ be another set of variables, of equal cardinality as $`X`$. For $`1in`$ we define the divided difference $`_i`$ by
$$_i:f\frac{ff^{\sigma _i}}{x_iz_i},$$
where $`\sigma _i`$ is the transposition which interchanges $`x_i`$ and $`z_i`$. It is easy to see that
$$_i\frac{1}{z_iy}=\frac{1}{(x_iy)(z_iy)}.$$
Since the operator $`_i`$ acts only on one row of the matrix $`C(Z,Y)`$ (to be precise, the $`i`$-th row), it follows that
$$_1_2\mathrm{}_n(det(C(Z,Y))|_{Z=X}=det(B(X,Y)).$$
$`(1)`$
Now, to generalize Lascoux’s theorem to our case, one simply readsthrough the proof of Theorem q in , on introducing slight modifications if necessary. The result is:
Theorem (Lascoux+).— Let $`H(X)`$ be the $`n\times (m+n1)`$ matrix defined by
$$H(X):=\left(h_{ji}(X)\right)_{1in,1jm+n1},$$
where $`h_s(X)`$ denotes the complete homogeneous symmetric function of degree $`s`$ in the variables $`X`$ (cf. \[12, Ch. 1\]), and $`E(Y)`$ be the $`(m+n1)\times n`$ matrix defined by
$$E(Y):=\left((j2k+2)(1)^{mj+k1}e_{mj+k1}(Y)\right)_{1jm+n1,1kn},$$
where $`e_s(Y)`$ denotes the elementary symmetric function of degree $`s`$ in the variables $`Y`$ (cf. \[12, Ch. 1\]). Then
$$_1_2\mathrm{}_n(\mathrm{\Delta }(Z)R(X,Y))|_{Z=X}=\mathrm{\Delta }(X)det\left(H(X)E(Y)\right).$$
Now we are in the position to state our main theorem, which will enable us, in Section 5, to evaluate numerous Scott-type permanents in closed form. The theorem implies immediately that any Scott-type permanent $`\text{PER}(P(x),Q(y))`$ is a rational function in the coefficients of the two polynomials $`P(x)`$ and $`Q(y)`$.
Theorem 1.— Let $`m`$ and $`n`$ be arbitrary positive integers, and let $`X=\{x_1,x_2,\mathrm{},x_n\}`$ and $`Y=\{y_1,y_2,\mathrm{},y_m\}`$ be two sets of variables. Then
$$\text{per}\left(\frac{1}{x_iy_j}\right)=\frac{det\left(H(X)E(Y)\right)}{R(X,Y)}.$$
Proof.—First let $`mn`$. By combining Theorems (Cauchy+) and (Borchardt+), and Equation (1), we obtain
$$(1)^{n(n1)/2}\frac{R(X,Y)}{\mathrm{\Delta }(X)\mathrm{\Delta }(Y)}_1_2\mathrm{}_n(det(C(Z,Y))|_{Z=X}$$
for the permanent. Now we apply Theorem (Cauchy+) again in order to replace $`C(Z,Y)`$ by the corresponding product form guaranteed by the theorem. After having used that the divided differences $`_i`$ commute with $`\mathrm{\Delta }(Y)/R(X,Y)R(Z,Y)`$ (because the latter expression is symmetric in $`x_i`$ and $`z_i`$), Theorem (Lascoux+) applies and yields the desired result.
If $`m<n`$, the permanent clearly vanishes. According to Theorem (Lascoux+), it suffices to establish that
$$U:=_1_2\mathrm{}_n(\mathrm{\Delta }(Z)R(X,Y))=0.$$
To begin with, we rewrite the Vandermonde determinant evaluation as $`\mathrm{\Delta }(Z)R(X,Y)=(1)^{n(n1)/2}det(z_i^{j1}R(x_i,Y))`$. Writing $`Q(x_i):=R(x_i,Y)=a_mx_i^m+\mathrm{}+a_1x_i+a_0`$, we obtain for $`U`$ the expression
$$U=det(_iz_i^{j1}R(x_i,Y))=det\left(z_i^{j1}Q(x_i)x_i^{j1}Q(z_i)\right)/(x_iz_i).$$
Since for $`m<n`$, we have
$$\underset{j=0}{\overset{m}{}}a_m\left(z_i^jQ(x_i)x_i^jQ(z_i)\right)=0,$$
the $`m+1`$ elements $`z_i^jQ(x_i)x_i^jQ(z_i)`$, $`0jm`$, are linearly dependent. Hence, $`U=0`$.
In , the first author obtained another expression for the permanent, in form of a certain weighted sum over involutions. To state and explain this formula, for $`sX`$ define
$$L(s;X,Y):=\underset{xs}{}\frac{1}{xs}+\underset{yY}{}\frac{1}{sy}.$$
Let us denote by $`(n)`$ the set of involutions on $`\{1,2,\mathrm{},n\}`$. Given an involution $`\sigma (n)`$, we define the weight $`\mathrm{\Psi }(\sigma )`$ of $`\sigma `$ by
$$\mathrm{\Psi }(\sigma ;X,Y):=\underset{(ij)\sigma }{}\frac{1}{(x_ix_j)^2}\underset{(k)\sigma }{}L(x_k;X,Y)$$
where the first product is over all transpositions $`(ij)`$ in the disjoint cycle decomposition of $`\sigma `$, and where the second product is over all fixed points $`k`$ of $`\sigma `$. Then the result from is the following.
Theorem 2.— Let $`m`$ and $`n`$ be arbitrary positive integers, and let $`X=\{x_1,x_2,\mathrm{},x_n\}`$ and $`Y=\{y_1,y_2,\mathrm{},y_m\}`$ be two sets of variables. Then
$$\text{per}\left(\frac{1}{x_iy_j}\right)=\underset{\sigma (n)}{}\mathrm{\Psi }(\sigma ;X,Y).$$
Example 1.— Let $`n=1`$, $`X=\{x\}`$, $`Y=\{y_1,y_2,\mathrm{},y_m\}`$. Then
$$\begin{array}{cc}\hfill H(X)& :=\left(h_{j1}(X)\right)_{i=1,1jm}=(h_0(x),h_1(x),\mathrm{},h_{m1}(x)),\hfill \\ \hfill E(Y)& :=\left(j(1)^{mj}e_{mj}(Y)\right)_{1jm,k=1}\hfill \\ & =(1(1)^{m1}e_{m1}(Y),2(1)^{m2}e_{m2}(Y),\mathrm{},me_0(Y))^t.\hfill \end{array}$$
We have
$$H(X)E(Y)=\underset{j=0}{\overset{m1}{}}(j+1)(1)^{mj1}e_{mj1}(Y)x^j.$$
Therefore,
$$\begin{array}{cc}\hfill \text{per}\left(\frac{1}{xy}\right)& =\frac{1}{xy_1}+\frac{1}{xy_2}+\mathrm{}+\frac{1}{xy_m}\text{[Definition, Th. 2]}\hfill \\ & =\frac{\underset{j=0}{\overset{m1}{}}(j+1)(1)^{mj1}e_{mj1}(Y)x^j}{(xy_1)(xy_2)\mathrm{}(xy_m)}.\text{[Th. 1]}\hfill \end{array}$$
Example 2.— For $`n=4`$ and $`m=3`$, Theorem 1 yields the following identity:
$$det\left(\left(\begin{array}{cccccc}h_0& h_1\hfill & h_2\hfill & h_3\hfill & h_4\hfill & h_5\hfill \\ 0& h_0\hfill & h_1\hfill & h_2\hfill & h_3\hfill & h_4\hfill \\ 0& 0\hfill & h_0\hfill & h_1\hfill & h_2\hfill & h_3\hfill \\ 0& 0\hfill & 0\hfill & h_0\hfill & h_1\hfill & h_2\hfill \end{array}\right)\times \left(\begin{array}{cccc}e_2& e_3\hfill & 0\hfill & 0\hfill \\ 2e_1& 0\hfill & 2e_3\hfill & 0\hfill \\ 3e_0& e_1\hfill & e_2\hfill & 3e_3\hfill \\ 0& 2e_0\hfill & 0\hfill & 2e_2\hfill \\ 0& 0\hfill & e_0\hfill & e_1\hfill \\ 0& 0\hfill & 0\hfill & 0\hfill \end{array}\right)\right)=0.$$
For $`n=2`$ and $`m=1`$, Theorem 2 yields
$$\frac{1}{(x_1x_2)^2}+\left(\frac{1}{x_2x_1}+\frac{1}{x_1y}\right)\left(\frac{1}{x_1x_2}+\frac{1}{x_2y}\right)=0.$$
3. The case of $`P(x)=x^n1`$ and of $`P(x)=x^{n1}+\mathrm{}+x+1`$
In this section we specialize Theorem 1 to the case that the $`x_i`$’s are the zeroes of the polynomial $`P(x)=x^n1`$ or of $`P(x)+x^{n1}+\mathrm{}+x+1`$, and the $`y_i`$’s are the zeroes of an arbitrary other polynomial. (This covers, for example, the case of Scott’s identity). For the remainder of this section, we fix $`m`$ and $`n`$, $`mn`$.
Let $`X=\{x_1,x_2,\mathrm{},x_n\}`$ be the set of zeroes of $`x^n1`$, and let $`Y=\{y_1,y_2,\mathrm{},y_m\}`$ be the set of zeroes of $`Q(x)=a_my^m+a_{m1}y^{m1}+\mathrm{}+a_1y^1+a_0`$, with $`a_m=1`$. We write
$$\text{PER}(P,Q):=\text{per}\left(\frac{1}{x_iy_j}\right)=\frac{det(H(X)E(Y))}{R(X,Y)}.$$
$`(2)`$
Since
$$h_i(X)t^i=\frac{1}{_i(1tx_i)}=\frac{1}{(1t^n)}=1+t^n+t^{2n}+\mathrm{},$$
we have
$$h_k(X)=\{\begin{array}{cc}1,\hfill & \text{if }k=0(\text{mod }n)\text{,}\hfill \\ 0,\hfill & \text{if }k0(\text{mod }n)\text{.}\hfill \end{array}$$
We denote by $`\text{I}_k`$ the $`k\times k`$ identity matrix, and by $`\text{0}_{l,c}`$ the $`l\times c`$ matrix with all entries equal to $`0`$. For all $`r`$, we write $`r\%n`$ for the number between $`1`$ and $`n`$ that satisfies $`r(\text{mod }n)=r\%n(\text{mod }n)`$. Then we have
$$H(X)=\left(\begin{array}{ccccc}\text{I}_n& \text{I}_n& \mathrm{}& \text{I}_n& \begin{array}{c}\text{I}_m^{}\\ \text{0}_{nm^{},m^{}}\end{array}\end{array}\right),$$
with $`m^{}=(m+n1)\%n`$.
Furthermore, let $`\mathrm{diag}_n^i(c_1,c_2,\mathrm{},c_n)`$ denote the $`n\times n`$ “diagonal” matrix, in which the (broken) diagonal starts in the $`i`$-th row,
$$\left(\begin{array}{cc}\text{0}_{i1,ni+1}& \mathrm{diag}_{i1}(c_{ni+2},c_{ni+3},\mathrm{},c_n)\\ \mathrm{diag}_{ni+1}(c_1,c_2,\mathrm{},c_{ni+1})& \text{0}_{ni+1,i1}\end{array}\right).$$
According to the definition of $`E(Y)`$, a simple calculation yields that
$$H(X)E(Y)=\underset{r=0}{\overset{m}{}}\mathrm{diag}_n^{r\%n}(ra_r,(r1)a_r,\mathrm{},(rn+1)a_r).$$
Example 3.— For $`n=3`$ and $`m=4`$, let $`P(x)=x^31`$ and $`Q(y)=y^4+a_3y^3+a_2y^2+a_1y+a_0`$. Then $`H(X)E(Y)`$ is the sum of the following matrices:
$$\left(\begin{array}{ccc}0& a_0\hfill & 0\hfill \\ 0& 0\hfill & 2a_0\hfill \\ 0& 0\hfill & 0\hfill \end{array}\right)+\left(\begin{array}{ccc}a_1& 0\hfill & 0\hfill \\ 0& 0\hfill & 0\hfill \\ 0& 0\hfill & a_1\hfill \end{array}\right)+\left(\begin{array}{ccc}0& 0\hfill & 0\hfill \\ 2a_2& 0\hfill & 0\hfill \\ 0& a_2\hfill & 0\hfill \end{array}\right)+\left(\begin{array}{ccc}0& 2a_3\hfill & 0\hfill \\ 0& 0\hfill & a_3\hfill \\ 3a_3& 0\hfill & 0\hfill \end{array}\right)+\left(\begin{array}{ccc}4& 0\hfill & 0\hfill \\ 0& 3\hfill & 0\hfill \\ 0& 0\hfill & 2\hfill \end{array}\right).$$
In the theorem below, we summarize our findings.
Theorem 3.— Let $`P(x)=x^n1`$ and $`Q(y)=a_my^m+\mathrm{}+a_1y^1+a_0`$, $`a_m`$ not necessarily $`1`$. Writing
$$\text{Fes}(Q)=det\left(\underset{r=0}{\overset{m}{}}\mathrm{diag}_n^{r\%n}(ra_r,(r1)a_r,\mathrm{},(rn+1)a_r)\right),$$
$`(3)`$
we have
$$\text{PER}(P,Q)=\frac{\text{Fes}(Q)}{\text{Res}(P,Q)},$$
where Res is the classical resultant of two polynomials.
Proof.—We have already seen that the theorem is true for $`a_m=1`$. On the other hand, there hold $`\text{Fes}(\lambda Q)=\lambda ^n\text{Fes}(Q)`$ and $`\text{Res}(P,\lambda Q)=\lambda ^n\text{Res}(P,Q)`$, as is easily verified.
Now let us consider the case that $`X=\{x_1,x_2,\mathrm{},x_{n1}\}`$ is the set of zeroes of $`P(x)=x^{n1}+\mathrm{}+x+1`$. We perform an analysis very similar to the one before, using the fact that we have
$$h_i(X)t^i=\frac{1}{_i(1tx_i)}=\frac{1t}{_i(1t^n)}=1t+t^nt^{n+1}+\mathrm{}.$$
In order to state the result, we introduce the following notation: We write $`\stackrel{𝑒}{\mathrm{diag}}{}_{n1}{}^{i}(c_1,c_2,\mathrm{},c_{n1})`$ for the $`(n1)\times (n1)`$ matrix
$$\left(\begin{array}{ccc}\text{0}_{i1,ni}& \text{0}_{i1,1}& \mathrm{diag}_{i2}(c_{ni+2},c_{ni+3},\mathrm{},c_{n1})\\ \mathrm{diag}_{ni}(c_1,c_2,\mathrm{},c_{ni})& \text{0}_{ni,1}& \text{0}_{ni+1,i2}\end{array}\right)$$
if $`i>1`$, respectively $`\mathrm{diag}_{n1}(c_1,c_2,\mathrm{},c_{n1})`$ if $`i=1`$. This is again a matrix with a (possibly broken) diagonal, in which the diagonal “jumps over” one row and column in the case that it is broken. (Note the slight discrepancy in dimension between the diagonal and the zero matrices.)
Theorem 4.— Let $`P(x)=x^{n1}+\mathrm{}+x+1`$ and $`Q(y)=a_my^m+\mathrm{}+a_1y^1+a_0`$, $`a_m`$ not necessarily $`1`$. Writing
$$\begin{array}{cc}\hfill \stackrel{𝑒}{\mathrm{Fes}}(Q)=& det(\underset{r=0}{\overset{m}{}}\stackrel{𝑒}{\mathrm{diag}}{}_{n1}{}^{r\%n}(ra_r,(r1)a_r,\mathrm{},(rn+2)a_r)\hfill \\ & \underset{r=0}{\overset{m}{}}\stackrel{𝑒}{\mathrm{diag}}{}_{n1}{}^{(r1)\%n}(ra_r,(r1)a_r,\mathrm{},(rn+2)a_r)),\hfill \end{array}$$
$`(4)`$
we have
$$\text{PER}(P,Q)=\frac{\stackrel{𝑒}{\mathrm{Fes}}(Q)}{\text{Res}(P,Q)},$$
where, again, Res is the classical resultant of two polynomials.
According to Theorems 3 and 4, for accomplishing the evaluation of the permanent, it is necessary to evaluate the numerator $`\text{Fes}(Q)`$, respectively $`\stackrel{𝑒}{\mathrm{Fes}}(Q)`$, and the denominator $`\text{Res}(P,Q)`$. For the evaluation of the resultant $`\text{Res}(P,Q)`$, we make use of the following lemma, which, for example, appears explicitly in . (In fact, it follows from a special case of Proposition 6 in the next section.)
Lemma 5.—Let $`d`$ be the greatest common divisor of $`m`$ and $`n`$, and let $`A`$ and $`C`$ be two nonzero constants. Then the resultant of the two polynomials $`Ax^mB`$ and $`Cx^nD`$ is given by
$$\text{Res}(Ax^mB,Cx^nD)=(1)^m\left(A^{n/d}D^{m/d}B^{n/d}C^{m/d}\right)^d.$$
What concerns the evaluation of $`\text{Fes}(Q)`$, respectively $`\stackrel{𝑒}{\mathrm{Fes}}(Q)`$, we refer the reader to the next section for the determinant evaluations that we are going to use. In combination with Lemma 5, these will allow us to evaluate Scott-type permanents in numerous special cases, see Section 5.
4. Determinant evaluations
Proposition 6.— Let $`n`$ and $`r`$ be positive integers, $`rn`$, and $`x_1,x_2,\mathrm{},x_n`$, $`y_1,y_2,\mathrm{},y_n`$ be indeterminates. Then, with $`d=\text{gcd}(r,n)`$, we have
$$\begin{array}{cc}& det\left(\begin{array}{ccccccc}x_1& 0& \mathrm{}& 0& y_{nr+1}& 0& \\ 0& x_2& 0& & 0& y_{nr+2}& 0\\ & & \mathrm{}& & & & \mathrm{}& 0\\ 0& & & & & & 0& y_n\\ y_1& 0& \\ 0& y_2& 0\\ & 0& \mathrm{}& 0& & & \mathrm{}& 0\\ & & 0& y_{nr}& 0& & 0& x_n\end{array}\right)\hfill \\ & =\underset{i=1}{\overset{d}{}}\left(\underset{j=1}{\overset{n/d}{}}x_{i+(j1)d}(1)^{n/d}\underset{j=1}{\overset{n/d}{}}y_{i+(j1)d}\right).\hfill \end{array}$$
$`(5)`$
(I.e., in the matrix there are only nonzero entries along two diagonals, one of which is a broken diagonal.)
Proof.—Let $`A=(A_{ij})_{1i,jn}`$ be the matrix of which the determinant is taken in (5). By definition, we have
$$detA=\underset{\sigma S_n}{}\text{sgn}\sigma \underset{i=1}{\overset{n}{}}A_{i\sigma (i)},$$
$`(6)`$
where $`S_n`$ denotes the symmetric group of order $`n`$. Since the matrix $`A`$ is a sparse matrix, only those permutations $`\sigma `$ contribute to the sum which have the property $`\sigma (i)=i`$ or $`\sigma (i)=i+r`$ mod $`n`$ for all $`i`$. Thus, the decomposition into disjoint cycles of such a permutation consists only of cycles of length $`1`$, and of cycles of length $`n/d`$ of the form $`(i,i+r,\mathrm{},i+(n/d1)r)`$ (where, again, all integers have to be taken modulo $`n`$). Using these observations and the fact that the sign of any cycle of length $`n/d`$ is $`(1)^{n/d1}`$ in (6) yields (5) immediately.
Theorem 7.— Let $`a,b,c,d,e`$ be indeterminates. For any positive integer $`n`$ and integers $`i,j`$ let $`n(i,j)`$ denote 1 plus (the representative between $`0`$ and $`n1`$ of) the residue class of $`ij`$ mod $`n`$. Then
$$\begin{array}{c}\text{ }\underset{1i,jn}{det}\left((n(i,j)+c)(n(i,j)a+b)+d(j1)(n(i,j)a+e)\right)\text{ }\hfill \\ \text{ }=det\left(\begin{array}{cccc}\left(\mathrm{c}+1\right)\left(\mathrm{a}+\mathrm{b}\right)& \left(\mathrm{c}+\mathrm{n}\right)\left(\mathrm{na}+\mathrm{b}\right)& & \left(\mathrm{c}+2\right)\left(2\mathrm{a}+\mathrm{b}\right)\\ +\mathrm{d}& +\mathrm{d}\left(\mathrm{na}+\mathrm{e}\right)& \mathrm{}& +\mathrm{d}\left(\mathrm{n}1\right)\left(2\mathrm{a}+\mathrm{e}\right)\\ \\ \left(\mathrm{c}+2\right)\left(2\mathrm{a}+\mathrm{b}\right)& \left(\mathrm{c}+1\right)\left(\mathrm{a}+\mathrm{b}\right)& & \left(\mathrm{c}+3\right)\left(3\mathrm{a}+\mathrm{b}\right)\\ +\mathrm{d}& +\mathrm{d}\left(\mathrm{a}+\mathrm{e}\right)& \mathrm{}& +\mathrm{d}\left(\mathrm{n}1\right)\left(3\mathrm{a}+\mathrm{e}\right)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \\ \left(\mathrm{c}+\mathrm{n}\right)\left(\mathrm{na}+\mathrm{b}\right)& \left(\mathrm{c}+\mathrm{n}1\right)& & \left(\mathrm{c}+1\right)\left(\mathrm{a}+\mathrm{b}\right)\\ +\mathrm{d}& \times \left(\left(\mathrm{n}1\right)\mathrm{a}+\mathrm{b}\right)+\mathrm{d}& \mathrm{}& +\mathrm{d}\left(\mathrm{n}1\right)\left(\mathrm{a}+\mathrm{e}\right)\\ & \left(\left(\mathrm{n}1\right)\mathrm{a}+\mathrm{e}\right)& & \end{array}\right)\text{ }\hfill \\ \text{ }=(n)^{n1}U_n(a,b,c,d,e)\underset{i=3}{\overset{n}{}}(ia+b+ca),(7)\text{ }\hfill \end{array}$$
where $`U_n(a,b,c,d,e)`$ is the polynomial
$$\begin{array}{cc}& \mathrm{U}_\mathrm{n}(\mathrm{a},\mathrm{b},\mathrm{c},\mathrm{d},\mathrm{e})=\frac{\left(\mathrm{n}+1\right)\left(\mathrm{n}+2\right)}{3}\mathrm{a}^2+\frac{\left(\mathrm{n}+1\right)\left(2\mathrm{n}+7\right)}{6}\mathrm{ab}+\frac{\left(\mathrm{n}+1\right)}{2}\mathrm{b}^2\hfill \\ & +\frac{\left(\mathrm{n}+1\right)\left(2\mathrm{n}+7\right)}{6}\mathrm{a}^2\mathrm{c}+\frac{\left(3\mathrm{n}+5\right)}{2}\mathrm{abc}+\mathrm{b}^2\mathrm{c}+\frac{\left(\mathrm{n}+1\right)}{2}\mathrm{a}^2\mathrm{c}^2+\mathrm{abc}^2+\frac{\left(\mathrm{n}+3\right)}{2}\mathrm{ad}\hfill \\ & +\mathrm{bd}+\mathrm{acd}\frac{\left(\mathrm{n}1\right)\left(2\mathrm{n}+5\right)}{6}\mathrm{ae}\frac{\left(\mathrm{n}1\right)}{2}\mathrm{be}\frac{\left(\mathrm{n}1\right)}{2}\mathrm{ace}.\hfill \end{array}$$
In the case that $`n=2`$, the product in (7) has to be read as $`1`$, and in the case that $`n=1`$, the product has to be interpreted as $`1/(2a+b+ca)`$.
Proof.—In the cases $`n=1`$ and $`n=2`$, the claim can be verified directly. For $`n3`$, we use the “identification of factors” method as explained in \[10, Sec. 2.4\] or \[9, Sec. 2\].
We proceed in several steps. An outline is as follows. In the first step we show that $`_{i=3}^n(ia+b+ca)`$ is a factor of the determinant as a polynomial in $`a,b,c,d,e`$. In the second step we prove that $`U_n(a,b,c,d,e)`$ is a factor of the determinant. Then, in the third step, we determine the maximal degree of the determinant as a polynomial in $`a`$, and also in $`b`$, $`c`$, $`d`$, and in $`e`$. It turns out that the maximal degree is $`n`$ as a polynomial in $`a`$, the same being true as a polynomial in $`b`$ and as a polynomial in $`c`$, while it is 1 as a polynomial in $`d`$, the same being true as a polynomial in $`e`$. On the other hand, the degree in $`a`$, and also in $`b`$ and in $`c`$, of the product on the right-hand side of (7), which by the first two steps divides the determinant, is exactly $`n`$. It is exactly 1 in $`d`$ and also in $`e`$. Therefore we are forced to conclude that the determinant equals
$$C(n)U_n(a,b,c,d,e)\underset{i=3}{\overset{n}{}}(ia+b+ca),$$
$`(8)`$
where $`C(n)`$ is a constant independent of $`a,b,c,d,e`$. Finally, in the fourth step, we determine the constant $`C(n)`$, which turns out to equal $`(n)^{n1}`$. Clearly, this would finish the proof of theorem.
Step 1. For $`i=3,\mathrm{},n`$ the term $`(ia+b+ca)`$ is a factor of the determinant. We claim that, if $`b=iaca`$, we have
$$\begin{array}{cc}& \left(\text{row }(ni)\right)+3\left(\text{row }(ni+1)\right)\hfill \\ & 3\left(\text{row }(ni+2)\right)+\left(\text{row }(ni+3)\right)=0\hfill \end{array}$$
as long as $`n>i`$. (Here, (row $`i`$) denotes the $`i`$-th row of the matrix underlying the determinant in (7).) In the case that $`n=i`$, we claim that we have
$$3\left(\text{row }1\right)+3\left(\text{row }2\right)\left(\text{row }3\right)+\left(\text{row }n\right)=0$$
as long as $`n>3`$, and that we have $`\left(\text{row }1\right)+\left(\text{row }2\right)=0`$ if $`n=3`$. All these claims are easily verified by an obvious case-by-case analysis.
Step 2. The polynomial $`U_n(a,b,c,d,e)`$ is a factor of the determinant. We claim that if $`d`$ is chosen so that $`U_n(a,b,c,d,e)`$ vanishes, we have
$$\underset{j=1}{\overset{n}{}}((j+1)a+b+ca)\left(\text{column }j\right)=0.$$
Again, it is a routine task to verify this identity.
Step 3. The determinant is a polynomial in $`a`$ (in $`b`$, respectively in $`c`$) of maximal degree $`n`$, and a polynomial in $`d`$ (respectively in $`e`$) of maximal degree $`1`$. The first claim follows from the fact that each term in the defining expansion of the determinant has degree $`n`$ in $`a`$ (as well as in $`b`$, respectively in $`c`$). To establish the second claim, we simply subtract the first row of the determinant from all other rows, with the effect that only the entries in the first row contain $`d`$ and $`e`$ after these transformations. Since the right-hand side of (7), which by Steps 1 and 2 divides the determinant as a polynomial in $`a,b,c,d,e`$, also has degree $`n`$ in $`a`$, in $`b`$, and in $`c`$, and degree $`1`$ in $`d`$, and in $`e`$, the determinant and the right-hand side of (7) differ only by a multiplicative constant.
Step 4. The evaluation of the multiplicative constant. By the preceding steps we know that the determinant equals (8). In particular, if we set $`a=c=d=e=0`$ and $`b=1`$, we have
$$\underset{1i,jn}{det}\left(n(i,j)\right)=C(n)(n+1)/2.$$
$`(9)`$
The matrix on the left-hand side of (9) is a circulant matrix with entries $`1,2,\mathrm{},n`$. Hence, its determinant equals
$$\underset{\omega \text{ : zero of }x^n1}{}(1+2\omega +3\omega ^2+\mathrm{}+n\omega ^{n1}).$$
The sum is easily evaluated by observing that it is the derivative of a geometric series. It turns out to be equal to $`n/(1\omega )`$. The resulting product simplifies by the observation
$$\underset{\omega \text{ : zero of }x^n1,\omega 1}{}(1\omega )=(1+x+\mathrm{}+x^{n1})|_{x=1}=n.$$
$`(10)`$
Thus, the determinant in (9) equals $`(n)^{n1}(n+1)/2`$. Therefore $`C(n)`$ is equal to $`(n)^{n1}`$.
This finishes the proof of (7) and thus of the theorem.
Theorem 8.— Let $`a`$ be an indeterminate. For any positive integer $`n`$ and integers $`i,j`$ let $`s(i,j)`$ denote (the representative between $`0`$ and $`n1`$ of) the residue class of $`ij+1`$ mod $`n`$. Then
$$\begin{array}{c}\text{ }\underset{1i,jn1}{det}\left(\left\{\begin{array}{cc}\left(\mathrm{n}\mathrm{m}1\right)+\mathrm{j}\left(1+\mathrm{a}\mathrm{n}\right)& \text{if }\mathrm{i}=\mathrm{j}2\left(\text{mod }\mathrm{n}\right)\\ \\ \left(\mathrm{n}1\right)\left(\mathrm{m}1\right)+\mathrm{j}\left(1\mathrm{a}\right)& \text{if }\mathrm{i}=\mathrm{j}3\left(\text{mod }\mathrm{n}\right)\\ \\ \left(\mathrm{n}\mathrm{m}32\mathrm{s}(\mathrm{i},\mathrm{j})\right)+\mathrm{j}& \text{otherwise}\end{array}\right\}\right)\text{ }\hfill \\ \text{ }=(1)^{n1}\frac{1}{n}\underset{i=2}{\overset{n}{}}(nmia).(11)\text{ }\hfill \end{array}$$
Proof.—The matrix underlying this determinant is a matrix whose elements have a uniform definition, except for two (broken) diagonals, the one with $`i=j2(\text{mod }n)`$, and the one with $`i=j3(\text{mod }n)`$.
To begin with, we reorder the rows so that the next-to-last row becomes the first row, the last row becomes the second row, and then follow the remaining rows in their original order.
Now we add all the rows to the (new) last row. In the resulting matrix, we change the sign of the last row and, subsequently, move it up so that it becomes the third row. As a result of these manipulations, we obtain
$$(1)^{n1}\underset{1i,jn1}{det}\left(\left\{\begin{array}{cc}(n1)(m1)+j(1a)& \text{if }i=j\\ (nm1)+j(an+1)& \text{if }i=j+1\\ mn+12i+3j& \text{if }i<j\\ m+n+12i+3j& \text{if }i>j+1\end{array}\right\}\right).$$
$`(12)`$
Next we apply further row operations. We subtract the second row from the first, the third from the second, …, the $`(n1)`$-st row from the $`(n2)`$-nd row, in that order. Subsequently, we repeat the same kind of operations, but stop before the last row, i.e., we subtract the second row from the first, the third from the second, …, the $`(n2)`$-nd row from the $`(n3)`$-rd row, in that order. As a result, the above determinant is converted into the determinant of the following matrix
$$\left(\left\{\begin{array}{cc}nmja& \text{if }i=j2\\ 2nm+2nj(3an)& \text{if }i=j1\text{ and }i<n2\\ nm2nj(3a2n)& \text{if }i=j\text{ and }i<n2\\ j(an)& \text{if }i=j+1\text{ and }i<n2\\ 2& \text{if }i=n2\text{ and }jn4\\ (n3)(an)+2& \text{if }i=n2\text{ and }j=n3\\ nm2n+2(n2)(2an)& \text{if }i=n2\text{ and }j=n2\\ nnm+1+(n1)(a1)& \text{if }i=n2\text{ and }j=n1\\ mn+3+3j& \text{if }i=n1\text{ and }jn3\\ nm1+(n2)(an+1)& \text{if }i=n1\text{ and }j=n2\\ (n1)(ma)& \text{if }i=n1\text{ and }j=n1\\ 0& \text{otherwise}\end{array}\right\}\right).$$
$`(13)`$
This is a matrix with four “special” diagonals (the diagonals with $`i=j2`$, $`j1`$, $`j`$, $`j+1`$, respectively) and two “special” rows (the last two rows). All other entries are zero.
Now we factor $`(nm3a)`$ out of the first row. Subsequently, we substract $`(an)`$ times the (new) first row from the second. Now one is able to factor $`(nm4a)`$ out of the (new) second row. Next, we substract $`2(an)`$ times the (new) second row from the third row. Etc. We stop this procedure in the $`(n3)`$-rd row. Thus, the determinant of the matrix in (13) is equal to $`_{i=3}^{n1}(nmia)`$ times the determinant of the matrix
$$\left(\begin{array}{ccccc}1& 2& 1& 0& \mathrm{}\\ 0& 1& 2& 1& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}\\ \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& 0\\ & & & 0& 1& 2& 1\\ a_{n2,1}& a_{n2,2}& \multicolumn{4}{c}{}& a_{n2,n1}\\ a_{n1,1}& a_{n1,2}& \multicolumn{4}{c}{}& a_{n1,n1}\end{array}\right),$$
$`(14)`$
where the entries in the last two rows are still the same as in (13).
We now perform the final set of transformations. We add column 1 through column $`n2`$ to column $`n1`$, and then we add
$$\underset{j=1}{\overset{n3}{}}(nj1)\left(\text{column }j\right)$$
to column $`n2`$. The effect is that a block matrix is obtained of the form
$$\left(\begin{array}{cc}U& 0\\ & M\end{array}\right),$$
where $`U`$ is an $`(n3)\times (n3)`$ upper-triangular matrix with 1s on the diagonal, and where $`M`$ is the $`2\times 2`$ matrix
$$\left(\begin{array}{cc}nmn2a+2& n2\\ \frac{1}{2}(n2)(n+m+2anm3)& \frac{5}{2}n\frac{n^2}{2}+ma3\end{array}\right).$$
Clearly, the determinant of $`U`$ is 1, while the determinant of $`M`$ is $`(ma)(nm2a)`$. Putting everything together, we have completed the proof of (11).
Corollary 9.— Let $`a`$ be an indeterminate. For any positive integer $`n`$ and integers $`i,j`$ let $`s(i,j)`$ denote (the representative between $`0`$ and $`n1`$ of) the residue class of $`ij+1`$ mod $`n`$. Then
$$\begin{array}{c}\text{ }\underset{1i,jn1}{det}\left(\left\{\begin{array}{cc}2s(i,j)a+j1& \text{if }ij2(\text{mod }n)\\ (n1)(n+aj1)& \text{if }i=j2(\text{mod }n)\end{array}\right\}\right)\text{ }\hfill \\ \text{ }=n^{n2}\underset{i=0}{\overset{n2}{}}(i+a).(15)\text{ }\hfill \end{array}$$
Proof.—In the determinant, we move the last row on top, replace the (now) last row by the sum of all the rows, factor $`(1)`$ out of the resulting row, and finally move it up so that it becomes the second row, retaining the order of all the other rows. These operations did not change the value of the determinant. However, the resulting determinant is exactly the determinant in (12) with $`a=n`$ and $`m=a+n`$. As we have shown in the proof of Theorem 8, the latter determinant differs from the determinant in (11) just by a sign of $`(1)^{n1}`$. Thus, we obtain the right-hand side of (15).
5. Closed form evaluations for Scott-type permanents
Theorem 10.— Let $`n`$, $`m`$ and $`r`$ be positive integers and $`d=\text{gcd}(n,r)`$. Then
$$\begin{array}{cc}& \text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{m}{}}a_{\mathrm{}}y^\mathrm{}n+\underset{\mathrm{}=0}{\overset{m}{}}b_{\mathrm{}}y^{\mathrm{}n+r})\hfill \\ & =\frac{\begin{array}{cc}& d^n{\displaystyle \underset{i=1}{\overset{d}{}}}(\left({\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}a_{\mathrm{}}\right)^{n/d}({\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}(in\mathrm{}1)a_{\mathrm{}}/{\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}da_{\mathrm{}})_{n/d}\hfill \\ & ({\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}b_{\mathrm{}})^{n/d}({\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}(irn\mathrm{}1)b_{\mathrm{}}/{\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}db_{\mathrm{}})_{n/d})\hfill \end{array}}{\left(\left({\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}a_{\mathrm{}}\right)^{n/d}\left({\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}b_{\mathrm{}}\right)^{n/d}\right)^d},\hfill \end{array}$$
where $`(\alpha )_k`$ is the standard notation for shifted factorials, $`(\alpha )_k:=\alpha (\alpha +1)\mathrm{}(\alpha +k1)`$, $`k1`$, and $`(\alpha )_0:=1`$.
Proof.—Let us first consider the case that $`nr`$. According to Theorem 3, we have to compute the quotient $`\text{Fes}(Q)/\text{Res}(x^n1,Q)`$, where $`Q=_{\mathrm{}=0}^ma_{\mathrm{}}y^\mathrm{}n+_{\mathrm{}=0}^mb_{\mathrm{}}y^{\mathrm{}n+r}`$. In order to compute $`\text{Fes}(Q)`$, in (3) we replace $`a_\mathrm{}n`$ by $`a_{\mathrm{}}`$ and $`a_{\mathrm{}n+r}`$ by $`b_{\mathrm{}}`$, $`\mathrm{}=0,1,\mathrm{},m`$, and set all other $`a_{\mathrm{}}`$’s equal to zero. In the resulting determinant we move the last row to the top, thus creating a sign of $`(1)^{n1}`$, and finally apply Proposition 6 with $`x_j=_{\mathrm{}=0}^m(n\mathrm{}j+1)a_{\mathrm{}}`$ and $`y_j=_{\mathrm{}=0}^m(n\mathrm{}+rj+1)b_{\mathrm{}}`$, $`j=1,2,\mathrm{},n`$. For the evaluation of $`\text{Res}(x^n1,Q)`$ we note that
$$\begin{array}{cc}& \text{Res}(x^n1,\underset{\mathrm{}=0}{\overset{m}{}}a_{\mathrm{}}y^\mathrm{}n+\underset{\mathrm{}=0}{\overset{m}{}}b_{\mathrm{}}y^{\mathrm{}n+r})\hfill \\ & =\underset{\omega \text{ : zero of }x^n1}{}\left(\underset{\mathrm{}=0}{\overset{m}{}}a_{\mathrm{}}+\omega ^r\underset{\mathrm{}=0}{\overset{m}{}}b_{\mathrm{}}\right)=\text{Res}(x^n1,y^r\underset{\mathrm{}=0}{\overset{m}{}}b_{\mathrm{}}+\underset{\mathrm{}=0}{\overset{m}{}}a_{\mathrm{}}).\hfill \end{array}$$
Now we can apply Lemma 5.
If on the other hand $`nr`$, then $`d=n`$. It can be verified directly that the claimed formula remains valid in that case, too.
Remark.— Given a polynomial $`Q(y)`$, there is no unique way to write it in the form $`_{\mathrm{}=0}^ma_{\mathrm{}}y^\mathrm{}n+_{\mathrm{}=0}^mb_{\mathrm{}}y^{\mathrm{}n+r}`$. For example, we may write $`Q(y)=y^n+a+b`$ as $`Q(y)=(y^n+ay^0)+b`$, or as $`Q(y)=(y^n+(a+b))`$ (i.e., either with $`a_1=1`$, $`a_0=a`$, $`r=0`$, $`b_0=b`$, or with $`a_1=1`$, $`a_0=a+b`$, $`b_{\mathrm{}}=0`$ for all $`\mathrm{}`$). Regardless which choice we make, Theorem 10 yields
$$\text{PER}(x^n1,y^n+a+b)=(1)^{n+1}\frac{\underset{i=1}{\overset{n}{}}(i(ni)(a+b))}{(a+b+1)^n}.$$
Corollary 11.— Let $`n`$ and $`m`$ be positive integers. Then
$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{m}{}}a_{\mathrm{}}y^\mathrm{}n)=\left(n\underset{\mathrm{}=0}{\overset{m}{}}\mathrm{}a_{\mathrm{}}/\underset{\mathrm{}=0}{\overset{m}{}}a_{\mathrm{}}\right)_n.$$
Corollary 12.— We have
$$\text{PER}(x^n1,y^{mn}+\mathrm{}+y^{2n}+y^n+1)=(mn/2)_n.$$
Corollary 13.— If $`m`$ is even then
$$\text{PER}(x^n+1,y^{mn}+\mathrm{}+y^{2n}+y^n+1)=(mn/2)_n.$$
Proof.—We use the case $`a_{\mathrm{}}=(1)^{\mathrm{}}`$ in Corollary 11, and the fact that
$$\text{per}\left(\frac{1}{x_i\sqrt[n]{1}y_j}\right)_{1in,\mathrm{\hspace{0.17em}1}jmn}=\text{per}\left(\frac{1}{x_iy_j/\sqrt[n]{1}}\right)_{1in,\mathrm{\hspace{0.17em}1}jmn}.$$
Corollary 14.— We have
$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{m}{}}\mathrm{}y^\mathrm{}n)=(n(2m+1)/3)_n.$$
Corollary 15.— We have
$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{m}{}}\mathrm{}y^{\mathrm{}^2n})=(nm(m+1)/2)_n.$$
Corollary 16.— We have
$$\text{PER}(x^n1,y^{mn}+ay^{rn}+b)=((m+ra)n/(a+b+1))_n.$$
Corollary 17.— We have
$$\text{PER}(x^n1,y^{mn}+1)=(mn/2)_n.$$
For $`m=1`$ one recovers Scott’s identity stated at the beginning of the introduction.
Corollary 18.— If $`m+ra=a+b+10`$, then
$$\text{PER}(x^n1,y^{mn}+ay^{rn}+b)=(1)^{n+1}n!.$$
Corollary 19.— If $`a2`$ then
$$\text{PER}(x^n1,y^{2n}+ay^n+1)=(1)^{n+1}n!.$$
Corollary 20.— If $`m+ra=0`$ and $`a+b+10`$, then
$$\text{PER}(x^n1,y^{mn}+ay^{rn}+b)=0.$$
Corollary 21.— If $`b1`$ then
$$\text{PER}(x^n1,y^{2n}2y^n+b)=0.$$
Corollary 22.— Let $`n`$ and $`m`$ be positive integers and $`d=\text{gcd}(n,m)`$. Then
$$\text{PER}(x^n1,y^m+b)=\frac{d^n{\displaystyle \underset{i=1}{\overset{d}{}}}\left(\left(\frac{im1}{d}\right)_{n/d}(b)^{n/d}\left(\frac{i1}{d}\right)_{n/d}\right)}{\left(1(b)^{n/d}\right)^d}.$$
Proof.—In Theorem 10, set $`m=0`$, $`a_0=b`$, $`b_0=1`$, and replace $`r`$ by $`m`$, in this order.
Corollary 23.— If $`\text{gcd}(m,n)=1`$, then
$$\text{PER}(x^n1,y^m+b)=(1)^{n+1}\frac{m(m1)\mathrm{}(mn+1)}{1(b)^n}.$$
Corollary 24.— If $`\text{gcd}(m,n)=1`$, then
$$\begin{array}{cc}& \text{PER}(x^{s(n1)}+\mathrm{}+x^{2s}+x^s+1,y^{s(m1)}+\mathrm{}+y^{2s}+y^s+1)\hfill \\ & =\frac{{\displaystyle \underset{i=0}{\overset{s1}{}}}\left({\displaystyle \underset{\mathrm{}=0}{\overset{n1}{}}}(i+\mathrm{}s){\displaystyle \underset{\mathrm{}=0}{\overset{n1}{}}}(i+\mathrm{}sms)\right)}{(mns)^s}.\hfill \end{array}$$
Proof.—Consider the Scott-type permanent $`\text{PER}(x^{sn}1,y^{sm}q^{sm})`$ (expressed in terms of its definition). When we multiply it by $`(1q)^s`$ and then perform the limit $`q1`$, then the permanent reduces to $`\text{PER}(x^{s(n1)}+\mathrm{}+x^s+1,y^{s(m1)}+\mathrm{}+y^s+1)`$, as is straightforward to see. On the other hand, the permanent that we started with is the permanent in Corollary 22 with $`n`$ replaced by $`sn`$, $`m`$ replaced by $`sm`$, and $`b=q^{sm}`$. Indeed, if we multiply the right-hand side from Corollary 22 (with these choices for the parameters) by $`(1q)^{sm}`$, and then perform the limit $`q1`$, we obtain exactly the claimed result.
Corollary 25.— If $`\text{gcd}(m,n)=1`$, then
$$\begin{array}{cc}& \text{PER}(x^{n1}+\mathrm{}+x+1,y^{m1}+\mathrm{}+y+1)\hfill \\ & =(1)^{n+1}\frac{(m1)\mathrm{}(mn+1)}{n}.\hfill \end{array}$$
Corollary 26.— If $`\text{gcd}(m,n)=1`$ and $`n`$ is odd, then
$$\text{PER}(x^n1,y^m+1)=\frac{m(m1)\mathrm{}(mn+1)}{2}.$$
Corollary 27.— If $`n`$ is odd, then
$$\text{PER}(x^n1,y^{n+1}+1)=\frac{(n+1)!}{2}.$$
Corollary 28.— Let $`n`$ and $`r`$ be positive integers (not necessarily $`n>r`$) and $`d=\text{gcd}(n,r)`$. Then
$$\begin{array}{cc}& \text{PER}(x^n1,y^n+ay^r+b)\hfill \\ & =\frac{d^n{\displaystyle \underset{i=1}{\overset{d}{}}}\left((b+1)^{\frac{n}{d}}\left({\displaystyle \frac{ibb+in1}{d(b+1)}}\right)_{\frac{n}{d}}(a)^{\frac{n}{d}}\left({\displaystyle \frac{ir1}{d}}\right)_{\frac{n}{d}}\right)}{\left((b+1)^{\frac{n}{d}}(a)^{\frac{n}{d}}\right)^d}.\hfill \end{array}$$
Proof.—In Theorem 10, set $`m=1`$, $`a_0=b`$, $`a_1=1`$, $`b_0=a`$, and $`b_1=0`$.
Corollary 29.— Let $`n`$ and $`r`$ be positive integers (not necessarily $`n>r`$) and $`\text{gcd}(n,r)=1`$. Then
$$\text{PER}(x^n1,y^n+ay^r+b)=(1)^{n+1}\frac{\underset{i=1}{\overset{n}{}}(i(ni)b))a^n(r)_n}{(b+1)^n(a)^n}.$$
If $`1rn1`$ then $`(r)_n=0`$. Thus, one recovers the results in . On the other hand, if we set $`r=n+1`$, we obtain, for example, the following result.
Corollary 30.— We have
$$\text{PER}(x^n1,y^{n+1}+y^n1)=n^n(1)^n(n+1)!.$$
Corollary 31.— We have
$$\text{PER}(x^n1,y^n+ny1)=1.$$
Theorem 32.— Let $`n`$ and $`m`$ be positive integers and $`a`$ be an arbitrary number. Then
$$\begin{array}{cc}& \text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{mn1}{}}(\mathrm{}+a)y^{\mathrm{}})\hfill \\ & =(1)^{n1}\frac{n(m1)V_n(a,m)}{6\left(mn+2a1\right)}\left(a+(m1)n+1\right)_{n2},\hfill \end{array}$$
where $`V_n(a,m)`$ is the polynomial
$$V_n(a,m)=16a+6a^2+n2an5mn+10amnmn^2+4m^2n^2.$$
Proof.—Let $`Q=_{\mathrm{}=0}^{mn1}(\mathrm{}+a)y^{\mathrm{}}`$. Using Theorem 3 again, we have to compute $`\text{Fes}(Q)/\text{Res}(x^n1,Q)`$. In order to compute $`\text{Fes}(Q)`$, in (3) set $`a_{\mathrm{}}=\mathrm{}+a`$, $`\mathrm{}=0,1,\mathrm{},mn1`$. In the resulting determinant we move the last row to the top, thus creating a sign of $`(1)^{n1}`$, and finally apply Theorem 7 with $`a=1`$, $`b=(m1)n1`$, $`c=a1`$, $`d=n^2(m1)(2m1)/6an(m1)/2`$ and $`e=a+n(m1)/21`$.
For the computation of $`\text{Res}(x^n1,Q)`$ we note that
$$\begin{array}{cc}& \text{Res}(x^n1,Q)=\underset{\omega \text{ : zero of }x^n1}{}\left(\underset{\mathrm{}=0}{\overset{mn1}{}}(\mathrm{}+a)\omega ^{\mathrm{}}\right)\hfill \\ & =\left(\left(\genfrac{}{}{0pt}{}{mn}{2}\right)+mna\right)\underset{\omega \text{ : zero of }x^n1,\omega 1}{}\left(\underset{\mathrm{}=1}{\overset{mn}{}}\mathrm{}\omega ^\mathrm{}1\right).\hfill \end{array}$$
The sum in the last line is the derivative of a geometric series, and is therefore easily evaluated. The result of the summation turns out to be $`mn/(1\omega )`$. The computation is completed by the observation (10), and some simplification.
Corollary 33.— Let $`n`$ and $`m`$ be positive integers, $`n2`$. Then
$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{mn1}{}}\mathrm{}y^{\mathrm{}})=(1)^{n1}\frac{(4mnn1)(mn2)!}{6(mnn1)!}.$$
Corollary 34.— Let $`n`$ and $`m`$ be positive integers, $`n2`$. Then
$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{mn1}{}}(\mathrm{}+1)y^{\mathrm{}})=(1)^{n1}\frac{(4mnn+1)(mnn)(mn1)!}{6(mnn+1)!}.$$
Corollary 35.— Let $`n`$ and $`m`$ be positive integers, $`n2`$. Then
$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{mn1}{}}(mn\mathrm{})y^{\mathrm{}})=\frac{(m1)(n+1)!}{6}.$$
Corollary 36.— Let $`n`$ and $`m`$ be positive integers, $`n2`$. Then
$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{mn1}{}}(mn\mathrm{}1)y^{\mathrm{}})=\frac{(m1)n!}{6}.$$
Remark.— It is also possible to move forward and derive formulas for $`\text{PER}(x^n1,_{\mathrm{}=0}^{mn/s1}(\mathrm{}+a)y^\mathrm{}s)`$, where $`s`$ is some positive integer. This would require to find analogues of Theorem 7 in which $`n(i,j)`$ is replaced by $`Dn(i,j)`$, where $`D`$ is the inverse of $`s/\text{gcd}(n,s)`$ modulo $`n/\text{gcd}(n,s)`$. As calculations aided by the computer indicate, the resulting determinant evaluations have forms very similar to (7). That is, the result shows a product of linear factors in $`a`$ and $`b`$ as the one on the right-hand side of (7), and one irreducible polynomial of higher degree (such as $`U_n(a,b,c,d,e)`$ in (7)). However, as $`D`$ increases, the degree of the irreducible polynomial also increases, whereas the amount of linear factors in $`a`$ and $`b`$ decreases, so that the results become more and more unwieldy. We therefore content ourselves with stating the result when $`s`$ divides $`n`$.
Theorem 37.— Let $`n`$, $`m`$ and $`s`$ be positive integers so that $`sn`$, and let $`a`$ be an arbitrary number. Then
$$\begin{array}{cc}& \text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{\frac{mn}{s}1}{}}(\mathrm{}+a)y^\mathrm{}s)=(1)^{n1}\frac{s^{n2s}}{6^s(mn+2ass)^s}\hfill \\ & \times \underset{k=0}{\overset{s1}{}}\left((a+\frac{1}{s}(nmnk)+1)_{n/s2}V_{n,s}(a,m,k)\right),\hfill \end{array}$$
where $`V_{n,s}(a,m,k)`$ is the polynomial
$$\begin{array}{cc}\hfill \mathrm{V}_{\mathrm{n},\mathrm{s}}& (\mathrm{a},\mathrm{m},\mathrm{k})=6\mathrm{k}^2\mathrm{mn}+6\mathrm{k}\mathrm{m}\mathrm{n}^210\mathrm{k}\mathrm{m}^2\mathrm{n}^2+\mathrm{mn}^35\mathrm{m}^2\mathrm{n}^3+4\mathrm{m}^3\mathrm{n}^36\mathrm{k}^2\mathrm{s}\hfill \\ & +12\mathrm{a}\mathrm{k}^2\mathrm{s}6\mathrm{k}\mathrm{n}\mathrm{s}+12\mathrm{a}\mathrm{k}\mathrm{n}\mathrm{s}+12\mathrm{k}\mathrm{m}\mathrm{n}\mathrm{s}24\mathrm{a}\mathrm{k}\mathrm{m}\mathrm{n}\mathrm{s}\mathrm{n}^2\mathrm{s}+2\mathrm{a}\mathrm{n}^2\mathrm{s}+6\mathrm{m}\mathrm{n}^2\mathrm{s}\hfill \\ & 12\mathrm{a}\mathrm{m}\mathrm{n}^2\mathrm{s}5\mathrm{m}^2\mathrm{n}^2\mathrm{s}+10\mathrm{a}\mathrm{m}^2\mathrm{n}^2\mathrm{s}2\mathrm{k}\mathrm{s}^2+12\mathrm{a}\mathrm{k}\mathrm{s}^212\mathrm{a}^2\mathrm{ks}^2\mathrm{ns}^2\hfill \\ & +6\mathrm{a}\mathrm{n}\mathrm{s}^26\mathrm{a}^2\mathrm{ns}^2+\mathrm{mns}^26\mathrm{a}\mathrm{m}\mathrm{n}\mathrm{s}^2+6\mathrm{a}^2\mathrm{mns}^2.\hfill \end{array}$$
Proof.—Let $`Q=_{\mathrm{}=0}^{mn1}(\mathrm{}+a)y^\mathrm{}s`$. Again, according to Theorem 3, we have to compute $`\text{Fes}(Q)/\text{Res}(x^n1,Q)`$. In order to compute $`\text{Fes}(Q)`$, in (3) set $`a_\mathrm{}s=\mathrm{}+a`$, $`\mathrm{}=0,1,\mathrm{},mn/s1`$, and all other $`a_i`$’s to zero. In the resulting determinant we move the last row to the top, thus creating a sign of $`(1)^{n1}`$. Since only every $`s`$-th $`a_i`$ is nonzero, we are dealing with a determinant of a matrix in which a lot of entries are zero. If we permute rows and columns so that first come the rows and columns whose indices are congruent $`1`$ mod $`s`$, then come the rows and columns whose indices are congruent $`2`$ mod $`s`$, etc., then the matrix of which we want to compute the determinant assumes a block form, with $`(n/s)\times (n/s)`$ blocks on the diagonal, and zeroes otherwise. Therefore the determinant equals the product of the determinants of the $`s`$ matrices of dimension $`(n/s)\times (n/s)`$ on the diagonal. Each of these determinants can be evaluated by means of Theorem 7. We leave it to the reader to fill in the details.
For the computation of $`\text{Res}(x^n1,Q)`$ we proceed as in the proof of Theorem 32.
Theorem 38.— Let $`n`$ and $`m`$ be positive integers, $`n2`$, and $`a`$ be an arbitrary number. Then
$$\begin{array}{cc}& \text{PER}(x^{n1}+\mathrm{}+x+1,\underset{\mathrm{}=0}{\overset{mn1}{}}(\mathrm{}+a)y^{\mathrm{}})\hfill \\ & =(1)^{n1}\left(a+(m1)n+1\right)_{n1}.\hfill \end{array}$$
Proof.—Let $`P(x)=x^{n1}+\mathrm{}+x+1`$ and $`Q(y)=_{\mathrm{}=0}^{mn1}(\mathrm{}+a)y^{\mathrm{}}`$. This time we apply Theorem 4. According to that theorem, we have to compute $`\stackrel{𝑒}{\mathrm{Fes}}(Q)/\text{Res}(P,Q)`$. In order to compute $`\stackrel{𝑒}{\mathrm{Fes}}(Q)`$, in (4) we set $`a_{\mathrm{}}=\mathrm{}+a`$, $`\mathrm{}=0,1,\mathrm{},mn1`$. The resulting determinant is exactly $`m^{n1}`$ times the determinant in Corollary 9 with $`a`$ replaced by $`((m1)n+a+1)`$. For the computation of the resultant of $`P`$ and $`Q`$ we proceed as in the proof of Theorem 7. Simplification of the result yields the claimed expression.
Theorem 39.— Let $`n`$ and $`m`$ be positive integers, $`n2`$, and $`a`$ be an arbitrary number. Then
$$\begin{array}{cc}& \text{PER}(x^{n1}+\mathrm{}+x+1,\underset{\mathrm{}=0}{\overset{mn2}{}}(\mathrm{}+a)y^{\mathrm{}})\hfill \\ & =(1)^{n1}\frac{(nmn)_{n1}(nm+a1)^{n1}}{(mn+a1)^n(a1)^n}.\hfill \end{array}$$
Proof.—Let $`P(x)=x^{n1}+\mathrm{}+x+1`$ and $`Q(y)=_{\mathrm{}=0}^{mn1}(\mathrm{}+a)y^{\mathrm{}}`$. Using Theorem 4 again, we have to compute $`\stackrel{𝑒}{\mathrm{Fes}}(Q)/\text{Res}(P,Q)`$. In order to compute $`\stackrel{𝑒}{\mathrm{Fes}}(Q)`$, in (4) we set $`a_{\mathrm{}}=\mathrm{}+a`$, $`\mathrm{}=0,1,\mathrm{},mn2`$. The resulting determinant is exactly $`m^{n1}`$ times the determinant in Theorem 8 with $`m`$ replaced by $`mn+A1`$, $`a`$ replaced by $`(mn+A1)/m`$, and $`A`$ replaced by $`a`$, in that order.
For the computation of $`\text{Res}(x^n+\mathrm{}+x+1,Q)`$ we proceed similarly as in the proof of Theorem 32. Using an observation from that proof, we note that
$$\begin{array}{cc}& \text{Res}(x^n+\mathrm{}+x+1,Q)=\underset{\omega \text{ : zero of }x^n1,x1}{}\left(\underset{\mathrm{}=0}{\overset{mn2}{}}(\mathrm{}+a)\omega ^{\mathrm{}}\right)\hfill \\ & =\underset{\omega \text{ : zero of }x^n1,x1}{}\left(\frac{mn}{1\omega }(mn+a1)\omega ^{mn1}\right)\hfill \\ & =(1)^{n1}\underset{\omega \text{ : zero of }x^n1,x1}{}\frac{mn+a1+\omega (1a)}{\omega (1\omega )}\hfill \\ & =\frac{(mn+a1)^n(a1)^n}{n^2m}.\hfill \end{array}$$
The result follows now upon some simplification.
6. Sums of involutions
As in \[6, Sec. 4\], we may obtain interesting summation theorems by combining special evaluations of Scott-type determinants (see Section 5) with Theorem 2. For example, if we combine Corollary 19 with Theorem 2, we obtain the following result.
Proposition 40.—Let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of $`x^n1`$. Then
$$\underset{\sigma (n)}{}\underset{(ij)\sigma }{}\frac{1}{(x_ix_j)^2}\underset{(k)\sigma }{}\frac{n+1}{2x_k}=(1)^{n+1}n!,$$
where the first product is over all transpositions $`(ij)`$ of $`\sigma `$ in its disjoint cycle decomposition, and where the second product is over all fixed points $`k`$ of $`\sigma `$.
Similarly, if we combine Corollary 21 with Theorem 2, we obtain the result below.
Proposition 41.—Let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of $`x^n1`$. Then
$$\underset{\sigma (n)}{}\underset{(ij)\sigma }{}\frac{1}{(x_ix_j)^2}\underset{(k)\sigma }{}\frac{n1}{2x_k}=0,$$
where the first product is over all transpositions $`(ij)`$ of $`\sigma `$ in its disjoint cycle decomposition, and where the second product is over all fixed points $`k`$ of $`\sigma `$.
If we combine Corollary 31 with Theorem 2, then we obtain the following result.
Proposition 42.—Let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of $`x^n1`$. Then
$$\underset{\sigma (n)}{}\underset{(ij)\sigma }{}\frac{1}{(x_ix_j)^2}\underset{(k)\sigma }{}\frac{2n+(n+1)x_k}{2x_k^2}=1,$$
where the first product is over all transpositions $`(ij)`$ of $`\sigma `$ in its disjoint cycle decomposition, and where the second product is over all fixed points $`k`$ of $`\sigma `$.
Finally, if we combine Corollary 27 with Theorem 2, we obtain the result below.
Proposition 43.—Let $`n`$ be odd, and let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of $`x^n1`$. Then
$$\underset{\sigma (n)}{}\underset{(ij)\sigma }{}\frac{1}{(x_ix_j)^2}\underset{(k)\sigma }{}\frac{1n+(3+n)x_k}{2(1+x_k)x_k}=\frac{(n+1)!}{2},$$
where the first product is over all transpositions $`(ij)`$ of $`\sigma `$ in its disjoint cycle decomposition, and where the second product is over all fixed points $`k`$ of $`\sigma `$.
References
G. E. Andrews, I. P. Goulden, D. M. Jackson.—Generalizations of Cauchy’s summation theorem for Schur functions, Trans. Amer. Math. Soc. 310 (1988), pp. 805–820.
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G.-N. Han.—Interpolation entre Cauchy et Borchardt, in preparation, 1999.
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R. Kittappa.—Proof of a conjecture of 1881 on permanents, Linear and Multilinear Algebra 10 (1981), pp. 75–82.
C. Krattenthaler.—An alternative evaluation of the Andrews–Burge determinant, in: Mathematical Essays in Honor of Gian-Carlo Rota, B. E. Sagan, R. P. Stanley, eds., Progress in Math., vol. 161, Birkhäuser, Boston, 1999, pp. 263–270.
C. Krattenthaler.—Advanced determinant calculus, Séminaire Lotharingien Combin. 42, “The Andrews Festschrift” (1999), paper B42q, 67 pp.
A. Lascoux.—Square-ice Enumeration, Séminaire Lotharingien Combin. 42, “The Andrews Festschrift” (1999), paper B42p, 15 pp.
I. G. Macdonald.—Symmetric functions and Hall polynomials, second edition.—Clarendon Press, Oxford, 1995.
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H. Minc.—On a conjecture of R. F. Scott, Linear Algebra Appl. 28 (1979), pp. 141–153.
T. Muir.—The theory of determinants in the historical order of development, 4 vols.—Macmillan, London, 19061923.
R. F. Scott.—Mathematical notes, Messenger of Math. 10 (1881), pp. 142–149.
D. Svrtan.—Proof of Scott’s conjecture, Proc. Amer. Math. Soc. 87 (1983), pp. 203–207.
I.R.M.A. and C.N.R.S. Université Louis Pasteur 7, rue René-Descartes F-67084 Strasbourg, France guoniu@math.u-strasbg.fr Institut für Mathematik Universität Wien Strudlhofgasse 4 A-1090 Vienna, Austria kratt@ap.univie.ac.at
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# Kinetics of Ordering in Fluctuation-Driven First-Order Transitions: Simulations and Dynamical Renormalization
## I Introduction
Kinetics of growth can be broadly classified into two categories; nucleation and spinodal decomposition (or continuous ordering). The former applies to situations where the initial state is metastable and the latter to where the initial state is unstable. The distinction becomes unclear near the limit of metastability and in systems where the concept of metastability itself is ill-defined. A class of systems where the definition of metastability becomes ambiguous is one where there is a fluctuation-driven first-order phase transition. A notable example of a system undergoing such a transition is a symmetric di-block copolymer. These systems undergo microphase separation and there is a temperature-driven transition from a uniform phase to a lamellar phase. Another system which is governed by the same dynamical equation as block copolymers, and so exhibits a similar transition, is a Rayleigh-Benard Cell .
A theoretical model describing these transitions was proposed by Brazovskii in 1975. Di-block copolymers were shown to belong to the Brazovskii class by Leibler and the theoretical predictions regarding the equilibrium nature of the ordering transition have been experimentally verified. It was shown by Brazovskii, within a self-consistent Hartree approximation, that the fluctuations destroy the mean-field instability and lead to a first-order phase transition. Theories of nucleation and growth have been constructed based on the idea that the static, Brazovskii-renormalized theory can provide an effective potential for a stochastic Langevin equation. In this paper we examine the validity of this description by using numerical simulations to study the relaxational dynamics of a model described by the Brazovskii Hamiltonian and comparing the results to the predictions of the “static” nucleation theories and to the predictions of a dynamically renormalized theory.
## II The Equation of Motion and Numerical Simulations
Our starting point is the same as in Swift and Hohenberg and describes a system with model A dynamics and a Brazovskii Hamiltonian such as would apply to the diblock copolymers. The Brazovskii Hamiltonian is characterized by a fluctuation spectrum whose maximum occurs at a non-zero wave vector $`|𝐪|=q_0`$ and can be represented by a hypersphere in $`d`$ dimensions. The full form of the Hamiltonian is
$$H=𝑑𝐱[\frac{1}{2}\varphi (^2+q_0^2)^2\varphi +\frac{1}{2}\tau \varphi ^2+\frac{1}{4!}\lambda \varphi ^4]$$
(1)
The dynamics is taken to be relaxational and so the equation of motion is given by the Langevin Equation.
$$\frac{d\varphi }{dt}=M\frac{\delta H}{\delta \varphi }+\eta $$
(2)
where $`M`$ is the mobility (which sets the time scale for the problem) and $`\eta `$ is Gaussian noise ($`<\eta >=0,`$ $`<\eta (\stackrel{}{x},t)`$ $`\eta (\stackrel{}{x}^{},t^{})>=\delta (\stackrel{}{x}\stackrel{}{x}^{})\delta (tt^{})`$ ).
The stochastic Langevin equation derived from the Hamiltonian differs from the usual Ginzburg-Landau description because of the appearance of the unusual gradient term. This dynamical equation is usually referred to as the Swift-Hohenberg equation and falls into the Type $`I_s`$ classification of Cross and Hohenberg . In this classification the system is unstable to a static, spatially periodic structure.
The complete equation of motion is;
$$\frac{d\varphi }{dt}=M(q_0^2^2\varphi +^4\varphi +(q_0^4+\tau )\varphi +\frac{\lambda }{6}\varphi ^3)+\eta $$
(3)
In Eq. (3), the coupling constant, $`\lambda `$ has been rescaled by the noise strength. For systems where the noise strength is small, such as Rayleigh-Benard convection, the effective coupling constant is also small. In di-block copolymers, the noise strength is of the order of $`k_BT`$.
For numerical calculations, the Langevin equation must be approximated by a discrete equation.
$$\varphi (t+\mathrm{\Delta }t)=\varphi (t)M\mathrm{\Delta }t(\mathrm{\pounds }\varphi +\mathrm{\pounds }^2\varphi +(r+q_0^4)\varphi +\frac{\lambda }{6}\varphi ^3)+(M\mathrm{\Delta }t/(\mathrm{\Delta }x)^3)^{1/2}\eta $$
(4)
£ here is the discrete Laplacian. This can take on several forms though for this work we use the simplest form: $`\frac{d^2\varphi }{dx^2}\mathrm{\pounds }\varphi =\frac{1}{\mathrm{\Delta }x^2}(x+\mathrm{\Delta }x)+\varphi (x\mathrm{\Delta }x)2\varphi (x))`$ where the sum is over the dimensions of the lattice. Other choices which include next nearest or more complicated neighborhoods are possible. This is important if isotropy is of concern however the effect is small and can be ignored. The scaling of the noise in Eq. 4 reflects the fact that as the cell size or time steps become smaller (larger), the possible fluctuations become larger (smaller).
For these simulations we set $`q_0=1`$ and choose the lattice spacing such that one lamellar spacing spans six lattice sites ($`\mathrm{\Delta }x=6/2\pi `$) For most of this work the overall lattice size is $`60^3`$. The mobility, $`M`$, is set to one. In order for the simulations to be numerically stable the time scale must be smaller then some stability time, $`\mathrm{\Delta }t1/\mathrm{\Delta }x^4.`$ To satisfy this the timescale is chosen to be $`\mathrm{\Delta }t`$ = 0.01 and measurements are taken at intervals of $`100\mathrm{\Delta }t`$ or longer.
Previous theories of nucleation have analyzed the above equation under the approximation that the effects of fluctuations can be incorporated into a static renormalized free energy $`F_R`$ which replaces the bare free energy $`F`$ in Eq. (2) . We will compare our numerical results to the Brazovskii predictions for the static renormalized parameters and show that the static scenario is beautifully borne out by the simulations in three dimensions. We will then analyze the time-dependence of the structure factor as observed in the numerical simulations and show that this time dependence is inconsistent with the dynamical picture based on a static renormalized free energy. A theory based on a dynamical renormalization of Eq. (2) can provide a qualitatively correct description of the numerical simulations.
## III Statics
### A Brazovskii Theory
Brazovski’s treatment of the model within the Hartree approximation can be restated in terms of an expansion of the thermodynamic potential $`\mathrm{\Gamma }(\overline{\varphi })`$, the generating functional for the vertex functions. This approach has been described in detail by Fredrickson and Binder, and within the Hartree approximation leads to a renormalized mass term ($`r`$) in Eq.(1). The diagrams up to one loop are shown in Fig. 1, and the mass renormalization relation is given by;
$$r+(q^2q_0^2)^2=\tau +(q^2q_0^2)^2+\frac{\lambda }{2}\frac{d𝐪}{(2\pi )^d}(\tau +(q^2q_0^2)^2)^1$$
(5)
This approximation is made self consistent by replacing the bare parameter in the integrand by the renormalized parameter. Essentially the bare propagator on the loop in Fig. (1) is replaced by renormalized propagator. The leads to the Hartree result.
$$r=\tau +\frac{\lambda }{2}\frac{d𝐪}{(2\pi )^d}(r+(q^2q_0^2)^2)^1=\tau +\alpha \lambda /\sqrt{r}$$
(6)
where $`\alpha `$ is proportional to the surface area of a $`d`$ dimensional sphere of unit radius. According to Brazovski this approximation is good only for $`\lambda ^6<<1.`$
The interesting point about Eq. (6) is that even for negative values of the bare parameter, $`\tau `$, the renormalized parameter, $`r`$, is positive. This implies that the disordered phase is always metastable. Brazovskii went on to show that the bare coupling parameter $`\lambda `$ gets renormalized to a negative value leading to a $`\varphi ^6`$ theory and the possibility of a first-order phase transition.
### B Computational Results.
The numerical simulations can measure the static structure factor in the disordered phase which is predicted by the Brazovskii theory to be,
$$S(q)^1=r+(q^2q_0^2)^2$$
(7)
where $`r`$ is the renormalized control parameter . Thus, $`S(q_0)^1`$ is just $`r`$. The renormalized mass can, therefore, be measured by monitoring the peak of the structure factor. The Hartree calculations predict $`r`$ to be positive and asymptotically approaching zero as $`\tau \mathrm{}`$. Experiments in symmetric di-block copolymers have verified that the behavior of $`S(q_0)`$ is consistent with the Brazovskii predictions and is very different from the mean-field prediction ( $`S(q_0)=\tau ^1`$).
Fig. 2 shows results for $`r`$ obtained from our simulations for different values of the coupling constant, $`\lambda `$. The Gaussian case ($`\lambda =0`$) is also shown for comparison. The system was run for 1000 time steps (each time step being $`100\mathrm{\Delta }t`$) and 100 samples were taken. The data points plotted are the average of the samples while the error bars represent the standard deviation. For large positive values of the bare control parameter, $`\tau `$, the fourth order term becomes less important and $`r`$ approaches the Gaussian value. This is seen clearly for small values of $`\lambda `$. As $`\tau `$ approaches zero, the measured values of $`r`$ deviate from the Gaussian case and stay positive for all $`\tau `$. Some care was taken to normalize the values presented in Fig. 2. For the Gaussian case, $`S(q_0)^1`$ should be linearly related to $`\tau `$ and the slope of the line should be one. In calculating $`S(q_0)`$ several normalizations are needed including the normalization due to the Fourier Transform (FT) and the normalization due to circular averaging. While the FT normalization is just related to the system size, the normalization due to circular averaging is more complicated to calculate. Instead of calculating these normalizations directly, the slope of the raw result from the Gaussian case was used to normalize all of the data. Another concern is that for an infinite continuous system this line should have an x-intercept at zero; however, for the finite systems on a grid used in the simulations the x-intercept is slightly negative. To account for this, all of the data is shifted over by an amount equal to that intercept. The importance of this shift will become apparent in the next paragraph when a scaling is applied and the values of the results around zero are magnified. A small negative (positive) value becomes a much larger negative (positive) value and so shifts from negative to positive will become important.
In reference the authors show that, within the Hartree approximation, the $`r(\tau )`$curves for different $`\lambda ^{}s`$ could be described by a single functional form in terms of scaled variables $`\tau ^{}`$ and $`r^{}`$;
$`r^{}`$ $`=`$ $`r(\alpha \lambda )^{2/3}`$ (8)
$`\tau ^{}`$ $`=`$ $`\tau (\alpha \lambda )^{2/3}`$ (9)
Figure 2 shows plots of $`r^{}`$ against $`\tau ^{}`$obtained from simulations for the three different values of $`\lambda `$. The curves are seen to scale quite well but the scaled curve falls significantly below the theoretical prediction shown as the thick black line. The theoretical prediction that the value of $`r^{}`$(and $`r`$) never becomes negative and is asymptotic to zero is, however, clearly borne out by the simulations. The deviation from the theoretical line could be a system size effect, however, results for larger and smaller systems are consistent with the data shown in Fig.3. For the Gaussian case $`r^{}=\tau ^{}`$ and for large positive values of the bare parameter both the theory and simulations approach this. Since the theory is larger then the simulations and both are above the Gaussian, our simulations suggest that the contributions to the two point function from diagrams not included in the Hartree approximation serve to lower the overall correction to mean field theory. Another property that the simulation data exhibits in Fig. 3 which is not predicted by the theory is that the curves for different values of $`\lambda `$ diverge from each other at negative values of $`\tau `$, i.e. the scaling is not perfect. Smaller values of $`\lambda `$ lie closer to the theory as expected. The computations were done using values of $`\lambda 10^2`$ and larger while the theory is valid only for $`\lambda 10^6`$ so the theory’s scaling predictions seem to be far more robust then expected.
To further explore the nature of the phase transition in this model, we can compare scaled values of $`S(q_0)`$ obtained from a hot (random) and cold (purely modulated order) initial configurations. Fig. 4 shows just such a plot. As in the previous plot, the theory is shown as a thick black line. For each value of $`\lambda `$ there are two lines presented: one for a disordered start and one for an ordered start. The averaging was done as described above. For positive and small negative values of $`\tau `$ the results from the disordered start and the ordered start are nearly the same. For $`\tau ^{}`$ below some $`\tau _s^{}`$ the two differ with the ordered phase having a higher peak. It should be pointed out that for these values of $`\tau ^{}`$ the disordered state is not in equilibrium (metastable or stable)and there is a slow but definite time evolution of the structure factor over the time period in which the averages are taken. This time-evolution will be discussed in more detail in the section on dynamics. For now the disordered start points are included in the plot simply as a comparison to see where the ordered phase melts. The value where the two data sets diverge can be considered an estimate for the limit of stability (spinodal) of the lamellar phase. As a consistency check the average value of the wave vector was measured. For values of $`\tau `$ above $`\tau _s`$ the wave vector is small and points in a random direction, while for lower values of $`\tau `$ the average wave vector is large and and points along the direction that the system was prepared in. For all three sets of data the lamellar spinodal lies in the range of $`2.1<\tau _s^{}<1.9`$. This is consistent with the prediction from Hartree theory that the first-order transition to the lamellar phase occurs at a lower value of $`\tau ^{}`$, $`\tau _t^{}=2.74`$. That is to say that our measured value of the lamellar spinodal is above the predicted transition temperature as it should be. For values of $`\tau ^{}`$ below $`\tau _s^{}`$, however, systems prepared in the disordered phase always evolve towards the ordered phase. It is possible that the lamellar ordered phase is enhanced by the boundary conditions - the lattice size is always chosen to be commensurate with the lamellar wavelength - however, it should be noted that lamella form in many possible directions where the system size is not commensurate with the lamellar wavelength. Discussion in the dynamics sections should shed some light on this issue.
## IV Dynamics
### A Relaxational Dynamics Based on Static Renormalized Parameters
Previous analysis of nucleation and metastability, in the context of fluctuation-driven phase transitions, have been based on a Langevin equation with the force obtained from the static, Hartree-renormalized free energy function. Fredrickson and Binder used this approach to compute the nucleation barrier and the completion rate of nucleation and growth in di-block copolymers. The interfacial tension was found to be small leading to a small nucleation barrier and rapid nucleation for deep quenches. In this picture, there is no essential difference between the kinetics of a fluctuation-driven first-order transition and a weak-first order transition. The results of our simulations suggest a very different scenario for the growth of the lamellar structures in a Brazovskii model. The shape of the nucleating droplet was more carefully analyzed by Swift and Hohenberg by taking into account spatial inhomogeneities in the effective free energy function. Their analysis relied on constructing a coarse-grained free energy functional. The coarse graining was based on a momentum-shell renormalization idea where fluctuations with momenta far away from the shell defined by $`|𝐪|=𝐪_\mathrm{𝟎}`$ are successively integrated out. This analysis led to non-spherical droplets and were consistent with the picture of spinodal nucleation that had been obtained earlier.
As pointed out in the work of Swift and Hohenberg, a complete theory of the dynamics of fluctuation-driven first order transitions would have to be based on a coarse-graining of the full dynamics as expressed by the original Langevin equation with the bare Brazovskii Hamiltonian. In this paper, we compare the results of our numerical simulations to the predictions of a Hartree-renormalization of the full dynamical equation as expressed in Eq. (2). Our emphasis has been on understanding the nature of the dynamics of the metastable phase and we have not analyzed, in any detail, the spatial structures associated with the growth process.
### B Dynamical Renormalization and Simulation Results
In order to systematically analyze the effects of fluctuations on the kinetics of growth of the lamellar phase, perturbative techniques analogous to the static Hartree approximation have to be applied to the Langevin equation. This is most conveniently done through the dynamical-action formalism. The application of this method to the S-H equation is outlined by Ignatiev et. al. . In the dynamical-action formalism, the average value of an operator $`O(\varphi )`$ over the noise history is rewritten as a functional integral;
$$O\{\varphi (𝐱,t)\}=D\varphi \mathrm{exp}(S[\varphi ]),$$
(10)
where $`S[\varphi ]`$ is the dynamical action.
As in the static case, the calculation of dynamical correlation functions is most conveniently formulated through the construction of a generating function. To establish the closest analogy with the static calculations, it is useful to work with the generator of vertex functions, $`\mathrm{\Gamma }_{dyn}[\overline{\varphi }]`$, and establish a diagrammatic expansion which is the exact analog of the diagrams retained by Brazovskii in the static calculation. The simulation results demonstrate that the static properties of the S-H equation closely follow the Brazovskii predictions and, therefore, it seems appropriate to apply this approximation scheme to the dynamics. The correspondence between statics and dynamics becomes particularly transparent in a super-field formulation which shows that the dynamical perturbation theory in terms of the super fields has exactly the same structure as the static perturbation theory except for the appearance of a very different “kinetic” term which leads to a bare propagator that is distinct from the static theory. The super-field correlation functions can be calculated by constructing diagrams as in the static theory but replacing the static bare propagator by the appropriate dynamical one. The super-field correlation functions encode all the dynamical correlations of the field $`\varphi `$ and the latter can be extracted from well-defined relations. The details of this calculation and the results will be described in a separate publication. In this paper we present the main results which can be compared to the numerical simulations.
The two correlation functions which appear in the dynamical study are
$`G(t,t^{})=<\varphi (t)\eta (t^{})>`$
and
$`C(t,t^{})=<\varphi (t)\varphi (t^{})>`$
Both of these can be obtained from a single super-field correlation function ($`Q(1,2)`$) which is related to $`\mathrm{\Gamma }_{dyn}`$ through
$$Q^1(1,2)=\frac{\delta ^2\mathrm{\Gamma }_{dyn}}{\delta \mathrm{\Phi }_1\delta \mathrm{\Phi }_2}$$
(11)
where $`\mathrm{\Phi }`$’s are the super fields and the derivative is evaluated with the source term set to zero. Applying this formalism to Eq. (2), leads to an expression for $`\mathrm{\Gamma }_{dyn}`$ which involves two-point correlation functions of the fluctuations and is a natural extension of the static $`\mathrm{\Gamma }(\overline{\varphi })`$ to correlation functions involving time. The expression, however, does not lend itself to easy analysis except in two cases, early times when a variant of linear theory can be applied and late times when the system is stationary. The late time dynamics involves the analysis of the nonlinear terms and will be discussed in the concluding section. The early time dynamics is where one is justified in retaining only quadratic terms in the renormalized $`\mathrm{\Gamma }_{dyn}`$. In this limit we obtain the following equation for the equal-time correlation function $`C_q(t,t)=<\varphi _q(t)\varphi _q(t)>`$ which is the structure factor that is monitored in the simulations:
$`C_q(t,t)`$ $`=`$ $`{\displaystyle _0^t}[G(t,t^{})]^2𝑑t^{}`$ (12)
$`G_q(t,t^{})`$ $`=`$ $`{\displaystyle _t^t^{}}(r(t^{\prime \prime })+(q^2q_0^2)^2)𝑑t^{\prime \prime }.`$ (13)
The mass parameter, $`r(t)`$, which is now time-dependent, is renormalized in the dynamic theory using the same approximation as in the static theory. The diagrams used are given in fig. 7. These again are only to one loop. The mass term then becomes
$`r(t)+(q^2q_0^2)^2`$ $`=`$ $`\tau +(q^2q_0^2)^2+`$ (15)
$`{\displaystyle \frac{\lambda }{2}}[{\displaystyle \frac{d𝐪}{(2\pi )^d}\frac{1}{(\tau +(q^2q_0^2)^2)}}{\displaystyle \frac{d𝐪}{(2\pi )^d}\frac{\mathrm{exp}\{2(\tau +(q^2q_0^2)^2)t\}}{(\tau +(q^2q_0^2)^2)}}]`$
Again, the approximation is made self consistent by replacing the $`\tau `$ with $`r`$ in the integrand.
$$r(t)=r_{eq}\frac{\lambda }{2}\frac{d^3q}{(2\pi )^3}\frac{\mathrm{exp}(2D(q)t)}{D(q)}$$
(16)
$`D(q)`$ here is the renormalized propagator, $`D(q)=r_{eq}+(q^2q_0^2)^2`$ where $`r_{eq}`$ denotes the static renormalized mass parameter which is the solution to eq. 7.
In Eq. 16, at long times the integrand in the second term becomes small and $`r(t)`$ approaches $`r_{eq}.`$ which is confirmed in the section on statics above. At time zero the subtraction of the second term is just a subtraction of the Hartree level correction factor to the bare parameter introduced earlier Eq. (6 and, therefore, $`r(t=0)`$ is just the bare parameter $`\tau `$, while as $`t\mathrm{}`$, $`r(t)`$ approaches the renormalized, equilibrium value.
This scenario is confirmed by the simulations. The simulation results discussed in Section III B confirm that the equilibrium value, $`r_{eq}`$, is consistent with the Hartree prediction. Analyzing the early stage dynamics can verify the value of $`r(t)`$ at short times. Figure 5 shows the growth of $`S(q_0)`$ as a function of time for various values of the bare control parameter, $`\tau `$. To average over noise, five independent runs where taken for each value of the control parameter. If $`r(t)`$ is time-independent, the early-stage evolution (linear theory) is described by
$$S(q,t)=C_q(t,t)=S(q,0)\mathrm{exp}\{D(q)t\}+\frac{1}{D(q)}(1\mathrm{exp}\{D(q)t\})$$
(17)
where $`D(q)`$ is $`D(q)=r_0+(q^2q_0^2)^2`$ and $`r_0`$ is the time-independent value of $`r`$. At the peak of the structure factor, $`q=q_0`$ , and $`D(q_0)`$ is just $`r_0`$, the value of which can be estimated by fitting the simulation results shown in Fig. 5 to 17. The results of such fits are presented in Fig. 6. This figure shows that $`r_0`$ and $`\tau `$ are linearly related with a slope which is nearly one. This implies that for short times, the dynamically renormalized parameter is linearly related to the bare parameter and is not renormalized to positive values for negative values of the bare parameter. The very early-stage dynamics, therefore, is characteristic of a system exhibiting unstable growth.
These fits must be considered with some care. In the case where $`r`$ is not time dependent, the linear theory describes the system only for short times, that is times on the order of the natural time scale, $`r^1`$ \[See for example. For the fastest growing mode, when the data is fit to different time scales the results obtained are consistent for all time scales up to $`r^1`$. For the S-H equation, however, the results obtained are dependent on the time range. In particular, as the range of time increases the values obtained for $`r`$ become less negative. This is consistent with $`r(t)`$ growing as a function of time and the results of the fits merely represent some average value of $`r(t)`$ for the time range involved. If the system has been quenched to a negative value of $`\tau `$, Eq. 16 indicates that the early-time evolution will exhibit unstable growth. As time evolves the second term in Eq. (16) decreases and the value of $`r(t)`$ approaches $`r_{eq}`$ which is always positive. Though the integral in Eq. 16 is hard to evaluate analytically, numerics can provide some insight into its behavior. Figure 8 shows a numerical evaluation of the time evolution of $`r(t)`$ for various values of $`r_{eq}`$ and $`\lambda =0.01`$. As confirmed in the statics section above, $`r(t)`$ will eventually become positive no matter how deep the quench is, though the time for this to happen may become very long. As long as $`r(t)`$ is negative the system is expected to undergo unstable growth though, since the mass parameter is time-dependent, the time evolution will be more complicated than the usual spinodal decomposition scenario. We can define a crossover time, $`t_{cross}`$ at which $`r(t)`$ changes sign, and an average growth time, $`t_{ave}`$. The average growth time is just the inverse of $`r_{ave}`$ defined as
$$t_{ave}^1=r_{ave}=\frac{1}{t_{cross}}_0^{t_{cross}}r(t^{})𝑑t^{}$$
(18)
In Fig. 9, we show plots of $`t_{cross}`$ and $`t_{ave}`$ as a function of the bare control parameter. For shallow quenches $`t_{cross}`$ is small compared to $`t_{ave}`$ and the system reaches a well-defined metastable equilibrium state which is disordered. If $`t_{cross}`$ becomes larger then $`t_{ave}`$, metastability becomes difficult to define. The system takes longer to equilibrate in the metastable disordered phase than it takes to grow lamellar structures. We can use the condition $`t_{cross}=t_{ave}`$ to define a crossover temperature $`\tau _{dyn}^{}`$. Above this temperature the disordered phase quickly becomes locally stable and a nucleation event is needed for the formation of the lamellar structures. Below $`\tau _{dyn}^{}`$, the system is expected to evolve continuously towards the lamellar phase without any evidence of metastability of the disordered phase. Fig. 9 is in sharp contrast to standard Ginzburg-Landau (G-L) theory. For the time dependent G-L equation, the presence of noise suppresses the critical point. In the region where the bare parameter is negative but the renormalized parameter is positive, $`t_{cross}`$ is always smaller then $`t_{ave}`$ and so there is no unstable growth until the true critical point is reached.
For the system under consideration here, different scenarios are possible depending on the relationship of this dynamical crossover temperature to the first-order transition temperature, $`\tau _{trans}^{}`$ , at which the free energy of the lamellar phase becomes lower that that of the disordered phase. If $`\tau _{trans}^{}`$ is higher than $`\tau _{dyn}^{}`$, there will be a regime of temperatures over which the the system will undergo nucleation and growth. On the other hand if $`\tau _{trans}^{}`$ is lower than $`\tau _{dyn}^{}`$, then there is no nucleating regime and one would observe only unstable growth, albeit of an unusual nature since the free-energy surface is evolving with time. From figure 9 we can see that $`t_{cross}`$ is small only for very shallow quenches below the mean-field transition. The results of our simulations, presented in the next section, are in qualitative agreement with this scenario. We would like to emphasize that the growth dynamics described above are qualitatively different from rapid nucleation in which the system quickly reaches the stable equilibrium state. A large value of $`t_{cross}`$, on the other hand, indicates a type of ergodicity breaking as the system takes a very long time to reach equilibrium.
All of our numerical results indicate that $`\tau _{trans}^{}`$ predicted from Hartree is lower than $`\tau _{dyn}^{}`$ and they are remarkably close to each other. We have been unable to come up with a simple relationship between these two temperatures and, therefore, can only interpret the similarity of the two as a remarkable coincidence. The dynamical crossover is deduced from time scales which characterize the evolution of the free-energy surface while the transition temperature is deduced from a comparison of the depth of the two wells. It is not clear why the two temperatures should be similar in magnitude. It can be argued that there is only a single parameter, $`\lambda `$, controlling the scale of fluctuations and, therefore, the two temperatures should be related, however, there is no obvious argument to suggest that they should be identical.
The picture emerging from the dynamical renormalization is a natural extension of the effect of fluctuations on the static results of the Brazovskii model. Following a quench from a relatively high temperature to a temperature where $`\tau `$ is negative, the system is in a locally unstable region (top of a hill). As the fluctuations grow with time, the non-linear terms characterized by $`\lambda `$ become important and they renormalize the curvature of the hill. This scenario is quite different from the usual picture of evolution in an adiabatic potential.
### C Simulation Results for Late Times
Further evidence for the dynamical scenario presented in the previous section comes from examining the long time evolution of the peak of the structure factor obtained from our numerical simulations and correlating that to snap shots of the system as it evolves. Figure 10 is a plot of the amplitude of the structure factor peak as it evolves. For shallow quenches, $`\tau >0.04,`$the peak grows to an equilibrium value and does not evolve further. These values are consistent with the equilibrium values reported above in the section on statics. For deeper quenches the peak amplitude grows quickly and then appears to saturate; however careful examination shows that the value continues to grow very slowly. This is consistent with late stage domain growth which as been studied extensively by Elder et. al. for 2-D systems . The last feature in the system is a final rise to an equilibrium value. This rise is a finite size affect is caused by the majority domain finally taking over the entire system. This being the case, the peak would not be expected to grow much since the difference between the value before and after this final evolution should only represent the surface area between the different domains.
For a quench depth of $`\tau =0.07`$ this time evolution can be compared to a series of system snapshots shown in figures 11\- 15. These figures represent $`\varphi =0`$ isosurfaces which would be the boundaries between the different microphases of the system. For early times the system appears to be very disordered while at later times domains of ordered lamellar structures begin to appear. Just before the final convulsive growth, at a time $`t=4000`$, there still appear to be a few different ordered domains while just after the final growth, at $`t=5000`$, only one domain appears to be left in the system. Thus we still have a consistent picture of slow but continuous domain growth in the system which is governed by the evolution of the 2nd order term as it goes from negative values to positive values.
## V Future Work
In the work presented above we have shown that analysis of the static free energy does not always provide an adequate description of the fluctuation driven dynamics. Though figure 10 shows features which are reminiscent of a first order phase transition, dynamical analysis of the renormalized coefficients suggest that a more complicated evolution is taking place. A similar situation occurs for the superconducting transition which is also a fluctuation-driven first-order transition.
In the present work we have paid close attention to the time evolution of the mass parameter and its effect on the evolution of the system at early times. The dynamical renormalization results also become tractable at times large compared to $`t_{cross}`$ when the system has reached equilibrium. As pointed out earlier, one enters this regime fairly quickly for shallow quenches. At this stage, the higher order terms in $`\mathrm{\Gamma }_{dyn}`$ become important. The most interesting aspect of the renormalized theory is that the fourth-order term becomes non-local in time :
$`\mathrm{\Gamma }_{dyn}(4){\displaystyle d1d2\mathrm{\Phi }^2(1)K(12)\mathrm{\Phi }^2(2)}`$
and the integral of the kernel over all times leads to the Brazovskii result for the renormalized fourth-order term. One can derive an “effective” Langevin equation for the fields $`\varphi `$ from the renormalized $`\mathrm{\Gamma }_{dyn}`$, and the above equation implies that this Langevin equation has “memory” effects which appear in the non-linear term. This would lead to unusual behavior of the equilibrium correlations functions and might provide a distinctive signature of fluctuation-driven first-order transitions. We are in the process of investigating the effective equation for $`\varphi `$ numerically. One obvious consequence of the non-local term is that the barrier to nucleation is dynamic and it is not valid to think of nucleation as tunneling through a fixed barrier. The situation is closer to the problem of quantum tunneling for many degrees of freedom where one also finds an effective equation which is non-local in time.
Another area of interest is the role of defects in the phase transition of the system. These will be important for very shallow quenches to just below the transition point and we have noted there presence when looking at snapshots of the system. These defects are reminiscent of the perforated lamellar phase described in reference . One way to study this is to use a directional order parameter measure introduced by Bray et al. to study the 2D problem. This order parameter is similar to the order parameter for complex fluids and provides directional as well as density information. What it may show is that the local gradient terms are enhanced for negative values $`\tau `$ suggesting that perhaps there is a disordered to nematic transition in this region which others have suggested for 2D .
We would like to close with a discussion of possible experimental systems where this dynamical scenario can be studied experimentally. Since the effective $`\lambda `$ is extremely small in Rayleigh-Benard systems the regime where the disordered state is metastable will probably be unaccessible. In di-block copolymers, however, there should be a significant temperature range, below the mean-field transition, where the disordered state is locally stable and a study of the equilibrium, two-time correlation function should reveal the existence of the non-linear memory term. It might also be possible to observe unusual nucleation if the system parameters allow for $`\tau _{trans}^{}`$ to be higher than $`\tau _{dyn}^{}`$.
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# The Narrow Line Region of Narrow-Line Seyfert 1 GalaxiesBased on observations made at CASLEO. Complejo Astronómico El Leoncito (CASLEO) is operated under agreement between the Consejo Nacional de Investigaciones Científicas y técnicas de la República Argentina and the National Universities of La Plata, Córdoba and San Juán.
## 1 Introduction
Narrow-Line Seyfert 1 Galaxies (hereafter NLS1) are a peculiar group of AGNs where the permitted optical lines show full width half-maximum (FWHM) not exceeding 2000 km s<sup>-1</sup>, the \[O III\] $`\lambda `$5007/H$`\beta `$ ratio is $`<`$ 3 and the UV-VIS spectrum is usually very rich in high ionization lines and Fe II emission multiplets. In the soft X-ray band, NLS1s have generally much steeper continuum slopes and rapid variability (Boller, Brandt & Fink 1996, hereafter BBF96). Recently, Leighly (1999) found that the hard X-ray photon index is significantly steeper in NLS1s compared with that of normal Seyfert 1s, and that soft excess emission appears considerably more frequently in NLS1s than in Seyfert 1 (hereafter Sy1) galaxies with broad optical lines.
It is not known at present the origin of the narrowness of broad permitted lines in NLS1s. Osterbrock & Pogge (1985); Ulvestad, Antonucchi & Goodrich (1985) and Stephens (1989) suggest that if the velocities in the BLR of Seyfert 1s were largely confined to a plane, the NLS1 galaxies could be understood as cases in which the line of sight is nearly perpendicular to this plane. BBF96 state, on the other hand, that if the gravitational force from the central black hole is the dominant cause of the motions of Seyfert BLR clouds, narrower optical emission lines will result from smaller black hole masses provided the characteristics BLR distance from the central source does not change strongly with black hole masses.
However, Rodríguez-Pascual, Mass-Hesse & Santos-Lléo (1997) report the detection in NLS1 galaxies of broad components with FWHM around 5000 km s<sup>-1</sup> for the high ionization UV permitted lines such as Ly$`\alpha `$, C IV $`\lambda `$1550 and He II $`\lambda `$1640. This result indicates that gas moving at velocities comparable to those found in typical Sy1 galaxies does indeed exist in NLS1s. In the optical region, they found “broad” components with FWHM less than 3000 km s<sup>-1</sup>, narrower than the broadest UV component in the same objects. Nonetheless, deblending the optical permitted lines in NLS1s is difficult because no transition between the narrow and broad components is observed. This shortcoming has strong influences in, for example, the analysis of the narrow line region (NLR) due to the large uncertainties in determining the fraction of H$`\alpha `$ and H$`\beta `$ which originates from low ionization material.
Up to now, most studies of the NLR in NLS1s assume that the flux emitted by the narrow H$`\beta `$ equals 10% of the flux of \[O III\] $`\lambda `$5007 (Osterbrock & Pogge 1985; Leighly 1999). This assumption is based on the results obtained from Seyfert 2 and intermediate Sy1 galaxies (e.g. Koski 1978; Cohen 1983). But in recent years, growing observational evidence points out to the existence of differences between the NLR of normal Sy1 and Sy2 galaxies (Schmitt & Kinney 1996; Schmitt 1998), making the above assumption highly uncertain. In addition, fixing the \[O III\] $`\lambda `$5007/H$`\beta `$ ratio to 10 implies ignoring the large scatter in the value of this ratio observed in normal Sy1s (2 to 19, see for example Rodríguez-Ardila, Pastoriza & Donzelli 1999, hereafter Paper I) and overlook the influences that this ratio could have in the energetics and physical conditions of the NLR of these objects.
Due to the above reason, the main purpose of this paper is to seek additional constrains in order to estimate the actual contribution of the narrow H$`\beta `$ flux to the total H$`\beta `$ emission line and study the implications that the newly adopted values could have on narrow line ratios and the physics of the NLR of NLS1s.
The present work is organized as follows. In Section 2 we describe the sample of NLS1s used in this paper. Section 3 presents the decomposition into narrow and broad components carried out in the Balmer lines of the NLS1 galaxies. Photoionization models that successfully reproduce the observed line ratios of NLS1s are presented in Section 4. A discussion of the main results appears in Section 5 and the conclusions are presented in Section 6.
## 2 Observations
Long-slit spectroscopic observations of seven NLS1 galaxies covering the spectral region 3700 Å - 9500 Å were obtained with the 2.15 m telescope of the Complejo Astronomico El Leoncito (CASLEO) using a TEK 1024 $`\times `$ 1024 CCD detector and a REOSC spectrograph. Two gratings of 300 l/mm with blaze angles near 5500 Å and 8000 Å, respectively, were used in order to fully cover the spectral interval 3700 – 9500 Å. The spatial scale of this setup is 0.95″/pix, with an instrumental resolution of 7 Å FWHM. A slit width of 2″.5 oriented in the East-West direction and crossing the center of the galaxies was employed. The galaxies and standard stars were observed near the zenith (air masses $`<`$ 1.2). A complete log of the observations and reduction procedure are described in Paper I.
CTS J03.19, CTS J04.08, CTS J13.12 and CTS H34.06 were classified as NLS1s by us based on the appearance of their respective optical spectrum. The criterion used in this classification is that the permitted lines be slightly broader than the forbidden lines, following the definition of Osterbrock & Pogge (1985). MRK 1239, CTS R12.02 (=NGC 4748) and 1H 1934-063 are well known NLS1s objects. For this reason, we consider that our sample is optically selected, in contrast to most NLS1s published in the literature, which are based on X-ray selected objects. The luminosities of the galaxies are low to intermediate and their radial velocities are not larger than 16,000 km s<sup>-1</sup>. Figure 1 shows the spectra of the seven galaxies of the current study, already corrected for redshift.
Fe II emission, which is particularly strong in NLS1s, may alter the flux and width of H$`\beta `$ \+ \[O III\] $`\lambda `$$`\lambda `$4959,5007. For this reason we have carefully removed the Fe II multiplets following the method described in Boroson & Green (1992). It consists of constructing a Fe II template by removing the lines which are not of Fe II from the spectrum of IZw1, a NLSy1 galaxy widely known for the strength of the Fe II emission and the narrowness of its “broad lines.” For this purpose, a high signal to noise spectrum of IZw1 covering the spectral range 3500-6800 Å was taken in one of the observing runs. After isolating the Fe IIemission, the template was broadened by convolving it with a Gaussian profile having a FWHM similar to that of the H$`\beta `$ line and scaled to match the observed Fe II emission of the corresponding galaxy. Figure 2 shows the above procedure for 1H 1934-063 (upper panel), CTS R12.02 (middle panel) and MRK 1239 (lower panel). In each case it is shown the observed spectrum, the Fe II template and the resulting spectrum after removing the Fe II emission. It can be seen that the most important Fe II features at both sides of H$`\beta `$ have been cleanly removed leaving as a residual the true He II $`\lambda `$4686 line profile and, in MRK 1239, the \[FeIXV\] $`\lambda `$5302 line.
## 3 The Gaussian Description of The Emission Line Profiles
### 3.1 Results of line profile fitting
In order to characterize the emission line profiles of the NLS1s we have assumed that they can be represented by a single or a combination of Gaussian profiles. The liner routine (Pogge & Owen 1993) which is a $`\chi ^2`$ minimization algorithm that fits as many as eight Gaussians to a line profile, was used for this purpose.
As a first steep we tried to fit the H$`\alpha `$ emission line with a single Gaussian component. H$`\alpha `$ is the strongest optical permited line and is located in a spectral region with the highest S/N. It is therefore the best place to look for broad components to the permitted optical lines of NLS1s. The best fit obtained is shown in the left panels of Figure 3. The thick line represents the synthetically calculated profile and the dotted line the residuals of the fit. It is clear that this simplest representation cannot adjust adequately the wings of H$`\alpha `$ although its core is nicely fitted in most of the objects. We then tried to adjust a second Gaussian to H$`\alpha `$. The best solution found are shown in the right panels of the same figure. Now the residuals of the fit are quite similar to the noise level around H$`\alpha `$, implying that a narrow (FWHM $``$ 600 km s<sup>-1</sup>) plus a broad component (FWHM $``$ 2500 km s<sup>-1</sup>) give a convincingly better description of the observed profiles. In the above fits the only constraints applied were that \[N II\] $`\lambda `$6548 and \[N II\] $`\lambda `$6584 be of equal FWHM and that their flux ratio and wavelength separation be equal to their theoretical value (1:3 and 36 Å, respectively).
The same decomposition was applied to H$`\beta `$ \+ \[O III\] $`\lambda `$4959,5007. As for H$`\alpha `$, the fit of H$`\beta `$ with a single Gaussian gives a poor representation of this line and it was necessary to include an additional component to represent adequately the observed profile. Figure 4 shows the results of this decomposition. The residuals (dashed line) of the dual component fit (right panels) are significantly improved with respect to the fit with a single Gaussian (left panels).
Columns 2 to 7 of Table 1 list the FWHM (in km s<sup>-1</sup>) of the emission lines measured from the Gaussian fitting. These values were obtained from the subtraction, in quadrature, of the observed FWHM and that of the instrumental profile, measured from the comparison lamp lines (FWHM $``$ 360 km s<sup>-1</sup> at H$`\alpha `$). The flux ratio of the narrow to the broad component of H$`\beta `$ for each galaxy is listed in Column 8 and the \[O III\] $`\lambda `$5007 flux, relative to the narrow H$`\beta `$ component, is in Column 9. In all but one case (MRK 1239), the center of the narrow components was coincident with the systemic velocity of the corresponding galaxy. MRK 1239 presents a second blueshifted (with an outflow velocity of 600 km s<sup>-1</sup> with respect to the nucleus) broader (FWHM $``$ 1580 km s<sup>-1</sup>) component that is also present in H$`\alpha `$ and H$`\beta `$ (see Figures 3 and 4) and that we associated to the NLR .
It is important to stress that the same result, –i.e a narrow component plus a broad one in the permitted lines of H$`\alpha `$ and H$`\beta `$ – would be obtained if instead of using the Gaussian decomposition technique we had employed the \[O III\] $`\lambda `$5007 line as a template representative of the NLR profiles. This test was applied to the H$`\beta `$ line of the galaxy sample in order to see if the broad component found in the permitted lines using the Gaussian representation were not an artifact of the fitting procedure. For this purpose the \[O III\] $`\lambda `$5007 line was isolated and normalized to the peak intensity of the H$`\beta `$ line to represent the maximum allowed H$`\beta `$ contribution from the NLR to the observed profile. When subtracted, the residuals consist basically of a broad wing and a strong absorption coincident with the peak position of the template profile, such as shown in Figure 5 for the NLS1 galaxies 1H 1934-063, CTS R12.02, CTS J13.12 and MRK 1239.
We interpret the broad wings as a clear evidence of the presence of a broad component similar to that observed in normal Seyfert 1 galaxies but with a smaller FWHM, supportting the results of Rodríguez-Pascual, Mass-Hesse and Santos-Léo (1997); Gonçalves, Véron & Véron-Cetty (1999) and Nagao et al. (1999) who also report a broad component in the permitted lines of NLS1s. The absorption is interpreted as due to an overestimation of the NLR component. Scaling the template profile in order to eliminate the absorption in the residuals leaves a pure broad component very similar in intensity and width as that found in the Gaussian decomposition (dashed line of Figure 5). In fact, the values of FWHM and line fluxes obtained using this approach is essentially the same as those obtained formely.
We conclude from the above test that the broad feature observed in H$`\beta `$ and H$`\alpha `$ is real and it is interpreted here as the contribution from the BLR to the H$`\beta `$ line.
Some authors (Moran, Halpern & Helfand 1996, Gonçalves, Veron-Cetty & Veron 1999) have suggested that NLS1s have more nearly Lorentzian, rather than Gaussian, profiles, as evidenced by their cusped peaks and broad wings, mainly in H$`\beta `$ and \[O III\] $`\lambda `$5007 lines and even in H$`\alpha `$. We have tested this hypothesis by fitting one-component Lorentzian profiles to the Balmer and adjacent lines. The best solution, plotted in Figure 6 for CTS R12.02, 1H 1934-063 and MRK 1239 shows that this description does not provide a satisfactory fit to the emission lines. Very similar results were obtained for the other NLS1s of our sample. In all cases, the wings of the Lorentzians are more extended than the wings of the observed profiles. This effect is more pronounced in H$`\alpha `$ (left panels of Figure 6) than in H$`\beta `$ (right panels), giving the apparent impression that this latter line is better fitted by a Lorentzian profile. We attribute this effect to the lower S/N of the H$`\beta `$ region, which reduces the contrast of faint broad component wings relative to the narrow line. The poor fit obtained for the \[O III\] $`\lambda `$$`\lambda `$4959,5007 lines demonstrate that the Lorentzian profiles are unsuitable for representing the NLR profiles of NLS1s.
At this point it is important to note that the choice of the Gaussian profile as representative of the form of the observed emission lines was due to its simplicity and lack of physical reasons to adopt another particular form. We must be aware that the decomposition of the permitted lines into narrow and broad components, particularly in NLS1s, is a very uncertain task. As Evans (1988) noted, the fitting of symmetrical functional forms to the observed profiles may lead to the appereance of “components” in the center and wings of the lines but that, in certain circumstances, cannot bear any physical meaning. In addition, the observed profile can be the result of the superposition of many emitting regions located along the line of sight, each of them with a different intrinsic profile. For these reasons, one has to look with caution the deblending into broad and narrow components obtained for the permitted lines. Nonetheless, the fact of having obtained similar results using two independent methods give additional support to our interpretation.
### 3.2 Meaning of the Narrow and Broad Components in NLS1s
In the previous section we showed that the permitted lines of NLS1s can be decomposed into a narrow and a broad component, as is usually carried out in normal and intermediate Seyfert 1 galaxies. But what is the real meaning of this? A Gaussian representation as such, is only a mathematical way of fitting approximately the data. Additional evidence and a rationale are needed in order to identify each Gaussian component with physically different regions. Also, we must be aware that the number of Gaussian components adjusted to a given line basically depends on the spectral resolution and the S/N within the region of the fit.
Two important points call our attention regarding the values presented in Table 1. First, the FWHM of the “broad” components in H$`\alpha `$ and H$`\beta `$ are rather similar within the same galaxy and from object to object. Second, the FWHM of the narrow component of permited lines is similar to that of the forbidden (\[O III\] and \[N II\]) lines. They are also larger than the width of the instrumental profile (FWHM $``$ 360 km s<sup>-1</sup>).
It is generally accepted that the kinematics of the BLR clouds is largely dominated by the gravitational potential of the central black hole. The FWHM of the lines emitted in that region is a measure of the dispersion of velocities of the emitting gas, along the observer’s line of sight, which represents also the depth of the gravitational potential well in which the emitting clouds find themselves. For this reason, it is expected that lines formed in the same spatial region will have similar FWHM and emission line profiles.
A similar argument can be applied to the narrow lines. However, the NLR is much more extended and located farther out from the central source (1–100 pc) so these clouds are immersed in a shallower gravitational potential dominated by the galaxy bulge.
Therefore, the consistency of the width and profile form of the broad components of H$`\alpha `$ and H$`\beta `$ and the similarity in width of the narrow permitted and forbidden lines, within the same galaxy, allow us to associate the narrow and broad Gaussian components of the permitted lines to the integrated line emission from the NLR and BLR, respectively.
A comparison of the FWHM of the narrow lines of our sample with those measured with a similar setup in normal Seyfert 1s (Cohen 1983; Stephens 1989; Puchnarewicz et al. 1997) shows that no differences seem to exist in the NLR kinematics between these two groups of galaxies.
The main difference lies in the BLR. The FWHM of the broad component is not only significatively smaller than that of the normal Sy1s but also its relative contribution to the total flux of the line is greatly reduced in these objects. In Column 8 of Table 1 we have listed the ratio of the flux associated to the narrow component to that of the broad one. It can be seen that on average, this ratio is very near unity, meaning that 50% of the total line flux is due to the narrow component. For comparison, this ratio is around 0.1 in typical Sy1 galaxies.
The above results by themselves do not represent a major departure from our current picture of NLS1s. However, a significant difference emerges when we consider line ratios between the narrow component of H$`\beta `$ and \[O III\] $`\lambda `$5007. In effect, since Osterbrock & Pogge (1995) and up to very recently (Leighly 1999), it has been assumed that the contribution in flux of the narrow component of H$`\beta `$ to the NLR spectrum of a given object equals 10% of the flux of \[O III\] $`\lambda `$5007. This assumption is based on the fact that \[O III\] $`\lambda `$5007/H$`\beta `$ is, on average, $``$ 10 in Seyfert 2 galaxies (Veilleux and Osterbrock 1987). Nonetheless, growing evidence of important differences between the NLR of Sy1 and Sy2 galaxies have appeared in the literature (Schmitt & Kinney 1996, Schmitt 1998; Paper I). In addition, NLS1s are recognized as a subclass within the realm of Seyfert 1s due to their peculiar properties, making very unlikely that the above assumption really holds in these objects, as is shown in column 9 of Table 1, which lists the \[O III\] $`\lambda `$5007/H$`\beta `$(narrow) ratio found from our decomposition of line profiles. These values also show that there is a wide intrinsic dispersion in the \[O III\] $`\lambda `$5007/H$`\beta `$ ratio of NLS1s, ranging from 0.8 to 5.
Since the new values of the narrow \[O III\] $`\lambda `$5007/H$`\beta `$ ratio are drawn from the decomposition into narrow and broad components of the Balmer lines, it is important to discuss the uncertainties inherent to the deblending process. Strictly speaking, one should be inclined to treat the flux associated to the narrow component of the permitted lines as a lower limit, being the actual flux between this value and that obtained from the total flux of the line. The latter option would be the case if no contribution from the BLR were present, as was initially suggested by Osterbrock & Pogge (1985) in order to explain the absence of broad permitted lines in NLS1s. Under the last circumstance, the resulting \[O III\] $`\lambda `$5007/H$`\beta `$ ratio would be even smaller but would not change drastically. It would now fall into the interval 0.5 – 1.7, making the departure from the ratios found in normal Seyfert 1 galaxies even stronger.
One can argue that the deblending of the permitted lines overestimated the flux of the narrow H$`\beta `$ component, making the \[O III\] $`\lambda `$5007/H$`\beta `$ ratio appear smaller that it really does. That is, instead of being a lower limit, it represents an upper limit. Under this circumstance, the actual flux of the narrow H$`\beta `$ would range between this upper limit and a small fraction of the total line flux. Physically, this is highly improbable, at least for two reasons. First, H$`\beta `$ is emitted by every gas component, so the narrow H$`\beta `$ cannot be narrower than the narrowest NLR line. This is in accord to the values of Table 1, where the FWHM of H$`\beta `$ is similar to that of \[O III\] $`\lambda `$5007 and of the same order of \[N II\] $`\lambda `$6584.
Second, lets suppose that the actual narrow H$`\beta `$ flux corresponds to 50% of value obtained from the deblending process. With this in mind, the \[O III\] $`\lambda `$5007/H$`\beta `$ ratio would now fall in the interval 1.6 – 10. But now a new problem will arise: the reddening measured from the Balmer decrement would be increased up to 0.7 mag. Such an increase in the E(B-V) would have a strong influence in the resulting emission line spectrum. Probably new physical processes and ionization mechanisms should need to be invoked in order to explain the observed spectrum. In addition it would not be clear in which emitting region the remaining “narrow” flux would be produced. In the BLR? in the intermediate NLR?
In the following sections, additional evidence supporting the results obtained from the deblending process is given. Nonetheless, it is clear that a large spatial resolution, such as that obtained with the HST, is necessary in order to examine carefully the NLR of NLS1 galaxies.
## 4 Matter-Bounded and Ionization Bounded clouds in the NLR of NLS1s
We have found that the NLR of NLS1s are characterized by lower \[O III\] $`\lambda `$5007/H$`\beta `$ ratios, more typical of some starburst or H II galaxies than of canonical Seyfert 1s or 2s. In addition, it was shown in Paper I that i) NLS1s have intrinsically weak low ionization forbidden lines; ii) the dominant mechanism of the NLR emission is photoionization by a central source and iii) a wide range of densities (10<sup>3</sup> cm<sup>-3</sup> – 10<sup>6</sup> cm<sup>-3</sup>) must exist in the NLR of these objects, with the larger densities associated to the inner NLR regions where \[O III\] and the high ionization lines are emitted while the lowest values are associated to the \[O I\] and \[S II\] emitting regions.
A scenario which springs naturally from the above results invokes the presence of at least two types of clouds. A denser, inner set of high excitation matter-bounded (MB) clouds are required to enhance the temperature sensitive \[O III\] 4363Å line relative to 5007Å (due to collisional deexcitation) and also matter-bounded in order not to emit low excitation lines (from low ionization species such as O I, S II, N II…). The MB component is where most of the \[O III\] emission and high ionization lines such as He II, \[Ne V\], \[Fe VII\] and \[Fe X\] would be produced. To account for the lower excitation lines, it is necessary to consider a second component consisting of outer, less dense set of ionization-bounded (IB) clouds, responsible for the production of low-ionization lines such as \[O I\], \[N II\] and \[S II\] (whose doublet ratio points to low densities) and some \[O III\] as well. These clouds must lie at a much larger distance in order that the ionization parameter be much smaller than the MB component as discussed in Rodríguez-Ardila, Pastoriza & Maza (1998, hereafter RAPM). Our starting hypothesis is that the relative proportion of these two types of clouds combined with a suitable choice of the input ionizing continuum will reproduce the observed differences between the NLR of NLS1s and that of normal Seyfert 1 galaxies.
This dual-component model have been successfully applied before to account for the observed emission line ratios of the NLR and extended narrow line region of many Seyfert 1 and 2 galaxies (Binette, Wilson & Storchi-Bergmann 1996, Binette et al. 1997, RAPM, Aita-Fraquelli, Storchi-Bergmann & Binette 1999). Such a simple model is surely an over-simplification relative to the wide range in parameters needed to fully describe the NLR (to illustrate the complexity of the problem see Moore & Cohen 1996, Ferguson et al. 1997 and Komossa & Schulz 1997). However, in the absence of sufficient observational constraints on all possible models and geometries of the NLR, it appears to us to be the simplest mean for taking into account the particular emission line signature reported in RAPM, that is a high density high excitation spectrum observed in conjunction with the presence of low density low excitation lines.
### 4.1 The Photoionization Models
In this section we test whether the simplified dual-component MB-IB description can satisfactorily account for the anomalous line ratios observed in NLS1s. We are particularly interested in determining which kind of spectral energy distribution (SED) observed in AGNs can best produce a \[O III\] $`\lambda `$5007/H$`\beta `$ ratio in the range 1 – 5 and the strongest optical emission lines, using similar physical parameters to those employed to model the NLR of normal Seyfert 1 galaxies. A detailed photoionization modeling of the NLS1 emission line ratios is in preparation (Rodríguez-Ardila, Binette & Pastoriza 1999)
Although there is no direct way to determine the intrinsic shape of the ionizing continuum in the EUV domaine, recent work by Zheng et al. (1997) and Laor et al. (1997) show that the ionizing continuum, from the Lyman limit to soft X-ray energies, may be characterized by a power-law of index $`\alpha 2`$. This result was derived from QSOs of intermediate redshifts but may be extended to lower luminosity AGN such as Seyfert 1 galaxies. Korista, Ferland & Baldwin (1997) suggest flatter indices $`\alpha 1.5`$, which have been used by Binette, Wilson & Storchi-Bergmann (1996) and RAPM, for instance, in their modeling of Seyfert 1 objects. Meanwhile, ROSAT observations of NLS1s show that the soft X-ray photon index is systematically steeper than that of Seyfert 1 galaxies with broad optical lines (BBF96; Foster & Halpern 1996; Laor et al. 1997).
Due to the above reasons, we generated sequences of models assuming a SED of broken power-laws of the form F<sub>ν</sub> = K$`\nu ^\alpha `$ where:
$$\alpha =1.4,13.6eVh\nu 1300eV;\alpha =0.4,h\nu 1300eV$$
(1)
$$\alpha =2.2,13.6eVh\nu 2000eV;\alpha =1.1,h\nu 2000eV$$
(2)
Equation 1 (SED 1) corresponds to the fit made by Kraemer et al. (1998) to the observed SED of the Seyfert 1 galaxy NGC 5548. Equation 2 (SED 2) uses the median values of the spectral indices found from ROSAT and ASCA data for NLS1s (Leighly 1999).
For each SED, the $`A_{M/I}`$ parameter, which characterizes the relative proportion of MB and IB clouds, was varied from 0.04 to 11. Density estimates of the NLR of NLS1s based on density and temperature line ratios (cf Paper I) show that MB clouds should have $`n_e10^6`$ cm<sup>-3</sup> while in IB clouds $`n_e10^3`$ cm<sup>-3</sup>. The ionization parameter $`U_{MB}`$ at the illuminated face of the MB clouds was initially set to 10<sup>-1.5</sup>, following estimates by RAPM, Kraemer et al. (1998) and Kraemer et al. (1999) for normal Seyfert 1 galaxies.
The multipurpose code MAPPINGS Ic (Ferruit et al. 1997) was used to compute the dual-component photoionization models. Plane-parallel geometry was assumed given the relatively large distance of the ionized clouds from the central source compared to the geometrical depth of the clouds. The gas is assumed to be atomic and the gas abundances solar. In order to let the inner set of clouds be matter-bounded, a fraction F<sub>MB</sub> of the input ionizing continuum is left to escape from the back of the clouds. This fraction was initially set to $`60`$% (the value used in Ferruit et al. 1997) but was allowed to increase substantially for the models making use of the steep SED. This was necessary in order to keep a high value of the He II/H$`\beta `$ ratio emitted by the MB component, an essential ingredient of the dual-component models. The “filtered” continuum is later used to ionize the IB clouds further out (which do not emit any He II). The IB clouds integration is stopped when the electron temperature falls below 5000 K. Dust is considered to be mixed with the ionized gas, and the code includes heating by dust photoionization. The dust content of the photoionized plasma is described by the quantity $`\mu `$ which is the dust-to-gas ratio of the plasma expressed in units of the solar neighborhood dust-to-gas ratio ($`\mu =1`$). In this work we use a constant $`\mu `$ = 0.015, which has a negligible effect on line transfer and on the gas heating.
Figure 7 shows the predicted \[O III\] $`\lambda `$5007, He II $`\lambda `$4686, \[O I\] $`\lambda `$6300 and \[O III\] $`\lambda `$4363, relative to H$`\beta `$ as function of the relative proportion of MB and IB clouds, A<sub>M/I</sub>. The solid line corresponds to the predictions of SED 1 and the dashed line to the ratios obtained with SED 2. In order to compare the output with the observations we have plotted the dereddened observed ratios for the NLS1 galaxies CTS H34.06 (filled triangle), 1H 1934-063 (filled square), CTS R12.02 (filled circle) and MRK 1239 (open square) as taken from Paper I. NGC 5548 was chosen as representative of the normal Seyfert 1 galaxies. Its dereddened ratios were taken from Kraemer et al. 1998 and are shown as filled diamonds. Table 2 lists the predictions of the best-fit model for the most important optical lines observed in each galaxy and the parameters of the corresponding model.
A close inspection to Figure 7 and the values of Table 2 allows us to say that our approach is successful at reproducing the observed ratios overall. The fact of being able to predict simultaneously and accurately the \[O III\] $`\lambda `$5007/H$`\beta `$ and He II $`\lambda `$4686/H$`\beta `$ ratios for the NLS1 galaxies and NCG 5548 indicates that the general distribution of ionizing photons at least in the range 13.6 – 100 eV is correct. Since the model parameters were rather similar for the two SEDs, the results provide a strong support to the idea that the NLR of NLS1s and Sy1s have similar physical properties in terms of density, chemical abundance, ionization parameter and distance to the central source. The most probable cause of the observed spectral differences can be related to the variation in steepness of their ionizing continuum.
The largest discrepancies between models and observations, for some of the galaxies, are in low excitation lines ratios such as \[O I\] $`\lambda `$6300/H$`\beta `$, \[S II\] $`\lambda `$$`\lambda `$6717,6731/H$`\beta `$ and \[N II\] $`\lambda `$6584/H$`\beta `$. These are overpredicted by a factor of 2 in CTS H34.06, 1H 1034-063 and MRK 1239.
This systematic overprediction, nonetheless, might provide helpful constraints to the physical conditions of the IB clouds and the nature of the central ionizing source. In effect, it is well known that the low ionization lines of AGNs are produced by the flux of high energy EUV photons in the range 300–900 eV (cf fig 1 in Binette et al. 1997) which penetrate deep into the NLR gas. If these photons are not allowed to reach this region, intrinsically weaker low ionization lines will be emitted. A possible candidate for this shielding might plausibly be a high ionization warm absorber (hereafter WA), located near the BLR. Evidences for the existence of such an absorbing material have been found in several NLS1s, particularly in those with high polarization (Leighly 1997). MRK 1239, the NLS1 with the largest discrepancies between the predicted and observed low ionization lines, (cf. Table 2) is widely known for being highly polarized (Goodrich 1989) and for showing WA spectral features (Leighly et al. 1997; Grupe et al. 1998). In contrast, CTS R12.02 (= NGC 4748) shows one of the lowest percentage of polarization (0.12% against 2.89% of MRK1239) according to Goodrich (1989) and presents here the best agreement between the model and the observed line ratios, even in the low ionization lines. No report about X-rays WA features are found in the literature for this galaxy.
Although polarization data for the remaining NLS1s are not available at present, the case of MRK 1239 can be taken as a good evidence for atenuation of the hard X-ray continuum and its influence on the NLR emitted spectrum. We discard the possibility of higher densities for the IB clouds since the derived value using the \[S II\] $`\lambda `$6717/$`\lambda `$6731 ratio in the galaxies are near $`10^3`$ cm<sup>-3</sup>.
Although a finer tuning of the parameters might improve the agreement between the observations and the model predictions for the objects studied here, this is unwarranted considering the simplicity of the approach taken. What is important to note is that the same initial set of physical conditions (albeit different SED) derived independently for normal Seyfert 1s are also valid to reproduce to a first approximation the emission line ratios observed in NLS1s.
## 5 New Emission Line Ratios for the NLR of NLS1s
As was discussed in the Introduction and Section 3.2, one of the main purposes of this paper is to review the current assumption that the NLR contributes H$`\beta `$ with 0.1$`\times `$ the flux of \[O III\] $`\lambda `$5007, implying a universal constant \[O III\] $`\lambda `$5007/H$`\beta `$ ratio of 10 (Osterbrock & Pogge 1985, Leighly 1999). This assumption is usually made on account of the difficulty in deblending the narrow and broad components of the permitted lines in these objects and is based on trends observed in narrow emission lines of Seyfert 2s and intermediate Seyfert 1 galaxies. According to the results obtained in Section 3 the NLR of NLS1s contributes 50% of the total H$`\beta `$ emission. This has important implications for the intrinsic ratios of the narrow lines. For instance, \[O III\] $`\lambda `$5007/H$`\beta `$ (narrow) now falls in the interval 1–5 since narrow H$`\beta `$ sees its flux increased by up to ten times its previously assumed value.
If we had (erroneously) assumed that in NLS1s the \[O III\]$`\lambda `$5007/H$`\beta `$ (narrow) ratio took the canonical value of 10 of normal Seyferts, we would have inferred a spectral index $`\alpha 1.4`$ to describe the EUV-to-soft-X-ray SED since such hardness is favored when atempting to reproduce a high ratio. The ROSAT observations, however, show that the soft X-ray spectra from NLS1s is systematically steeper ($`\alpha 2`$ or smaller) than those from Seyfert 1 galaxies with broad optical lines. In addition, Table 2 shows that photoionization models with $`\alpha 2.2`$ always provide \[O III\]$`\lambda `$5007/H$`\beta `$ ratios smaller than 7 under the initial set of conditions assumed. This result is therefore fully consistent with our findings of a much lower ratio in NLSIs than in normal Seyferts. Beside, much lower NLR densities ($`n_B10^5`$ cm<sup>-3</sup>) would have been required to get \[O III\] $`\lambda `$5007/H$`\beta `$ around 10 using SED 2. Such low densities would be discrepant with the higher ones determined in Paper I and from other authors (Osterbrock & Pogge 1985).
One possibility is that the continuum in the 13.6–200 eV region that photoionizes the NLR of NLS1s is similar to that of NGC 5548 (SED 1), but that it is modified before reaching the NLR gas by intervening material located between the BLR and NLR. Kraemer et al. (1999) explored, for instance, the effects of UV absorbing material on the shape of the continuum radiation emitted from the AGN, and on the relative strengths of the ensuing emission lines formed in the NLR of Seyfert 1 galaxies exposed to this emerging continuum distribution. Their results indicated that a low ionization UV absorber with a large covering factor can indeed modify the intrinsic EUV continuum of the central source and produce significative variations in the NLR emission spectrum of the AGN compared to that produced by the unattenuated continuum. Nonetheless, amongst the various types of warm absorbers tested, no model could reproduce a low \[O III\] $`\lambda `$5007/H$`\beta `$ ratio. It should be mentioned that they assumed the NLR gas to be uniformly radiation-bounded unlike the models presented here which consist of a combination of MB and IB components.
We consider worthy to test whether the inclusion of a low-ionization WA with similar characteristics to that of Kraemer et al. (1999) and exposed to SED 1 might explain the strongest (narrow) line ratios observed in NLS1s. In this picture, differences in \[O III\] $`\lambda `$5007/H$`\beta `$ for example would be caused by the effects of the intervening WA rather than by changes in steepness of the intrinsic EUV – soft X-ray continuum, as was suggested in the preceding section. Results are presented in Fig. 7 where the short-dashed line shows the predicted line ratios for a NLR with the same physical parameters as employed in modeling NCG 5548 but photoionized by the continuum leaking from a WA with $`U_{WA}`$=0.01, $`n_H=1\times 10^7`$ cm<sup>-3</sup>, solar abundance and thickness N$`{}_{\mathrm{H}}{}^{}=10^{20}`$ cm<sup>-2</sup>. The intrinsic SED incident on the WA was SED 1.
Comparison of these results with those of the unabsorbed model (solid line) allows us to say that the variations of the emission line ratios introduced by an intervening low-ionization WA are relatively minor and clearly insufficient for explaining the differences in line ratios observed between NLS1s and normal Sy1s. On the other hand, models with $`U_{WA}<10^{2.5}`$ are found to alter the incident EUV-to-soft-X-ray distribution in such a way that few photons are left to ionize the NLR as in Kraemer et al. (1999). It can be seen that the \[O III\] $`\lambda `$5007/H$`\beta `$ ratio does not change significantly relative to the unabsorbed model. In either cases, low ionization lines such as \[O I\] $`\lambda `$6300 (Figure 7c) see their ratios increase relative to H$`\beta `$. The same trend is observed for \[N II\] $`\lambda `$6584 and \[S II\] $`\lambda `$$`\lambda `$6717,6731. We recall from the previous section that these lines were already overpredicted. The lower the ionization degree of the WA, the larger this effect. This is clearly understood if we consider that the net effect of a low-ionization WA is to modify the EUV distribution while the hard EUV portion is left unaltered. Since low ionization lines are produced by the harder EUV radiation, they are now augmented due to the reduction of ionizing photons which produce the Balmer and He II $`\lambda `$4686 lines.
We therefore conclude that a low-ionization WA, if it exists in NLS1s, cannot explain the low \[O III\] $`\lambda `$5007/H$`\beta `$ ratio observed in these objects.
## 6 Conclusions
We have analyzed long-slit spectral data of a sample composed of seven NLS1 galaxies. A decomposition of the H$`\alpha `$ and H$`\beta `$ emission line profiles into Gaussian components allowed us to separate the flux contribution of the NLR from the total flux of the line. Our results show that, on average, 50% of the total H$`\beta `$ flux is due to emission from the NLR. Using the \[O III\] $`\lambda `$5007 line profile as a template for the narrow lines in order to subtract this contribution in the permitted lines give very similar results to those obtain throught the Gaussian decomposition. This confirms the presence of a broad component in the permitted lines.
The FWHM of the broad components of H$`\alpha `$ and H$`\beta `$ in the NLS1s studied here seems to be rather uniform within the same galaxy after comparing H$`\alpha `$ and H$`\beta `$ and throughout the sample (2250 km s<sup>-1</sup> and 2560 km s<sup>-1</sup> for H$`\alpha `$ and H$`\beta `$, respectively). No evidence of Lorentzian profiles was observed neither in the narrow nor the broad lines. The narrow components of H$`\alpha `$ and H$`\beta `$ present FWHM comparable to those of the forbidden lines which are typical of the NLR in any Seyfert.
The resulting \[O III\] $`\lambda `$5007/H$`\beta `$ ratios fall in the interval 1–5, significantly lower than the value currently assumed ($``$10). This entails that the emission line ratios from the NLR are different in NLS1s from those observed in normal and intermediate Sy1 galaxies.
We test photoionization models that consider a NLR composed of a combination of matter-bounded clouds and ionization-bounded clouds. The former, with typical densities $`10^6`$ cm<sup>-3</sup> and photoionized by the intrinsic continuum from the central source, are responsible for the emission of most of the \[O III\] and high ionization lines. This component should be located in the inner regions of the NLR. The latter, located farther out than the MB clouds and characterized by a lower density ($`n_e10^3`$ cm<sup>-3</sup>), are photoionized by the continuum filtered from the MB clouds and emits most of the low ionization lines. Assuming similar physical parameters in the NLR of NLS1s and normal Sy1s, we show that the observed differences in emission line ratios between these two groups of galaxies can be explained in terms of differences in the form of the input ionizing spectra. NLS1s ratios are better reproduced with a steep power-law continuum, with spectral index $`\alpha <2`$ while flatter spectral indices ($`\alpha 1.5`$) match the observed line ratios in normal Sy1s. This scenario reproduces with very good agreement the line ratios of NLS1s. It is furthermore consistent with ROSAT observations of NLS1s, which show that these objects are characterized by steeper power-law indices than those of Sy1 galaxies with broad optical lines. Our modeling therefore support the view that the NLR is directly photoionized by the unaltered SED distribution emitted by the central engine.
ARA gratefully acknowledges the staff of the CASLEO Observatory for instrumental and observing assistance. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, Californian Institute of Technology, under contract with the National Aeronautics and Space Administration. The work of Luc Binette was supported by the CONACyT grant 27546-E and the work of ARA and MGP was supported by the PRONEX/FINEP grant 76.97.1003.00
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# The approximation numbers of Hardy–type operators on trees.
## 1 Introduction.
In ,and results were established for the Hardy operator
$$Tf(x)=v(x)_0^xf(t)u(t)𝑑t$$
(1. 1)
as a map from $`L^p(0,\mathrm{})`$ to $`L^p(0,\mathrm{})`$, for $`1p\mathrm{}`$. When $`p(1,\mathrm{})`$, it was proved in that under appropriate conditions on $`u`$ and $`v`$ the approximation numbers $`a_n(T)`$ of $`T`$ satisfy
$$\underset{n\mathrm{}}{lim}na_n(T)=\frac{1}{\pi }_0^{\mathrm{}}|u(t)v(t)|𝑑t$$
(1. 2)
1991 Mathematics Subject Classification : 47G10, 47B10
when p=2, and when $`1<p<\mathrm{}`$
$`{\displaystyle \frac{\alpha _p}{4}}{\displaystyle _0^{\mathrm{}}}|u(t)v(t)|𝑑t`$ $``$ $`\underset{n\mathrm{}}{lim\; inf}na_n(T)`$ (1. 3)
$``$ $`\underset{n\mathrm{}}{lim\; sup}na_n(T)\alpha _p{\displaystyle _0^{\mathrm{}}}|u(t)v(t)|𝑑t`$
for a specified constant $`\alpha _p`$ depending on $`p`$. For the cases $`p=\mathrm{}`$ and $`p=1`$ similar estimates were derived in but with $`v_s(t)=lim_{\epsilon 0}v_{\mathrm{},(t\epsilon ,t+\epsilon )}`$ instead of $`v(t)`$ when $`p=\mathrm{}`$ and $`u_s(t)`$ instead of $`u(t)`$ in the case $`p=1`$. Furthermore, in and two-sided estimates are given for the $`l^\alpha `$ and weak-$`l^\alpha `$ norms of the sequence of approximation numbers in the case when the Hardy operator is compact.
A special case of the main result in this paper is that the counterpart of (1.2) in $`L^p(0,\mathrm{})`$, namely
$$\underset{n\mathrm{}}{lim}na_n(T)=\alpha _p_0^{\mathrm{}}|u(t)v(t)|𝑑t,$$
(1. 4)
holds for all $`p(1,\mathrm{})`$; the general result is that the analogue of (1.4) when the interval $`(0,\mathrm{})`$ is replaced by a tree $`\mathrm{\Gamma }`$ is true. Such Hardy operators on trees have already been investigated in where it was shown that they occur naturally in spectral problems defined on domains with irregular boundaries. Necessary and sufficient criteria for the boundedness of Hardy operators between various Lebesgue spaces on $`\mathrm{\Gamma }`$ are established in , but the complex nature of the problem is such that the neat abstract result is difficult to apply even for the most elementary of trees. It is therefore to be expected that the problems of compactness and estimating the approximation numbers are likely to be much more complicated than in the interval case. This, confirmed in this paper, is what makes it so surprising that the analogue of (1.4) for a tree is established here when $`p2`$ before it was known for an interval. Estimates for $`l^q`$ and weak-$`l^q`$ norms of the approximation numbers of $`T`$ are also obtained.
In the case $`p=2`$ of the problem subsequently studied in and for general $`p[1,\mathrm{}]`$ was considered and (1.2) proved, using Hilbert space methods which do not extend to general values of $`p`$. The same problem on a tree $`\mathrm{\Gamma }`$ is the subject of where an intensive study is made of problems on trees which are closely related to those here, but in the case $`p=2`$ only, and using methods which are very different from those in this paper. The conditions imposed to ensure the validity of the analogue of (1.4) for a tree in are similar to those here, but a comparison seems difficult in general (see Remark 6.12). The main difference is that in they relate to arbitrary partitions of $`\mathrm{\Gamma }`$ into intervals, whereas our partitions are into connected subsets specifically determined by functions which have a fundamental role in the analysis.
## 2 Preliminaries.
In this section we recall the definition of a tree $`\mathrm{\Gamma }`$, introduce a Hardy–type operator on the tree and quote from the criterion for the boundedness of the operator as a map from $`L^p(\mathrm{\Gamma })`$ into $`L^p(\mathrm{\Gamma })`$.
A tree $`\mathrm{\Gamma }`$ is a connected graph without loops or cycles, where the edges are non-degenerate closed line segments whose end–points are the vertices. Each vertex of $`\mathrm{\Gamma }`$ is of finite degree, i.e. only a finite number of edges emanate from each vertex. For every $`x,y\mathrm{\Gamma }`$ there is a unique polygonal path in $`\mathrm{\Gamma }`$ which joins $`x`$ and $`y`$. The distance between $`x`$ and $`y`$ is defined to be the length of this polygonal path and in this way $`\mathrm{\Gamma }`$ is endowed with a metric topology.
###### Lemma 2.1
Let $`\tau (\mathrm{\Gamma })`$ be the metric topology on $`\mathrm{\Gamma }`$. Then
a set $`A\mathrm{\Gamma }`$ is compact if and only if it is closed and meets only a finite number of edges;
$`\tau (\mathrm{\Gamma })`$ is locally compact;
$`\mathrm{\Gamma }`$ is the union of a countable number of edges. Thus if $`\mathrm{\Gamma }`$ is endowed with the natural 1-dimensional Lebesgue measure it is a $`\sigma `$-finite measure space.
Proof. See . $`\mathrm{}`$
Let $`x,y\mathrm{\Gamma }`$ and denote by $`(x,y)`$ the unique path joining $`x,y`$ in $`\mathrm{\Gamma }`$. For $`a\mathrm{\Gamma }`$ we define $`t_ax`$ (or $`x_at`$) to mean that $`x`$ lies on the path $`(a,t)`$. We write $`x_at`$ for $`x_at`$ and $`xt`$. This is a partial ordering on $`\mathrm{\Gamma }`$ and the ordered graph so formed is referred to us a tree rooted at $`a`$ and denoted by $`\mathrm{\Gamma }_a`$. If $`a`$ is not a vertex we make it one by replacing the edge on which it lies by two edges. In this way $`\mathrm{\Gamma }_a`$ is the unique finite union of subtrees $`\mathrm{\Gamma }_{a,i}`$ which intersect only at $`a`$.
Note that if $`x(a,b)`$ then $`x_ay`$ if and only if $`x_by`$.
We shall use the following notation. For a subtree $`K`$ of $`\mathrm{\Gamma }`$, $`V(K)`$, $`E(K)`$ will denote respectively the sets of vertices, edges of $`K`$ and $`K`$ will denote the set of boundary points of $`K`$ in $`\mathrm{\Gamma }`$. The notation $`K\mathrm{\Gamma }`$ will be used to mean that the closure of $`K`$ is a compact subset of $`\mathrm{\Gamma }`$; note that, from Lemma 2.1 (i) this implies that $`K`$ meets only a finite number of edges of $`\mathrm{\Gamma }`$. The characteristic function of a set $`E`$ will be denoted by $`\chi _E`$. The integral is interpreted in the following sense :
$$_Eg=\underset{e}{}_{eE}g$$
where
$$_{eE}g=_c^dg(x)\chi _E(x)𝑑x,$$
the integral $`_c^d`$ being over the set of points lying in the path $`(c,d)`$. For a measurable subset $`K`$ of $`\mathrm{\Gamma }`$ we define the norm
$$f_{p,K}=\left(_K|f|^p\right)^{1/p}$$
on $`L^p(K)`$. The $`L^p(\mathrm{\Gamma })`$ norm will be denoted by $`_p`$ if there is no chance of confusion. Also, if the value of $`p`$ is clear from the context, we shall write $`_K,`$ for the $`L^p`$ norms on $`K,\mathrm{\Gamma }`$ respectively. If $`A`$ is a bounded map between normed spaces $`X,Y`$ we denote its norm by
$$A|XY$$
This will be simplified to $`A`$ if the spaces $`X,Y`$ are unambiguous.
A connected subset of $`\mathrm{\Gamma }`$ is a subtree if we add its boundary points to the set of vertices of $`\mathrm{\Gamma }`$, and hence form new edges from existing ones. Hereafter we shall always adopt this convention when we refer to subtrees.
###### Definition 2.2
Let $`K`$ be a subtree of $`\mathrm{\Gamma }`$ containing $`a`$. A point $`tK`$ is said to be maximal if every $`x_at`$ lies in $`\mathrm{\Gamma }K`$. We denote by $`𝐈_a(\mathrm{\Gamma })`$ (or simply $`𝐈_a`$) the set of subtrees $`K`$ of $`\mathrm{\Gamma }`$ containing a whose boundary points are all maximal.
We assume throughout, unless mentioned otherwise, that $`u,v`$ satisfy the following conditions:
$$uL^p^{}(K),vL^p(\mathrm{\Gamma }),\text{ for every }K\mathrm{\Gamma }.$$
(2. 5)
We may assume, without loss of generality, that $`u,v0`$. This is because multiplication by $`sgnu`$ and $`sgnv`$ are isometries on $`L^p(\mathrm{\Gamma })`$; recall that $`sgnu=u/|u|`$ when $`u0`$ and $`1`$ otherwise.
###### Definition 2.3
Let $`\mathrm{\Gamma }`$ be a tree, $`1p\mathrm{}`$, and let $`u`$ and $`v`$ be measurable functions on $`\mathrm{\Gamma }`$ which satisfy (2. 5). For $`x\mathrm{\Gamma }`$ and $`fL^p(\mathrm{\Gamma })`$ we define the Hardy operator by
$$T_af(x):=v(x)_a^xf(t)u(t)𝑑t,a\mathrm{\Gamma }.$$
(2. 6)
In the following necessary and sufficient condition for the boundedness of $`T_a`$ was obtained.
###### Theorem 2.4
Let $`1p\mathrm{}`$, $`a\mathrm{\Gamma }`$, and suppose $`u`$ and $`v`$ satisfy (2. 5). For $`K𝐈_a`$ define
$$\alpha _K:=inf\{f_p:_a^t|f||u|=1\text{ for all }tK\}.$$
(2. 7)
Then $`T_a`$ is bounded from $`L^p(\mathrm{\Gamma })`$ into $`L^p(\mathrm{\Gamma })`$ if and only if
$$A:=\underset{K𝐈_a}{sup}\frac{v\chi _{\mathrm{\Gamma }K}_p}{\alpha _K}<\mathrm{}.$$
(2. 8)
Moreover, $`AT_a4A`$.
## 3 Bounds for the approximation numbers
We recall that, given any $`m𝐍`$, the $`m`$–th approximation number of a bounded operator $`T:L^p(\mathrm{\Gamma })L^p(\mathrm{\Gamma })`$, $`a_m(T)`$, is defined by
$$a_m(T):=infTF|L^p(\mathrm{\Gamma })L^p(\mathrm{\Gamma }),$$
where the infimum is taken over all bounded linear maps $`F:L^p(\mathrm{\Gamma })L^p(\mathrm{\Gamma })`$ with rank less than $`m`$.
A measure of non-compactness of $`T`$ is given by
$$\beta (T):=infTP|L^p(\mathrm{\Gamma })L^p(\mathrm{\Gamma }),$$
where the infimum is taken over all compact linear maps $`P:L^p(\mathrm{\Gamma })L^p(\mathrm{\Gamma })`$. Since $`L^p(\mathrm{\Gamma })`$ has the approximation property for $`1p\mathrm{}`$, $`T`$ is compact if and only if $`a_m(T)0`$ as $`m\mathrm{}`$, and $`\beta (T)=lim_n\mathrm{}a_n(T)`$.
###### Definition 3.1
Let $`K`$ be a subtree of $`\mathrm{\Gamma }`$ and $`a\mathrm{\Gamma }`$. We define:
$$A(K)A(K,u,v):=\{\begin{array}{cc}sup_{fL^p(K),f0}inf_{\alpha 𝐂}\frac{T_{a,K}f\alpha v_{p,K}}{f_{p,K}}\hfill & \text{ if }\mu (K)>0,\hfill \\ 0\hfill & \text{ if }\mu (K)=0,\hfill \end{array}$$
where
$$T_{a,K}f(x):=v(x)\chi _K(x)_a^xu(t)f(t)\chi _K(t)𝑑t,$$
and
$$\mu (K):=\{\begin{array}{cc}_K|v(t)|^p𝑑t\hfill & 1p<\mathrm{},\hfill \\ _K|v(t)|𝑑t\hfill & p=\mathrm{}.\hfill \end{array}$$
###### Lemma 3.2
The number $`A(K,u,v)`$ in Definition 3.1 is independent of $`a\mathrm{\Gamma }`$.
Proof. Denote by $`S`$ the canonical map of $`L^p(K)`$ into its quotient by the space of scalar multiples of $`v`$. Then $`A(K)=ST_{a,K}|L^p(K)L^p(K)`$. For $`b\mathrm{\Gamma }`$ we have $`T_{b,K}f=v\chi _K_b^afu\chi _K𝑑t+T_{a,K}Uf`$, where $`Uf(t)=f(t)`$ if $`t`$ lies on the path $`(a,b)`$ and $`f(t)`$ otherwise. Clearly $`U`$ is a linear isometry of $`L^p(K)`$ onto itself and $`ST_{b,K}=ST_{a,K}U`$. $`\mathrm{}`$
###### Corollary 3.3
For all subtrees $`K\mathrm{\Gamma }`$
$$A(K)\underset{a\mathrm{\Gamma }}{inf}T_{a,K}|L^p(K)L^p(K).$$
Note that if $`\mathrm{\Lambda }`$ is a subtree of $`\mathrm{\Gamma }`$, $`a\mathrm{\Gamma }`$ and $`b`$ the nearest point of $`\mathrm{\Lambda }`$ to $`a`$ then $`T_{b,\mathrm{\Lambda }}=T_{a,\mathrm{\Lambda }}`$ and $`T_{b,\mathrm{\Lambda }}:=T_{b,\mathrm{\Lambda }}|L^p(\mathrm{\Lambda })L^p(\mathrm{\Lambda })T_a=:T_a|L^p(\mathrm{\Gamma })L^p(\mathrm{\Gamma })`$. Moreover if $`\mathrm{\Lambda }^{}\mathrm{\Lambda }`$, with $`c`$ the nearest point of $`\mathrm{\Lambda }^{}`$ to $`a`$ and $`(b,c)`$ a subinterval of an edge of $`\mathrm{\Lambda }`$ then $`T_{b,\mathrm{\Lambda }}f=T_{b,\mathrm{\Lambda }}(f\chi _{(b,c)})+T_{c,\mathrm{\Lambda }^{}}(f\chi _\mathrm{\Lambda }^{}),`$ whence $`0T_{b,\mathrm{\Lambda }}T_{c,\mathrm{\Lambda }^{}}u_{p^{},(b,c)}v_{p,(b,c)}`$. This remark yields
###### Lemma 3.4
For $`1p\mathrm{},T_{x,K}|L^p(K)L^p(K)`$ is continuous in $`x`$.
Let $`x\mathrm{\Gamma }`$. Denote by $`\mathrm{\Gamma }_{x,i}`$ $`i=1,\mathrm{},n_x`$ the non-overlapping subtrees of $`\mathrm{\Gamma }`$ which are the closures of the connected components of $`\mathrm{\Gamma }\{x\}`$, and set $`T_{x,i}T_{x,\mathrm{\Gamma }_{x,i}}`$ and $`T_{x,i}T_{x,i}|L^p(\mathrm{\Gamma }_{x,i})L^p(\mathrm{\Gamma }_{x,i})`$ . We suppose that the numbering is done in the order of descending norms of the $`T_{x,i}:L^p(\mathrm{\Gamma }_{x,i})L^p(\mathrm{\Gamma }_{x,i})`$. Clearly $`T_x=\mathrm{max}_{i=1,\mathrm{},n_x}T_{x,i}`$.
Call a point $`x\mathrm{\Gamma }`$ simple if there is just one $`T_{x,i}`$ with maximal norm, so that $`T_{x,1}>T_{x,2}.`$ If $`a`$ is a simple point and $`(a,y)`$ the first edge of $`\mathrm{\Gamma }_{a,1}`$ then by continuity either there is a point $`z`$ of $`(a,y)`$ which is not simple or $`a\mathrm{\Gamma }_{y,1}`$. If the latter, continue the path beginning with $`(a,y)`$ along the initial edge of $`\mathrm{\Gamma }_{y,1}`$. By induction thus define a path $`l`$ in $`\mathrm{\Gamma }`$ satisfying one of the following:
$`l`$ is finite and its end $`b`$ is not simple;
$`l`$ is finite, its end $`b`$ is simple and $`\{x:x_ab\}=\mathrm{}`$;
$`l`$ is infinite.
Now (ii) is impossible since $`lim_{xb}T_{x,1}=0`$, and $`T_{x,1}A(\mathrm{\Gamma })`$. Also (iii) implies $`T`$ is not compact. For if $`x`$ is in $`l`$, $`T_{x,1}A(\mathrm{\Gamma })`$ and hence there is a compact subset $`K`$ of $`\mathrm{\Gamma }_{x,1}`$ and a function $`f`$ supported in $`K`$ with $`f1`$ and $`T_af_K\frac{1}{2}A(\mathrm{\Gamma })`$. It follows that there is a sequence of disjoint compact sets $`K_n`$ and functions $`f_n`$ with the same property. Then, if $`m>n,T_a(f_nf_m)_\mathrm{\Gamma }T_af_n_{K_n}\frac{1}{2}A(\mathrm{\Gamma })`$. Thus, if $`T_a`$ is compact, (i) holds. Moreover, $`T_b=\mathrm{min}_{x\mathrm{\Gamma }}T_x.`$ For if $`xb`$ then $`x`$ one of $`\mathrm{\Gamma }_{b,1}`$, $`\mathrm{\Gamma }_{b,2}`$, say $`\mathrm{\Gamma }_{b,2}.`$ Then if $`\mathrm{\Gamma }_{x,j}`$ is the subtree containing $`b`$, $`T_{x,j}T_{b,2}=T_{b,1}.`$ From this we have the following result which will be an important tool for determining a lower bound for $`A(K)`$ once Theorem 3.8 below is available.
###### Lemma 3.5
Suppose $`T_a`$ is compact and that there exist $`ij`$ such that $`T_{a,i},T_{a,j}T_a`$. Then
$$\mathrm{min}\{T_{a,i},T_{a,j}\}\underset{x\mathrm{\Gamma }}{\mathrm{min}}T_x.$$
Proof. The result is a consequence of the discussion preceding the lemma if $`a`$ is not simple. If $`a`$ is simple, then, with $`b`$ the non-simple end-point of the path $`l`$ in (i) above, and $`\mathrm{min}\{T_{a,i},T_{a,j}\}=T_{a,i}`$, say, we have $`T_{a,i}<T_a`$ and $`\mathrm{\Gamma }_{a,i}`$ is a subtree of some tree $`\mathrm{\Gamma }_{b,k}`$. Thus
$$T_{a,i}T_{b,k}T_b=\underset{x\mathrm{\Gamma }}{\mathrm{min}}T_x.$$
$`\mathrm{}`$
In the next two lemmas $`_{p,\mu }`$ denotes the norm in $`L^p(\mathrm{\Gamma },d\mu )`$, where $`d\mu (t)=|v(t)|^pdt`$.
###### Lemma 3.6
If $`1<p\mathrm{}`$ there is a unique scalar $`c_f`$ such that $`fc_fe_{p,\mu }=inf_{c𝐂}fce_{p,\mu }`$ for $`e0,eL^p(\mathrm{\Gamma },d\mu )`$.
Proof. Since $`fce_{p,\mu }`$ is continuous in $`c`$ and tends to $`\mathrm{}`$ as $`c\mathrm{}`$, the existence of $`c_f`$ is guaranteed by the local compactness of $`𝐂`$. For $`1<p<\mathrm{}`$ the uniqueness follows from the uniform convexity of $`L^p(\mathrm{\Gamma },d\mu )`$. Let $`p=\mathrm{}`$, and suppose that there are two values of $`c_f,c_1c_2`$. This yields the contradiction $`f(1/2)(c_1+c_2)_{p,\mu }<fc_1_{p,\mu }`$.$`\mathrm{}`$
###### Lemma 3.7
The map $`fc_f:L^p(\mathrm{\Gamma },d\mu )𝐂`$ is continuous for $`1<p\mathrm{}`$.
Proof. Suppose that $`c_{g_n}c`$ as $`g_nf`$. Then
$$g_nc_f_{p,\mu }g_nc_{g_n}_{p,\mu }$$
and so
$$fc_f_{p,\mu }fc_{p,\mu }$$
which gives $`c=c_f`$ $`\mathrm{}`$
###### Theorem 3.8
Let $`1<p\mathrm{}`$. If $`T_a`$ is compact $`A(\mathrm{\Gamma })=\mathrm{min}_{x\mathrm{\Gamma }}T_x|L^p(\mathrm{\Gamma })L^p(\mathrm{\Gamma })`$.
Proof. There is a non-simple point $`b`$ at which $`T_x`$ attains its minimum. If $`\alpha <T_b`$ there exist $`f_i,i=1,2`$, supported in $`\mathrm{\Gamma }_{b,i}`$ with $`f_i=1,T_bf_i>\alpha `$, and $`f_1`$ positive, $`f_2`$ negative. Clearly the same is true of the corresponding values of $`c_f`$, say $`c_1,c_2`$. Then by continuity there is a $`\lambda [0,1]`$ such that $`c_g=0`$ for $`g=\lambda f_1+(1\lambda )f_2`$, and $`T_bg^p=\lambda ^pT_bf_1^p+(1\lambda )^pT_bf_2^p>\alpha ^pg^p`$. Then, by Lemma 3.6,
$$A(\mathrm{\Gamma })\underset{c}{inf}(T_bcv)g/g=T_bg/g>\alpha .$$
Since $`\alpha <T_b`$ is arbitrary, $`A(\mathrm{\Gamma })T_b`$ and the result follows from Corollary 3.3.$`\mathrm{}`$
The next lemma establishes an important geometrical property of a tree which is an essential ingredient of the subsequent analysis. First we make some observations.
Suppose $`w`$ is a non–negative function defined on the set of all closed subtrees of a tree $`\mathrm{\Gamma }`$, satisfying
$$XYw(X)w(Y).$$
(3. 9)
Define
* $$N_\epsilon (\mathrm{\Gamma })=\underset{𝒮_\epsilon (\mathrm{\Gamma })}{\mathrm{min}}\mathrm{\#}$$
where $`𝒮_\epsilon :=\{;\text{ is a set of non–overlapping closed subtrees of }\mathrm{\Gamma }\text{ such that }`$ $`i)_XX=\mathrm{\Gamma },ii)Xw(X)\epsilon \}`$;
* $$M_\epsilon (\mathrm{\Gamma })=\underset{𝒢_\epsilon (\mathrm{\Gamma })}{\mathrm{max}}\mathrm{\#}𝒢$$
where $`_\epsilon :=\{𝒢;𝒢\text{ is a set of non–overlapping closed subtrees of }\mathrm{\Gamma }\text{ such that }`$ i) $`_{X𝒢}X=\mathrm{\Gamma }`$ ii) $`\mathrm{\#}\{X;X𝒢,w(X)\epsilon \}1\}`$
Two non–overlapping closed subtrees of $`\mathrm{\Gamma }`$ can have at most one point in common, for otherwise $`\mathrm{\Gamma }`$ would contain a cycle. A chain $`𝒞`$ of closed subtrees is a sequence $`X_1,\mathrm{},X_l`$ of closed subtrees such that $`X_iX_{i+1}=\{x_i\}`$ $`(i=1,\mathrm{},l1)`$ where the $`x_i`$ are distinct. The length of $`𝒞`$ is $`l`$.
There is a set $`𝒮_\epsilon (\mathrm{\Gamma })`$ with $`\mathrm{\#}=N_\epsilon (\mathrm{\Gamma })`$ (possibly $`\mathrm{}`$). Let $`𝒞`$ be a chain of elements of $``$ of maximal length $`l`$. Then we have the following:
If $`l=1`$ then $`\mathrm{\#}=1`$ and so $`w(\mathrm{\Gamma })\epsilon `$.
If $`l=2`$ define $`Y=X_1X_2`$. Then if $`\mathrm{\Gamma }Y`$, $`\mathrm{\Gamma }Y^o`$ is a closed subtree and $`N_\epsilon (\mathrm{\Gamma }Y^o)=N_\epsilon (\mathrm{\Gamma })2`$. Moreover $`Y`$ is a closed subtree of $`\mathrm{\Gamma }`$ and $`w(Y)>\epsilon `$. For since $`l`$ is maximal $`x_1`$ lies in every member of $``$ so that $`\mathrm{\Gamma }Y^o`$ is a closed subtree. Also $`\mathrm{\Gamma }Y^o=_X^{}X`$, where $`^{}=\{X_1,X_2\}`$, which implies that $`N_\epsilon (\mathrm{\Gamma }Y^o)N_\epsilon (\mathrm{\Gamma })2.`$ If $`N_\epsilon (\mathrm{\Gamma }Y^o)<N_\epsilon (\mathrm{\Gamma })2`$, there exists $`^{\prime \prime }`$, a suitable covering of $`\mathrm{\Gamma }Y^o`$ with $`\mathrm{\#}^{\prime \prime }=N_\epsilon (\mathrm{\Gamma }Y^o)`$, and then $`^{\prime \prime }\{X_1,X_2\}𝒮_\epsilon (\mathrm{\Gamma })`$ which contradicts the definition of $`N_\epsilon (\mathrm{\Gamma })`$. Finally $`w(Y)>\epsilon `$ for if not, on taking $`^{\prime \prime \prime }=^{}\{Y\}`$ we have a contradiction.
If $`l3,`$ suppose $`Z_1,\mathrm{},Z_k`$ are the sets in $``$ which meet $`X_{l1}`$ in a point different from $`x_{l2}`$. Then since $`l`$ is maximal $`X_1,X_2,\mathrm{},X_{l1},Z_i`$, is a chain of maximal length for $`i=1,2,\mathrm{},k`$. Either a) $`k=1`$ or b) $`k>1`$. If a), we have $`Z_i=X_l`$ and we take $`Y=X_lX_{l1}`$, so $`Y`$ is a closed subtree. Then $`\mathrm{\Gamma }Y^o`$ is a closed subtree of $`\mathrm{\Gamma }`$. Moreover $`N_\epsilon (\mathrm{\Gamma }Y^o)=N_\epsilon (\mathrm{\Gamma })2`$, $`w(Y)>\epsilon `$ by an argument similar to that of (i). If b), define $`a_i`$ by $`\{a_i\}=X_{ł1}Z_i`$ and, for $`ij,a_{ij}=a_ia_j=\mathrm{max}\{u:ua_i\text{ and }ua_j\}`$ in the ordering of $`\mathrm{\Gamma }`$ arising from taking $`x_1`$ as root. Then $`a_{ij}X_{l1}.`$ Define $`\varrho _i:=\mathrm{\#}\{a_{jk};a_{jk}a_i\}.`$ Without loss of generality we may suppose $`\varrho _1=\mathrm{max}_i\varrho _i`$ and that $`a_{jk}a_{12}`$ for all $`a_{jk}a_1`$. Define $`Y=\{x;(u)a_{12}ua_1\text{ or }a_{12}ua_2\text{ and }ux\}\{a_{12}\}`$. Then $`Y`$ is a closed subtree of $`\mathrm{\Gamma }`$ and $`YZ_1Z_2X_{l1}.`$ If $`\mathrm{\Gamma }^{}=\mathrm{\Gamma }Y^o`$, $`\mathrm{\Gamma }^{}`$ is a closed subtree of $`\mathrm{\Gamma }`$ and it follows by arguments similar to those in (ii) that, since $`𝒞`$ is a chain of maximal length in $``$ and $``$ is a covering of $`\mathrm{\Gamma }`$, $`w(Y)>\epsilon `$ and
$$N_\epsilon (\mathrm{\Gamma })3N_\epsilon (\mathrm{\Gamma }^{})N_\epsilon (\mathrm{\Gamma })2.$$
(3. 10)
###### Lemma 3.9
If $`N_\epsilon (\mathrm{\Gamma })<\mathrm{}`$ then $`M_\epsilon (\mathrm{\Gamma })\frac{1}{3}N_\epsilon (\mathrm{\Gamma })`$.
Proof. The proof is by induction on $`N_\epsilon (\mathrm{\Gamma })`$. If $`N_\epsilon (\mathrm{\Gamma })=1,w(\mathrm{\Gamma })\epsilon `$ and $`M_\epsilon (\mathrm{\Gamma })=1`$. If $`N_\epsilon (\mathrm{\Gamma })=2,l=2`$ and $`\mathrm{\Gamma }=X_1X_2`$ and $`M_\epsilon (\mathrm{\Gamma })1`$. If $`N_\epsilon (\mathrm{\Gamma })>2`$ the arguments of i) ii) iii) above and induction prove the result. For by i) $`l1`$ and by ii) or iii) $`N_\epsilon (\mathrm{\Gamma }Y^o)<N_\epsilon (\mathrm{\Gamma })`$ so we may suppose $`M_\epsilon (\mathrm{\Gamma }Y^o)1/3N_\epsilon (\mathrm{\Gamma }Y^o)`$ and then there exists $`𝒢^{}_\epsilon (\mathrm{\Gamma }Y^o)`$ with $`\mathrm{\#}𝒢^{}=M_\epsilon (\mathrm{\Gamma }Y^o)`$. But then if $`𝒢=𝒢^{}\{Y\}`$, $`𝒢_\epsilon (\mathrm{\Gamma })`$ since $`\omega (Y)>\epsilon `$ and $`M_\epsilon (\mathrm{\Gamma })\mathrm{\#}𝒢=M_\epsilon (\mathrm{\Gamma }Y^o)+1\frac{1}{3}N_\epsilon (\mathrm{\Gamma }Y^o)+1\frac{1}{3}N_\epsilon (\mathrm{\Gamma }).`$ $`\mathrm{}`$
###### Lemma 3.10
Let $`K`$ be a compact tree, let $`w`$ be a function satisfying (3. 9) and, for every $`(c,d)E(K)`$, suppose that $`w(c,.)`$ is a continuous function on $`(c,d)`$. Then there exists a set $`𝒢`$ of non-overlapping subtrees of $`K`$ with $`\omega (X)\epsilon `$ for $`X𝒢`$ and
$$\mathrm{\#}𝒢N_\epsilon (K)3\mathrm{\#}E(K).$$
(3. 11)
Proof. Let $`S_\epsilon (K)`$ and $`\mathrm{\#}=N_\epsilon (K)`$. For $`(c,d)E(K)`$ define $`𝒦:=\{X:X,|X(c,d)|>0\}`$ to be such that $`\mathrm{\#}𝒦>3`$. Set $`𝒦=\{X_1,\mathrm{},X_n\}`$, where $`n=\mathrm{\#}𝒦`$ and $`X_1X_2\mathrm{}X_n`$, where $`XY`$ means that $`xy`$ for all $`xX,yY`$. Then $`X_i=(a_i,b_i)(c,d)`$ for $`1<i<n`$. By the continuity of $`\omega `$ on $`(c,d)`$ and the minimality of $`\mathrm{\#}`$, there exists $`b_2^{}(b_2,b_3)`$ such that $`\omega (a_2,b_2^{})=\epsilon `$. It follows that there are non-overlapping sub-intervals $`X_2^{},\mathrm{},X_k^{}`$ of $`(c,d)`$, where $`k=n2`$, for which $`\omega (X_j^{})=\epsilon `$. The lemma follows from this.$`\mathrm{}`$
###### Definition 3.11
Let $`K`$ be a subtree $`\mathrm{\Gamma }`$. We denote by $`𝒫(K)`$ the set of partitions $`\{\mathrm{\Gamma }_i:i=1,\mathrm{},n\}`$ of $`K`$ (i.e $`_{i=1}^n\mathrm{\Gamma }_i=K`$) by subtrees $`\mathrm{\Gamma }_i`$ of $`K`$ such that $`|\mathrm{\Gamma }_i\mathrm{\Gamma }_j|=0`$ for $`ij`$. For a given $`\epsilon >0`$ we define
$`N(K,\epsilon )N(K,\epsilon ,u,v):=\mathrm{min}\{n:`$ $`\{\mathrm{\Gamma }_i:i=1,\mathrm{},n\}𝒫(K)`$
$`\text{and }A(\mathrm{\Gamma }_i,u,v)\epsilon \};`$
$`M(K,\epsilon )M(K,\epsilon ,u,v):=`$ $`\mathrm{max}\{m:\text{ non-overlapping subtrees }\mathrm{\Gamma }_iK,`$
$`i=1,\mathrm{},m,\text{ such that }A(\mathrm{\Gamma }_i,u,v)>\epsilon \}.`$
Note that, with $`\omega ()=A()`$, $`M_\epsilon (\mathrm{\Gamma })M(\mathrm{\Gamma },\epsilon )+1`$.
Hereafter we shall write $`T,T_K`$ for $`T_a,T_{a,K}`$ respectively, unless there is a possibility of confusion. The following lemma will yield a one-dimensional approximation to $`T`$ on a subtree $`K`$ of $`\mathrm{\Gamma }`$. We recall the notation
$$\mu (K)=_K𝑑\mu ,d\mu =v^pdx.$$
###### Lemma 3.12
Let $`K`$ be a subtree of $`\mathrm{\Gamma }`$ and $`vL^p(K)`$, $`1p\mathrm{}`$, with $`\mu (K)0`$. Then there exists $`w_K\{L^p(K;d\mu )\}^{}`$ (in case $`p=\mathrm{}`$, $`w_K`$ is also from $`\{L^{\mathrm{}}(K)\}^{}`$ ) such that:
$$w_K(1)=1,$$
$$w_K_{\{L^p(K,d\mu )\}^{}}=\frac{1}{v_{p,K}},(w_K_{\{L^{\mathrm{}}(K)\}^{}}=1\text{ for }p=\mathrm{})$$
and
$$\underset{c𝐂}{inf}(\phi c)v_{p,K}(\phi w_K(\phi ))v_{p,K}2\underset{c𝐂}{inf}(\phi c)v_{p,K}$$
(3. 12)
for all $`\phi L^p(K,d\mu )`$. In the case $`p=2`$
$$\underset{c𝐂}{inf}(\phi c)v_{L^p(K)}=(\phi w_K(\phi ))v_{L^p(K)}$$
(3. 13)
where
$$w_K(\phi )=\frac{1}{\mu (K)}_K\phi 𝑑\mu .$$
Proof. Define the linear function $`w`$ on the constants in $`L^p(K,d\mu )`$ with $`w(c.1)=c`$. Then $`w(1)=1`$ and $`w^p=1/\mu (K)`$ for $`1p<\mathrm{}`$ while $`w=1`$ when $`p=\mathrm{}`$. The existence of $`w_K`$ follows by the Hahn-Banach theorem, and (3.4) is immediate. The case $`p=2`$ is obvious. $`\mathrm{}`$
###### Remark 3.13
For $`1p<\mathrm{}`$ and $`\mathrm{\Gamma }=(0,\mathrm{})`$ the choice $`w_K(\phi )=\frac{1}{\mu (K)}_K\phi 𝑑\mu `$ is appropriate, and was that used in . In Lemma 2.4, when $`p=\mathrm{}`$ and $`\mathrm{\Gamma }=(a,b)`$, $`w_K`$ was defined as the limit along a filter base of subsets $`w_\beta `$ of the unit ball in $`\{L^{\mathrm{}}(K)\}^{}`$ defined by
$$w_\gamma (\phi ):=\frac{1}{|A_\gamma |}_{A_\gamma }\phi (x)𝑑x,\phi L^{\mathrm{}}(K).$$
###### Lemma 3.14
Let $`K`$ be a subtree of $`\mathrm{\Gamma },aK,1<p\mathrm{}`$, and suppose that $`T_{a,K}`$ is compact. Then
$$A(K)=T_{a,K}v\varpi |L^p(K)L^p(K),$$
where $`\varpi `$ is the bounded linear functional
$$\varpi (f)=_a^bfu$$
and $`bK`$ is such that $`A(K)=T_{b,K}|L^p(K)L^p(K)`$.
Proof. We know from Theorem 3.8 that there exists a $`bK`$ such that
$$A(K,u,v)=T_{b,K}|L^p(K)L^p(K)=T_{b,K}U|L^p(K)L^p(K),$$
where $`U`$ is the linear isometry defined in the proof of Lemma 3.2, and with respect to which
$$T_{b,K}f=v\chi _K_b^afu\chi _K+T_{a,K}Uf.$$
Thus, if we define
$$Pf(x)=v(x)\chi _K(x)_a^bUfu\chi _K,$$
we have
$$T_{a,K}P|L^p(K)L^p(K)=A(K,u,v),$$
and the lemma follows.$`\mathrm{}`$
###### Lemma 3.15
Let $`\epsilon >0`$, $`1p\mathrm{}`$ and let $`K`$ be a subtree of $`\mathrm{\Gamma }`$. If $`N:=N(K,\epsilon ,u,v)<\mathrm{}`$ then
$$a_{N+1}(T_K)\gamma _p\epsilon ,$$
where $`\gamma _p=2`$ when $`p2`$ and $`\gamma _2=1`$.
Proof. Let $`\{\mathrm{\Gamma }_i\}_1^N`$ be the partition of $`K`$ which defines $`N:=N(K,\epsilon ,v,u)`$ in Definition 3.11 and set $`Pf=_{i=1}^NP_if`$ where
$$P_if(x):=\chi _{\mathrm{\Gamma }_i}(x)v(x)\left[_a^{a_i}uf+w_i\left(_{a_i}^xuf\chi _{\mathrm{\Gamma }_i}\right)\right],$$
$`T_KT_{a,K},a_i`$ is the point in $`\mathrm{\Gamma }_i`$ nearest $`a`$ and $`w_iw_{\mathrm{\Gamma }_i}`$ is the linear function from Lemma 3.12. Then $`rankPN`$ and, on using Lemma 3.12, and setting $`T_iT_{\mathrm{\Gamma }_i}`$,
$`(T_KP)f_{p,K}^p`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}(T_KPf)_{p,\mathrm{\Gamma }_i}^p`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}T_iw_i\left({\displaystyle _{a_i}^{}}uf\chi _{\mathrm{\Gamma }_i}𝑑t\right)v()_{p,\mathrm{\Gamma }_i}^p`$
$``$ $`\gamma _p^p\left(\underset{i=1,\mathrm{},N}{\mathrm{max}}A(\mathrm{\Gamma }_i,u,v)\right)^pf_{p,K}^p`$
$``$ $`(\gamma _p\epsilon )^pf_{p,K}^p,`$
whence the lemma.$`\mathrm{}`$
###### Lemma 3.16
Let $`1<p\mathrm{},\epsilon >0,`$ and suppose that $`T_K`$ is compact. Then, with $`N=N(K,\epsilon ,u,v)<\mathrm{}`$,
$$a_{N+1}(T_K)\epsilon .$$
Proof. The same argument as in Lemma 3.15 applies but with
$$P_if(x)=\chi _{\mathrm{\Gamma }_i}(x)v(x)\left[_a^{a_i}uf+\varpi _i\left(_{a_i}^{b_i}uf\chi _{\mathrm{\Gamma }_i}\right)\right],$$
where $`\varpi _i,b_i`$ are as in Lemma 3.14 corresponding to $`K=\mathrm{\Gamma }_i`$.$`\mathrm{}`$
###### Lemma 3.17
Let $`\epsilon >0`$, $`1p\mathrm{}`$ and let $`K`$ be a subtree of $`\mathrm{\Gamma }`$. Let $`\{\mathrm{\Gamma }_i;i=1,\mathrm{},M\}`$ be a set of non-overlapping subtrees of $`K`$ such that $`A(\mathrm{\Gamma }_i)\epsilon `$ for all $`1iM`$. Then
$$a_M(T_K)\epsilon .$$
Proof. Let $`\lambda (0,1).`$ From the definition of $`A(\mathrm{\Gamma }_i)`$ we know that for $`i=1,\mathrm{},M`$ there exist $`\varphi _iL^p(\mathrm{\Gamma }),\varphi _i_p=1`$, with support in $`\mathrm{\Gamma }_i`$ such that
$$\underset{\alpha 𝐂}{inf}T_K\varphi _i\alpha v_{p,\mathrm{\Gamma }_i}>\lambda A(\mathrm{\Gamma }_i)\lambda \epsilon .$$
Let $`P:L^p(K)L^p(K)`$ be bounded with $`rank(P)<M`$. Then, there are constants $`\lambda _1,\mathrm{},\lambda _M`$ not all zero, such that
$$P\varphi =0,\varphi :=\underset{i=1}{\overset{M}{}}\lambda _i\varphi _i.$$
Then, noting that the following summation is over $`\lambda _i0`$, and denoting by $`a_i`$ the point of $`\mathrm{\Gamma }_i`$ nearest $`a`$, where $`T_K=T_{a,K}`$,
$`a_M^p\varphi _{p,K}^p`$ $``$ $`T_K\varphi P\varphi _{p,K}^p=T_K\varphi _{p,K}^p={\displaystyle \underset{i=1}{\overset{M}{}}}(T_K\varphi )\chi _{\mathrm{\Gamma }_i}_{p,K}^p`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{M}{}}}\chi _{\mathrm{\Gamma }_i}(x)v(x)({\displaystyle _{a_i}^x}\lambda _i\varphi _i(t)u(t)\chi _{\mathrm{\Gamma }_i}(t)dt`$
$`+{\displaystyle _a^{a_i}}\varphi (t)u(t)dt)^p_{p,K}`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{M}{}}}\left(T_K\varphi _i(x)+v(x){\displaystyle \frac{\eta _i}{\lambda _i}}\right)\lambda _i_{p,\mathrm{\Gamma }_i}^p,`$
$`\text{where }\eta _i:={\displaystyle _a^{a_i}}\varphi (t)u(t)𝑑t,`$
$``$ $`{\displaystyle \underset{i=1}{\overset{M}{}}}\underset{\alpha 𝐂}{inf}(T_K\varphi _i(x)v(x)\alpha )_{p,\mathrm{\Gamma }_i}^p|\lambda _i|^p`$
$``$ $`(\lambda \epsilon )^p{\displaystyle \underset{i=1}{\overset{M}{}}}\varphi _i_{p,\mathrm{\Gamma }_i}^p|\lambda _i|^p(\lambda \epsilon )^p\varphi _{p,K}^p.`$
$`\mathrm{}`$
###### Theorem 3.18
Let $`1p\mathrm{},\epsilon >0`$ and $`N=N(\mathrm{\Gamma },\epsilon ,u,v)<\mathrm{}`$(see Definition 3.11). Then
$$a_{N+1}(T)\gamma _p\epsilon $$
where $`\gamma _p=2`$ when $`p2`$ and $`\gamma _2=1`$, and
$$a_{[\frac{N}{3}]1}(T)>\epsilon .$$
The measure of non-compactness of $`T,\beta (T)`$, satisfies
$$\beta (T):=\underset{n\mathrm{}}{lim}a_n(T)inf\{\epsilon :N(\mathrm{\Gamma },\epsilon ,u,v)<\mathrm{}\},$$
where the symbol $``$ means that the quotient of the two sides lies between positive constants. Hence, $`T`$ is compact if and only if $`N(\mathrm{\Gamma },\epsilon ,u,v)<\mathrm{}`$ for all $`\epsilon >0`$. If $`T`$ is compact, $`\gamma _p=1`$ for $`1<p\mathrm{}`$.
Proof. The first inequality follows from Lemma 3.15. The second inequality follows from Lemmas 3.9 and 3.17 on putting $`w():=A()`$. The two inequalities imply the result about the measure of non-compactness, and hence the compactness, of $`T`$. The last statement is a consequence of Lemma 3.16.$`\mathrm{}`$
From Lemmas 3.10, 3.16, 3.17 and 3.18 with $`w(X):=A(X)`$ we derive
###### Lemma 3.19
Let $`1<p<\mathrm{}`$, $`\epsilon >0`$, and let $`K`$ be a compact tree. Then
$$a_{N+1}(T_K)\epsilon $$
and
$$a_{M1}(T_K)>\epsilon ,$$
where $`N=N(K,\epsilon ,u,v),and`$ $`MN3\mathrm{\#}E(K).`$
Since, by Lemma 3.2, $`A(\mathrm{\Gamma }_a)`$ is independent of $`a\mathrm{\Gamma }`$, the approximation numbers are independent of $`a\mathrm{\Gamma }`$. Note that the above proof of Lemma 3.19 requires $`A(c,.)`$ to be continuous (see Lemma 3.10). This is not true for $`p=1`$ or $`\mathrm{}`$.
## 4 Local properties of $`A`$
In this section we establish properties of the function $`A`$ in Definition 3.1 which will be needed for the local asymptotic results in the next section.
###### Lemma 4.1
Let $`u,v`$ be constants over a real interval $`I=(a_1,b_1)`$ and $`1p\mathrm{}`$. Then $`A(I,u,v)=|v||u||I|\alpha _p`$, where $`\alpha _p=A((0,1),1,1)`$.
Proof. We have
$`A(I,u,v)`$ $`=`$ $`\underset{f_{p,I}1}{sup}\underset{\alpha 𝐂}{inf}v\left({\displaystyle _{a_1}^x}uf(t)𝑑t\alpha \right)_{p,I}`$
$`=`$ $`|v||u|\underset{f_{p,I}1}{sup}\underset{\alpha 𝐂}{inf}{\displaystyle _{a_1}^x}f(t)𝑑t\alpha _{p,I}`$
$`=`$ $`|v||u||I|\underset{f_{p,(0,1)}1}{sup}\underset{\alpha 𝐂}{inf}{\displaystyle _0^x}f(t)𝑑t\alpha _{p,(0,1)}`$
$`=`$ $`|v||u||I|A((0,1),1,1).`$
$`\mathrm{}`$
Note that $`\alpha _1=\alpha _{\mathrm{}}=1/2`$ and $`\alpha _2=1/\pi `$.
###### Lemma 4.2
Let $`1p\mathrm{}`$, $`u_1,u_2L^p^{}(\mathrm{\Gamma })`$ and $`vL^p(\mathrm{\Gamma })`$. Then
$$|A(\mathrm{\Gamma },u_1,v)A(\mathrm{\Gamma },u_2,v)|v_{p,\mathrm{\Gamma }}u_1u_2_{p^{},\mathrm{\Gamma }}.$$
Proof. We have
$`A(\mathrm{\Gamma },u_1,v)`$ $`=`$ $`\underset{f_p=1}{sup}\underset{\alpha 𝐂}{inf}v(x)\left({\displaystyle _a^x}(u_1(t)u_2(t)+u_2(t))f(t)𝑑t\alpha \right)_p`$
$``$ $`\underset{f_p=1}{sup}\underset{\alpha 𝐂}{inf}[v(x){\displaystyle _a^x}(u_1(t)u_2(t))f(t)dt_{p,\mathrm{\Gamma }}`$
$`+`$ $`v(x)({\displaystyle _a^x}u_2(t)f(t)dt\alpha )_p]`$
$``$ $`\underset{f_p=1}{sup}\underset{\alpha 𝐂}{inf}\left[v_pu_1u_2_p^{}+v(x)\left({\displaystyle _{a_1}^x}u_2(t)f(t)𝑑t\alpha \right)_p\right]`$
$``$ $`v_pu_1u_2_p^{}+A(\mathrm{\Gamma },u_2,v).`$
The same holds with $`u_1,u_2`$ interchanged. $`\mathrm{}`$
###### Lemma 4.3
Let $`1p\mathrm{}`$, $`uL^p^{}(\mathrm{\Gamma })`$, and $`v_1,v_2L^p(\mathrm{\Gamma })`$. Then
$$|A(\mathrm{\Gamma },u,v_1)A(\mathrm{\Gamma },u,v_2)|2v_1v_2_pu_p^{}.$$
Proof. We have
$`A(\mathrm{\Gamma },u,v_1)`$ $`=`$ $`\underset{f_p=1}{sup}\underset{\alpha 𝐂}{inf}v_1(x)\left[{\displaystyle _a^x}u(t)f(t)𝑑t\alpha \right]_p`$
$`=`$ $`\underset{f_p=1}{sup}\underset{|\alpha |u_p^{}f_p}{inf}v_1(x)\left[{\displaystyle _a^x}u(t)f(t)𝑑t\alpha \right]_p.`$
Since
$$v_1(x)[_a^xu(t)f(t)𝑑t\alpha ]_pv_1v_2_pu_p^{}f_p+(v_1v_2)\alpha _p+v_2[_a^xu(t)f(t)𝑑t\alpha ]_p$$
it follows that
$`A(\mathrm{\Gamma },u,v_1)`$ $``$ $`2v_1v_2_pu_p^{}+\underset{f_p=1}{sup}\underset{|\alpha |u_p^{}}{inf}v_2(x)\left[{\displaystyle _a^x}u(t)f(t)𝑑t\alpha \right]_p`$
$`=`$ $`2v_1v_2_pu_p^{}+A(\mathrm{\Gamma },u,v_2).`$
Similarly with $`v_1`$ and $`v_2`$ interchanged. $`\mathrm{}`$
###### Lemma 4.4
Let $`1<p<\mathrm{}`$ and suppose that $`K_1,K_2,K_1K_2`$, are compact subtrees of $`\mathrm{\Gamma }`$. Then
$$|A(K_1)A(K_2)|0\text{ as }|K_2K_1|0$$
Proof. We see that $`A(K_1)=A(\mathrm{\Gamma },u\chi _{K_1},v\chi _{K_1})`$ and $`A(K_2)=A(\mathrm{\Gamma },u\chi _{K_2},v\chi _{K_2}).`$ The lemma then follows from Lemmas 4.2 and 4.3.$`\mathrm{}`$
In order to treat the cases $`p=1,\mathrm{}`$ we need the following terminology.
###### Definition 4.5
Let $`uL^{\mathrm{}}(\mathrm{\Gamma }).`$ Then
$$u_s:=\underset{\epsilon 0_+}{lim}u\chi _{B(t,\epsilon )}_{\mathrm{},\mathrm{\Gamma }}$$
where $`B(t,\epsilon )`$ is the ball center $`t`$, radius $`\epsilon `$ on $`\mathrm{\Gamma }`$.
###### Definition 4.6
Let $`g`$ be a function defined on a real interval $`I`$. Then
$$g^{}(x):=inf\{t;g_{}(t)x\},$$
where $`g_{}(t):=|\{xI;g(x)t\}|.`$ The function $`g^{}`$ is the non–increasing rearrangement of $`g`$.
Note that since we have $``$ in the definitions above, $`g_{}`$ and $`g^{}`$ are left–continuous functions. For the case $`p=\mathrm{}`$ we have the following two lemmas.
###### Lemma 4.7
Let $`I`$ be a bounded interval, $`\gamma ,\delta 𝐑`$ with $`\delta v_s(t)0`$ on $`I`$, and let $`p=\mathrm{}`$. Then
$$A(I;\gamma ,\delta )A(I;\gamma ,v_s)1/2|\gamma |(v_s\chi _I)^{}(t)t_{\mathrm{},(0,|I|)}.$$
Proof. See \[6; Lemma 4.5\]. $`\mathrm{}`$
###### Lemma 4.8
Let $`I`$ be a bounded interval, $`\gamma ,\delta 𝐑`$ with $`\delta v_s(t)0`$ on $`I`$ and let $`p=\mathrm{}`$. Then for any $`\alpha >1`$
$$A(I;\gamma ,\delta )A(I;\gamma ,v_s)\frac{\alpha }{2}|\gamma |(\delta v_s(t))𝑑t+\frac{|\gamma |\delta |I|}{2\alpha }.$$
Proof. See \[6, Lemma 4.6\]. $`\mathrm{}`$
For the case $`p=1`$ we have
###### Lemma 4.9
Let $`I`$ be a bounded interval, $`\gamma ,\delta 𝐑`$ with $`\delta u_s(t)0`$ and $`\mathrm{}>u_1(t)u_2(t)0`$ on $`I`$. Then for $`p=1`$
$$A(I;u_1,\gamma )A(I;u_2,\gamma )$$
and
$$A(I;u_s,\gamma )1/4|\gamma |(u_s\chi _I)^{}(t)t_{\mathrm{},(0,|I|)}.$$
Proof. In the first inequality
$$A(I;u_1,\gamma )=|\gamma |\underset{f_{1,I}1}{sup}\underset{\alpha 𝐂}{inf}\frac{_a^xu_1f\alpha _{1,I}}{f_{1,I}}.$$
For any $`f_2_{1,I}1`$ there exists $`f_1`$ such that $`f_1_{1,I}f_2_{1,I}1`$ and $`_a^xu_1f_1=_a^xu_2f_2`$. (Put $`f_1(t):=f_2(t)u_2(t)/u_1(t)`$ if $`u_1(t)0`$ and $`f_1(t):=f_2(t)`$ otherwise.) Then $`A(I;u_1,\gamma )A(I;u_2,\gamma )0`$.
In the second inequality, we have
$`A(I;u_s,\gamma )`$ $`=`$ $`|\gamma |\underset{f_{1,I}=1}{sup}\underset{\alpha 𝐂}{inf}{\displaystyle _a^x}u_sf\alpha _{1,I}`$
$``$ $`|\gamma |\underset{f_{1,I}=1}{sup}\underset{\alpha 𝐂}{inf}{\displaystyle _a^x}\chi _{M_\beta }\beta f\alpha _{1,I}`$
where $`M_\beta :=\{yI;u_s(y)\beta \}`$ and $`0\beta u_s_{\mathrm{},I}`$. Put $`f=\delta _{x_\beta }`$, where $`x_\beta I=(a,b)`$ and $`|M_\beta (a,x_\beta )|=|M_\beta (x_\beta ,b)|=1/2|M_\beta |.`$ Then
$`A(I;u_s,\gamma )`$ $``$ $`|\gamma |\underset{\alpha 𝐂}{inf}\chi _{(x_\beta ,b)}\beta \alpha _{1,I}`$
$`=`$ $`|\gamma |\underset{\beta >\alpha >0}{inf}(\alpha (x_\beta a),(\beta \alpha )(bx_\beta ))`$
$`=`$ $`|\gamma ||\beta |\underset{1>\alpha >0}{inf}(\alpha (x_\beta a),(1\alpha )(bx_\beta ))`$
$`=`$ $`|\gamma ||\beta |{\displaystyle \frac{bx_\beta }{ba}}(x_\beta a)`$
$``$ $`|\gamma ||\beta ||{\displaystyle \frac{M_\beta }{2}}|{\displaystyle \frac{ba|M_\beta /2|}{ba}}`$
$``$ $`|\gamma ||\beta ||M_\beta |{\displaystyle \frac{1}{2}}{\displaystyle \frac{ba(\frac{ba}{2})}{ba}}`$
$`=`$ $`{\displaystyle \frac{1}{4}}|\gamma ||\beta ||M_\beta |.`$
Hence, we have for every $`0\beta u_s_{\mathrm{},I}`$
$$A(I;u_s,\gamma )|\gamma ||\beta ||M_\beta |\frac{1}{4}$$
and so
$$A(I;u_s,\gamma )|\gamma |\frac{1}{4}t(u_s\chi _I)^{}(t)_{\mathrm{},(0,|I|)}.$$
$`\mathrm{}`$
###### Lemma 4.10
Let $`I`$ be a bounded interval, $`\gamma ,\delta 𝐑`$ with $`\delta u_s(t)0`$ on $`I`$ and $`p=1`$. Then for any $`\alpha >1`$
$$A(I;\delta ,\gamma )A(I;u_s,\gamma )\frac{\alpha }{2}_I|\gamma |(\delta u_s(t))𝑑t+\frac{|\gamma |\delta |I|}{2\alpha }.$$
Proof. From Lemmas 4.9 and 4.1 we have
$$0A(I;\delta ,\gamma )A(I;u_s,\gamma )\frac{1}{2}|\gamma ||\delta ||I|\frac{1}{4}|\gamma |(u_s\chi _I)^{}(t)t_{\mathrm{}}.$$
The rest of the proof is similar to that in \[6, Lemma 4.6\] on using Lemma 4.8 instead of \[6, Lemma 4.5\]. $`\mathrm{}`$
## 5 Local asymptotic results
5.1 Case $`1<p<\mathrm{}`$
###### Lemma 5.1
Let $`K`$ be a compact subtree of $`\mathrm{\Gamma }`$ and $`1<p<\mathrm{}`$. Then
$$\alpha _p_K|u||v|=\underset{\epsilon 0_+}{lim}\epsilon N(K,\epsilon ,u,v),$$
$$\alpha _p_K|u||v|=\underset{\epsilon 0_+}{lim}\epsilon M(K,\epsilon ,u,v),$$
where $`N(K,\epsilon ,u,v)`$ and $`M(K,\epsilon ,u,v)`$ are defined in Definition 3.11, and $`\alpha _p=A((0,1),1,1)`$ (see Lemma 4.1).
Proof. Since $`K`$ is a compact tree it has a bounded number of vertices, i.e. $`K`$ is a finite union of intervals. The argument in \[2, Theorem 5\], with A replacing the function $`L`$ there, continues to go through and yields the first equation of the lemma. The second identity follows from the first identity and Lemma 3.10 since $`A`$ is a continuous function on an interval. $`\mathrm{}`$
###### Lemma 5.2
Let $`1<p<\mathrm{}`$. Then
$$\alpha _p_\mathrm{\Gamma }|u||v|\underset{\epsilon 0+}{lim\; inf}\epsilon N(\mathrm{\Gamma },\epsilon ,u,v).$$
Proof. There exist compact subtrees $`\mathrm{\Gamma }_n`$ of $`\mathrm{\Gamma }`$, $`n=1,2,\mathrm{}`$ such that $`\mathrm{\Gamma }_n\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }_n\mathrm{\Gamma }`$ (i.e. $`|\mathrm{\Gamma }\mathrm{\Gamma }_n|0`$) as $`n\mathrm{}`$. By Lemma 5.1 we have
$$\alpha _p_{\mathrm{\Gamma }_n}|u||v|=\underset{\epsilon 0_+}{lim}\epsilon N(\mathrm{\Gamma }_n,\epsilon ,u,v)\underset{\epsilon 0_+}{lim\; inf}\epsilon N(\mathrm{\Gamma },\epsilon ,u,v),$$
whence the result. $`\mathrm{}`$
###### Theorem 5.3
Let $`1<p<\mathrm{}`$, $`uL^p^{}(\mathrm{\Gamma })`$ and $`vL^p(\mathrm{\Gamma })`$. Then
$$\underset{\epsilon 0_+}{lim}\epsilon N(\mathrm{\Gamma },\epsilon ,u,v)=\alpha _p_\mathrm{\Gamma }|u||v|,$$
$$\underset{\epsilon 0_+}{lim}\epsilon M(\mathrm{\Gamma },\epsilon ,u,v)=\alpha _p_\mathrm{\Gamma }|u||v|.$$
Proof. Let $``$ be a maximal such that $`\mathrm{\#}=M(\mathrm{\Gamma },\epsilon )M(\mathrm{\Gamma },\epsilon ,u,v)`$ and $`𝒮`$ a minimal cover such that $`\mathrm{\#}𝒮=N(\mathrm{\Gamma },\epsilon )N(\mathrm{\Gamma },\epsilon ,u,v)`$.
Given $`\eta >0`$, choose a compact subtree $`K\mathrm{\Gamma }`$ such that $`u\chi _{\mathrm{\Gamma }K}_p^{}\eta `$ and $`v\chi _{\mathrm{\Gamma }K}_p\eta `$.
Set
$`_1`$ $`=`$ $`\{\mathrm{\Gamma }^{}:\mathrm{\Gamma }^{},\mathrm{\Gamma }^{}K\},`$
$`_2`$ $`=`$ $`\{\mathrm{\Gamma }^{}:\mathrm{\Gamma }^{},\mathrm{\Gamma }^{}\mathrm{\Gamma }K\},`$
$`_3`$ $`=`$ $`\{_1_2\},`$
and similarly for $`𝒮_1,𝒮_2,𝒮_3`$ with respect to $`𝒮`$.
We know that $`\mathrm{\Gamma }K`$ is the union of disjoint connected components $`\{\mathrm{\Gamma }_i^{}\}`$ and $`\mathrm{\#}\{\mathrm{\Gamma }_i^{}\}\mathrm{\#}K`$. Also, if $`X_2𝒮_2`$ then $`X\mathrm{\Gamma }_i^{}`$ for some $`i`$. Thus $`\mathrm{\#}_2_iM(\mathrm{\Gamma }_i^{},\epsilon ).`$
Consider now the union of $`𝒮_1,𝒮_3`$ and those subtrees in the definition of the $`N(\mathrm{\Gamma }_i^{},\epsilon )`$. This covers $`\mathrm{\Gamma }`$ and so
$$\mathrm{\#}𝒮_2\underset{i}{}N(\mathrm{\Gamma }_i^{},\epsilon )3\underset{i}{}M(\mathrm{\Gamma }_i^{},\epsilon )+3\mathrm{\#}K$$
(5. 14)
by Lemma 3.9. Let $`\mathrm{\Gamma }_i^{}(j)`$ be the subtrees in the definition of $`M(\mathrm{\Gamma }_i^{},\epsilon )`$. Then
$$\epsilon \mathrm{\#}_2\underset{i,j}{}A(\mathrm{\Gamma }_i^{}(j))\underset{i,j}{}u_{p^{},\mathrm{\Gamma }_i^{}(j)}v_{p,\mathrm{\Gamma }_i^{}(j)}u\chi _{\mathrm{\Gamma }K}_p^{}v\chi _{\mathrm{\Gamma }K}_p\eta ^2.$$
Since $`\mathrm{\#}_1M(K,\epsilon )`$
$$M(\mathrm{\Gamma },\epsilon )M(K,\epsilon )+\mathrm{\#}_2+\mathrm{\#}_3,$$
$$0<\epsilon [M(\mathrm{\Gamma },\epsilon )M(K,\epsilon )]\epsilon (\mathrm{\#}_2+\mathrm{\#}_3)\eta ^2+\epsilon \mathrm{\#}K.$$
Then, by Lemma 5.1,
$$0\underset{\epsilon 0}{lim\; sup}\epsilon M(\mathrm{\Gamma },\epsilon )\alpha _p_K|uv|\eta ^2.$$
Now let $`K\mathrm{\Gamma }`$ ($`\eta 0`$) to get
$$\underset{\epsilon 0}{lim\; sup}\epsilon M(\mathrm{\Gamma },\epsilon )=\alpha _p_\mathrm{\Gamma }|uv|.$$
We get the same for $`lim\; inf`$ and so
$$\underset{\epsilon 0}{lim}\epsilon M(\mathrm{\Gamma },\epsilon )=\alpha _p_\mathrm{\Gamma }|uv|.$$
Since $`N(K,\epsilon )+\mathrm{\#}𝒮_2+\mathrm{\#}𝒮_3N(\mathrm{\Gamma },\epsilon ),`$
$$0N(\mathrm{\Gamma },\epsilon )N(K,\epsilon )\mathrm{\#}𝒮_2+\mathrm{\#}𝒮_33\underset{i}{}M(\mathrm{\Gamma }_i^{}(j),\epsilon )+3\mathrm{\#}K+\mathrm{\#}K$$
by (5.1). Hence, as before,
$$\underset{\epsilon 0}{lim}\epsilon N(\mathrm{\Gamma },\epsilon )=\alpha _p_\mathrm{\Gamma }|uv|.$$
$`\mathrm{}`$
###### Corollary 5.4
Let $`1<p<\mathrm{}`$, $`uL^p^{}(\mathrm{\Gamma })`$ and $`vL^p(\mathrm{\Gamma })`$. Then
$$\underset{n\mathrm{}}{lim}na_n(T)=\alpha _p_\mathrm{\Gamma }|u||v|.$$
Proof. Note that the application of Theorem 3.18 to Theorem 5.3 implies that $`lim_n\mathrm{}a_n(T)=0`$ and hence that $`T`$ is compact.
Let $`\{\mathrm{\Gamma }_l\}_{l=1}^{\mathrm{}}`$ be as in the proof of Lemma 5.2, and set $`T_l=T_{a,\mathrm{\Gamma }_l}`$ for some $`a\mathrm{\Gamma }`$. Since $`\mathrm{\Gamma }_l`$ is compact then from Lemma 3.19 and Theorem 5.1 we have
$$\underset{n\mathrm{}}{lim}na_n(T_l)=\alpha _p_{\mathrm{\Gamma }_l}|u||v|.$$
An operator of rank $`<n`$ on $`\mathrm{\Gamma }_l`$ can be considered as the restriction to $`\mathrm{\Gamma }_l`$ of such an operator on $`\mathrm{\Gamma }`$ and also $`T|L^p(\mathrm{\Gamma })L^p(\mathrm{\Gamma })T_l|L^p(\mathrm{\Gamma }_l)L^p(\mathrm{\Gamma }_l)`$ if $`T=T_a`$ and $`a\mathrm{\Gamma }_l`$. It follows that $`a_n(T)a_n(T_l)`$ and so
$$\underset{n\mathrm{}}{lim\; inf}na_n(T)\alpha _p_{\mathrm{\Gamma }_l}|u||v|.$$
But, we know from Lemma 5.3 that
$$\underset{\epsilon 0_+}{lim}\epsilon N(\mathrm{\Gamma },\epsilon )=\alpha _p_\mathrm{\Gamma }|u||v|.$$
and so, by Lemma 3.16
$$\underset{n\mathrm{}}{lim\; sup}na_n(T)\alpha _p_\mathrm{\Gamma }|u||v|.$$
Hence the corollary is proved.$`\mathrm{}`$
5.2 The cases $`p=\mathrm{}`$ and $`p=1`$.
The analogies of Lemma 5.1 for $`p=\mathrm{}`$ and $`p=1`$ are, respectively
###### Lemma 5.5
Let $`K`$ be a compact subtree of $`\mathrm{\Gamma },`$ and $`p=\mathrm{}`$. Then
$$\frac{1}{2}_K|u||v_s|\underset{\epsilon 0_+}{lim\; inf}\epsilon N(K,\epsilon ,u,v)\underset{\epsilon 0_+}{lim\; sup}\epsilon N(K,\epsilon ,u,v)\frac{3}{2}_K|u||v_s|$$
where $`v_s`$ is defined in Definition 5.4.
###### Lemma 5.6
Let $`K`$ be a compact subtree of $`\mathrm{\Gamma }`$ and $`p=1`$. Then
$$\frac{1}{2}_K|u_s||v|\underset{\epsilon 0_+}{lim\; inf}\epsilon N(K,\epsilon ,u,v)\underset{\epsilon 0_+}{lim\; sup}\epsilon N(K,\epsilon ,u,v)\frac{3}{2}_K|u_s||v|.$$
Both lemmas follow from the results for intervals in since $`K`$ is a finite union of intervals. Lemmas 5.5 and 5.6 yield, as in Lemma 5.2,
###### Lemma 5.7
For $`p=\mathrm{}`$
$$\frac{1}{2}_\mathrm{\Gamma }|u||v_s|\underset{\epsilon 0_+}{lim\; inf}\epsilon N(\mathrm{\Gamma },\epsilon ,u,v)$$
and for $`p=1`$
$$\frac{1}{2}_\mathrm{\Gamma }|u_s||v|\underset{\epsilon 0_+}{lim\; inf}\epsilon N(\mathrm{\Gamma },\epsilon ,u,v).$$
###### Lemma 5.8
Let $`uL^p^{}(\mathrm{\Gamma })`$ and $`vL^p(\mathrm{\Gamma })`$. Then for $`p=\mathrm{}`$
$$\frac{1}{2}_\mathrm{\Gamma }|u||v_s|\underset{\epsilon 0_+}{lim\; inf}\epsilon N(\mathrm{\Gamma },\epsilon ,u,v)\underset{\epsilon 0_+}{lim\; sup}\epsilon N(\mathrm{\Gamma },\epsilon ,u,v)\frac{3}{2}_\mathrm{\Gamma }|u||v_s|$$
and for $`p=1`$
$$\frac{1}{2}_\mathrm{\Gamma }|u_s||v|\underset{\epsilon 0_+}{lim\; inf}\epsilon N(\mathrm{\Gamma },\epsilon ,u,v)\underset{\epsilon 0_+}{lim\; sup}\epsilon N(\mathrm{\Gamma },\epsilon ,u,v)\frac{3}{2}_\mathrm{\Gamma }|u_s||v|.$$
Proof. Let $`p=\mathrm{}`$. We need only prove the last inequality. Let $`\{\mathrm{\Gamma }_l\}_{l=1}^{\mathrm{}}`$ be compact subtrees of $`\mathrm{\Gamma }`$ which are such that
$$\left|_\mathrm{\Gamma }|u||v_s|_{\mathrm{\Gamma }_l}|u||v_s|\right|\frac{1}{l}$$
and
$$u_{1,\mathrm{\Gamma }\mathrm{\Gamma }_l}\frac{1}{l}.$$
Fix $`l𝐍`$. There exist intervals $`W(j)`$ in $`\mathrm{\Gamma }_l`$ and step functions $`u_n,v_n`$ on $`\mathrm{\Gamma }_l`$,
$$u_n=\underset{j=1}{\overset{m}{}}\xi _j\chi _{W(j)},v_n=\underset{j=1}{\overset{m}{}}\eta _j\chi _{W(j)},$$
which are such that
$$uu_n_{1,\mathrm{\Gamma }_l}<\frac{1}{n},_{\mathrm{\Gamma }_l}|u(t)|(v_n(t)v_s(t))dt<\frac{1}{n}$$
and $`v_s_{\mathrm{},\mathrm{\Gamma }}v_n(t)v_s(t)`$ on $`\mathrm{\Gamma }_l`$; cf \[6, Theorem 4.7\].
Let $`M:=M(\mathrm{\Gamma },\epsilon )`$ and let $`\{\mathrm{\Gamma }_i^M\}_{i=1}^M`$ be a maximal set of subtrees of $`\mathrm{\Gamma }`$ in the definition of $`M`$ (see Definition 3.11). Then, because $`\mathrm{\Gamma }_l`$ is a compact subtree of $`\mathrm{\Gamma }`$, we have $`M2m\mathrm{\#}V(\mathrm{\Gamma }_l)\mathrm{\#}\mathrm{\Gamma }_l\mathrm{\#}𝐊,`$ where
$$𝐊:=\{\mathrm{\Gamma }_j:\mathrm{\Gamma }_j\{\mathrm{\Gamma }_k^M\}_{k=1}^M,\text{ and there exists }i\text{ such that }W(i)\mathrm{\Gamma }_j\text{ or }\mathrm{\Gamma }_j\mathrm{\Gamma }\mathrm{\Gamma }_l\}.$$
On using Lemmas 4.6 and 4.7, we have
$`\epsilon (M2m`$ $``$ $`\mathrm{\#}V(\mathrm{\Gamma }_l)\mathrm{\#}\mathrm{\Gamma }_l){\displaystyle }_{k𝐊}A(\mathrm{\Gamma }_k^M,u,v)`$
$``$ $`{\displaystyle \underset{k𝐊}{}}(A(\mathrm{\Gamma }_k^M,u_n,v_n)+[A(\mathrm{\Gamma }_k^M,u,v)A(\mathrm{\Gamma }_k^M,u_n,v)]`$
$`+[A(\mathrm{\Gamma }_k^M,u_n,v)A(\mathrm{\Gamma }_k^M,u_n,v_n)])`$
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{m}{}}}|\xi _j||\eta _j||W(j)|`$
$`+{\displaystyle \underset{j=1}{\overset{M}{}}}\left(uu_n_{1,\mathrm{\Gamma }_j^M}v_{\mathrm{},\mathrm{\Gamma }_J^M}\right)`$
$`+{\displaystyle \underset{k𝐊;i,W(i)\mathrm{\Gamma }_k^M}{}}\left[A(\mathrm{\Gamma }_k^M,u_n,v)A(\mathrm{\Gamma }_k^M,u_n,v_n)\right]`$
$`+{\displaystyle \underset{k𝐊;\mathrm{\Gamma }_k^M\mathrm{\Gamma }\mathrm{\Gamma }_l}{}}\left[A(\mathrm{\Gamma }_k^M,u_n,v)A(\mathrm{\Gamma }_k^M,u_n,v_n)\right]`$
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{m}{}}}|\xi _j||\eta _j||W(j)|+uu_n_1v_{\mathrm{}}`$
$`+{\displaystyle \underset{kK;i,W(i)\mathrm{\Gamma }_k^M}{}}\left[{\displaystyle \frac{\alpha }{2}}{\displaystyle _{\mathrm{\Gamma }_k^M}}(v_nv_s)|\xi _i|𝑑t+{\displaystyle \frac{|\xi _i||\eta _i|}{2\alpha }}|\mathrm{\Gamma }_k^M|\right]`$
$`+{\displaystyle \underset{k𝐊,\mathrm{\Gamma }_k^M\mathrm{\Gamma }\mathrm{\Gamma }_l}{}}\left[A(\mathrm{\Gamma }_k^M,0,v)A(\mathrm{\Gamma }_k^M,0,0)\right]`$
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{m}{}}}|\xi _j||\eta _j||W(j)|+({\displaystyle \frac{1}{n}}+{\displaystyle \frac{1}{l}})v_s_{\mathrm{}}`$
$`+{\displaystyle \frac{\alpha }{2}}{\displaystyle _\mathrm{\Gamma }}(v_nv_s)|u_n|𝑑t`$
$`+{\displaystyle \frac{1}{2\alpha }}{\displaystyle _\mathrm{\Gamma }}|u_n||v_n|𝑑t`$
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Gamma }}|u||v_s|+c\left(\alpha {\displaystyle \frac{1}{n}}+{\displaystyle \frac{1}{\alpha }}+{\displaystyle \frac{1}{n}}+{\displaystyle \frac{1}{l}}\right)`$
for some constant $`c`$ independent on $`\epsilon `$. We therefore conclude that
$$\frac{1}{3}\underset{\epsilon 0_+}{lim\; sup}\epsilon N(\mathrm{\Gamma },\epsilon ,u,v)\frac{1}{2}_\mathrm{\Gamma }|u||v_s|+K(\alpha \frac{1}{n}+\frac{1}{\alpha }),$$
whence
$$\frac{1}{3}\underset{\epsilon 0_+}{lim\; sup}\epsilon N(\mathrm{\Gamma },\epsilon ,u,v)\frac{1}{2}_\mathrm{\Gamma }|u||v_s|.$$
The case $`p=1`$ is similar. $`\mathrm{}`$
From Lemmas 5.8 and 3.18 we derive
###### Lemma 5.9
Let $`uL^p^{}(\mathrm{\Gamma })`$ and $`vL^p(\mathrm{\Gamma })`$. Then for $`p=\mathrm{}`$
$$\frac{1}{6}_\mathrm{\Gamma }|u||v_s|\underset{n\mathrm{}}{lim\; inf}na_n(T)\underset{n\mathrm{}}{lim\; sup}na_n(T)3_\mathrm{\Gamma }|u||v_s|$$
and for $`p=1`$
$$\frac{1}{6}_\mathrm{\Gamma }|u_s||v|\underset{n\mathrm{}}{lim\; inf}na_n(T)\underset{n\mathrm{}}{lim\; sup}na_n(T)3_\mathrm{\Gamma }|u_s||v|.$$
## 6 The main results for $`1<p<\mathrm{}`$.
We suppose throughout this section that $`T:=T_a`$ and $`\mathrm{\Gamma }:=\mathrm{\Gamma }_a`$ for some $`a\mathrm{\Gamma }`$. Also we write $`T_K`$ for $`T|L^p(K)L^p(K)`$, for any $`K\mathrm{\Gamma }`$.
With $`U(x):=_a^x|u(t)|^p^{}𝑑t`$ $`(x\mathrm{\Gamma })`$ we define $`Z_k`$ to be the closure of
$$\{x:x\mathrm{\Gamma },2^{\frac{kp^{}}{p}}U(x)<2^{\frac{(k+1)p^{}}{p}}\}.$$
(6. 15)
Here $`k`$ may be any integer if $`uL_{loc}^p^{}(\mathrm{\Gamma })L^p^{}(\mathrm{\Gamma })`$, while, if $`uL^p^{}(\mathrm{\Gamma })`$, $`2^ku_{p^{},\mathrm{\Gamma }}^p`$; we refer to these values of $`k`$ as the admissible values.
We have that $`Z_k=_{i=1}^{n_k}Z_{k,i}`$, where the $`Z_{k,i}`$ are the connected components of $`Z_k`$. Corresponding to each admissible $`k`$ we set
$$\sigma _{k,i}^p:=2^k\mu (Z_{k,i})\text{ for }i\{1,\mathrm{},n_k\}$$
(6. 16)
and
$$\sigma _k^p:=2^k\mu (Z_k).$$
(6. 17)
For non–admissible $`k`$ we set $`\sigma _k=0`$. We also set $`\sigma _{k,i}=0`$ for $`i\{1,\mathrm{},n_k\}`$.
Let
$$B_{k,i}:=\mathrm{\#}Z_{k,i}1;$$
(6. 18)
that is, $`B_{k,i}`$ is the number of boundary points of $`Z_{k,i}`$ excluding its root.
###### Lemma 6.1
$$\underset{k𝐙}{sup}\underset{1jn_k}{\mathrm{max}}\sigma _{k,i}T.$$
(6. 19)
Proof. This follows from \[7, Proposition 5.1\], which asserts that
$$\underset{x\mathrm{\Gamma }}{sup}u\chi _{(a,x)}_p^{}v\chi _{(a,x)^c}_pT,$$
where $`(a,x)^c=\{y\mathrm{\Gamma }:x_ay\}`$. For then, by (6. 15),
$`T`$ $``$ $`\underset{x\mathrm{\Gamma }}{sup}U(x)^{1/p^{}}\left({\displaystyle _{yx}}|v(y)|^p𝑑y\right)^{1/p}`$
$``$ $`\underset{k,i}{sup}2^{k/p}\mu (Z_{k,i})^{1/p}`$
$`=`$ $`\underset{k,i}{sup}\sigma _{k,i}.`$
$`\mathrm{}`$
###### Lemma 6.2
Let $`\mathrm{\Gamma }^{}`$ be a subtree of $`\mathrm{\Gamma }=\mathrm{\Gamma }_a`$ and $`b=b(\mathrm{\Gamma }^{})`$ the nearest point of $`\mathrm{\Gamma }^{}`$ to $`a`$. Then, for any $`c>4`$, there exist $`X=X(\mathrm{\Gamma }^{})𝐈_b(\mathrm{\Gamma }^{})`$ and $`k^{}=k^{}(\mathrm{\Gamma }^{})𝐙`$ such that, with $`T^{}=T_\mathrm{\Gamma }^{}`$ and $`Y=Y(\mathrm{\Gamma }^{})=\mathrm{\Gamma }^{}X,`$
$$T^{}2^{2/p}c\{\underset{iS}{}2^k^{}\mu (YZ_{k^{},i})\}^{1/p},$$
(6. 20)
where $`S=S(\mathrm{\Gamma }^{})=\{i:\mu (YZ_{k^{},i})>0\}.`$
Proof. From Theorem 2.4, for $`c>4`$, there exists $`X𝐈_b(\mathrm{\Gamma }^{})`$ such that
$`T^{}`$ $``$ $`c{\displaystyle \frac{\mu (Y)^{1/p}}{\alpha _X}}`$
$``$ $`c\underset{tX\{b\}}{\mathrm{min}}({\displaystyle _b^t}|u|^p^{}dx)^{1/p^{}}\mu (Y)^{1/p}`$
$``$ $`c\underset{tX\{b\}}{\mathrm{min}}[U(t)U(b)]^{1/p^{}}[{\displaystyle \underset{i,k}{}}\mu (YZ_{k,i})]^{1/p}`$
$``$ $`c\underset{t(X\{b\})Z_{\gamma _o}}{\mathrm{min}}[U(t)U(b)]^{1/p^{}}[{\displaystyle \underset{i,k}{}}\mu (YZ_{k,i})]^{1/p},`$
where $`\gamma _0=\mathrm{min}\{k:\mu (YZ_k)>0\}`$. Since $`\mu (\mathrm{\Gamma })<\mathrm{}`$, we may assume that $`Y`$ is compact and hence
$$\underset{k\gamma _0}{\mathrm{max}}\underset{i=1}{\overset{n_k}{}}2^k\mu (YZ_{k,i})$$
is attained, and so
$`{\displaystyle \underset{k,i}{}}\mu (YZ_{k,i})`$ $`=`$ $`{\displaystyle \underset{k\gamma _0}{}}2^k{\displaystyle \underset{i=1}{\overset{n_k}{}}}2^k\mu (YZ_{k,i})`$
$``$ $`2^{1\gamma _0}{\displaystyle \underset{i=1}{\overset{n_k^{}}{}}}2^k^{}\mu (YZ_{k^{},i})`$
say, for some $`k^{}\gamma _0.`$ Hence
$$T^{}c[2^{(\gamma _0+1)p^{}/p}2^{\gamma _0p^{}/p}]^{1/p^{}}[2^{1\gamma _0}\underset{i=1}{\overset{n_k^{}}{}}2^k^{}\mu (YZ_{k^{},i})]^{1/p},$$
whence (6. 20). $`\mathrm{}`$
###### Lemma 6.3
Let $`\{\mathrm{\Gamma }_i\}_{}`$ be a finite set of non-overlapping subtrees of $`\mathrm{\Gamma }`$ and set $`T_l=T_{\mathrm{\Gamma }_l}`$. Then,
$$\underset{l}{}T_l^q(2^{2/p+2})^q\underset{(k,i)\eta }{}B_{k,i}^{q/p^{}}\sigma _{k,i}^q\text{ if }1qp$$
(6. 21)
and
$$\underset{l}{}T_l^q(2^{2/p+2})^q\underset{k\eta _0}{}\sigma _k^q\text{ if }pq<\mathrm{},$$
(6. 22)
where $`\eta `$ and $`\eta _0`$ are finite sets.
Proof. Let $`\mathrm{\Gamma }_\lambda \{\mathrm{\Gamma }_l\}_{}`$, and, in the notation of Lemma 6.2, set $`b_l=b(\mathrm{\Gamma }_l),`$ $`k_l=k^{}(\mathrm{\Gamma }_l)`$, $`Y_l=Y(\mathrm{\Gamma }_l)`$ and $`S_l=S(\mathrm{\Gamma }_l)`$. There are two cases to consider for $`\mathrm{\Gamma }_\lambda `$:
$`b_\lambda Z_{k_\lambda }.`$ In this case $`b_\lambda Z_{k_\lambda ,i_\lambda },`$ $`S_\lambda =S(\mathrm{\Gamma }_\lambda )=\{i_\lambda \}`$ and, for any $`c>4`$,
$$T_\lambda (2^{2/p}c)\sigma _{k_\lambda ,i_\lambda }.$$
$`b_\lambda Z_{k_\lambda }.`$ Denote by $`\mathrm{\Lambda }`$ the subset of $``$ which is such that for $`l\mathrm{\Lambda }`$, $`b_lZ_{k_\lambda ,i_l}`$ for some unique $`i_lS_\lambda `$ and so $`S_l=\{i_l\}`$.
Set $`\mathrm{\Lambda }_i=\{l\mathrm{\Lambda }:i_l=i\}.`$ Then, by (6. 20), for $`q1`$,
$`{\displaystyle \underset{l\mathrm{\Lambda }}{}}T_l^q`$ $`=`$ $`{\displaystyle \underset{iS_\lambda }{}}{\displaystyle \underset{l\mathrm{\Lambda }_i}{}}T_l^q`$
$``$ $`(2^{2/p}c)^q{\displaystyle \underset{iS_\lambda }{}}{\displaystyle \underset{l\mathrm{\Lambda }_i}{}}\left\{2^{k_\lambda }\mu (Y_łZ_{k_\lambda ,i})\right\}^{q/p}`$
$``$ $`(2^{2/p}c)^q{\displaystyle \underset{iS_\lambda }{}}\left\{{\displaystyle \underset{l\mathrm{\Lambda }_i}{}}\left[2^{k_\lambda }\mu (Y_łZ_{k_\lambda ,i})\right]^{1/p}\right\}^q`$
$``$ $`(2^{2/p}c)^q{\displaystyle \underset{iS_\lambda }{}}\left[\left\{{\displaystyle \underset{l\mathrm{\Lambda }_i}{}}2^{k_\lambda }\mu (Y_łZ_{k_\lambda ,i})\right\}^{1/p}\left\{{\displaystyle \underset{l\mathrm{\Lambda }_i}{}}1\right\}^{1/p^{}}\right]^q`$
$``$ $`(2^{2/p}c)^q{\displaystyle \underset{iS_\lambda }{}}\left(B_{k_\lambda ,i}^{1/p^{}}\sigma _{k_\lambda ,i}\right)^q.`$ (6. 24)
Also in case $`(ii)`$, from (6. 20),
$$T_\lambda (2^{2/p}c)(\underset{iS_\lambda }{}\sigma _{k_\lambda ,i}^p)^{1/p}.$$
Hence, if $`1qp`$,
$`T_\lambda ^q`$ $``$ $`(2^{2/p}c)^q({\displaystyle \underset{iS_\lambda }{}}\sigma _{k_\lambda ,i}^q)`$ (6. 25)
$``$ $`(2^{2/p}c)^q{\displaystyle \underset{iS_\lambda }{}}\left(B_{k_\lambda ,i}^{1/p^{}}\sigma _{k_\lambda ,i}\right)^q.`$
If $`qp`$, then from (6),
$`{\displaystyle \underset{l\mathrm{\Lambda }}{}}T_l^q`$ $``$ $`(2^{2/p}c)^q{\displaystyle \underset{iS_\lambda }{}}\left\{{\displaystyle \underset{l\mathrm{\Lambda }_i}{}}2^{k_\lambda }\mu (Y_lZ_{k_\lambda ,i})\right\}^{q/p}`$ (6. 26)
$``$ $`(2^{2/p}c)^q{\displaystyle \underset{iS_\lambda }{}}\sigma _{k_\lambda ,i}^q`$
$``$ $`(2^{2/p}c)^q\left({\displaystyle \underset{iS_\lambda }{}}\sigma _{k_\lambda ,i}^p\right)^{q/p}`$
$``$ $`(2^{2/p}c)^q\sigma _{k_\lambda }^q.`$
Also, by (6. 20),
$`T_\lambda ^q`$ $``$ $`(2^{2/p}c)^q\left\{{\displaystyle \underset{lS_\lambda }{}}2^{k_\lambda }\mu (Y_\lambda Z_{k_\lambda ,i})\right\}^{q/p}`$ (6. 27)
$``$ $`(2^{2/p}c)^q\sigma _{k_\lambda }^q.`$
The lemma follows from (6. 24)-(6. 27) since $`c>4`$ is arbitrary. $`\mathrm{}`$
###### Theorem 6.4
For $`1<p<\mathrm{}`$, let $`u,v`$ satisfy (2. 5) and suppose that $`B_{k,i}^{1/p^{}}\sigma _{k,i}l^1(𝐙\times 𝐍)`$. Then
$$\underset{\epsilon 0}{lim}\epsilon M(\mathrm{\Gamma },\epsilon )=\alpha _p_\mathrm{\Gamma }|u||v|,$$
(6. 28)
$$\underset{\epsilon 0}{lim}\epsilon N(\mathrm{\Gamma },\epsilon ,)=\alpha _p_\mathrm{\Gamma }|u||v|,$$
(6. 29)
and
$$\underset{n\mathrm{}}{lim}na_n(T)=\alpha _p_\mathrm{\Gamma }|u||v|.$$
(6. 30)
Proof. Given $`\eta >0`$, we choose $`l`$ to be such that $`_{kl}_{i=1}^{n_k}B_{k,i}^{1/p^{}}\sigma _{k,i}\eta `$ and set $`K=_{kl}Z_k`$. Then, in the notation of the proof of Theorem 5.3, we have by Corollary 3.3, that
$`\epsilon \mathrm{\#}_2`$ $``$ $`{\displaystyle \underset{i,j}{}}A(\mathrm{\Gamma }_i^{(j)}){\displaystyle \underset{i,j}{}}T_{\mathrm{\Gamma }_i^{}(j)}`$
$``$ $`c{\displaystyle \underset{kl}{}}{\displaystyle \underset{i=1}{\overset{n_k}{}}}B_{k,i}^{1/p^{}}\sigma _{k,i}`$
for some positive constant $`c`$, by Lemma 6.3 with $`q=1`$. The proofs of the first two identities then follow that of Theorem 5.3. Theorem 3.18 and Lemma 3.16 complete the proof.
Note that the convergence of $`_{k𝐙}_{i=1}^{n_k}B_{k,i}^p^{}\sigma _{k,i}<\mathrm{}`$ implies that $`T`$ is compact. To see this, let $`K_j=_{kj}Z_k`$, and so
$`(T_{K_j}f)(x)`$ $`=`$ $`v(x)\chi _{K_j}(x){\displaystyle _a^x}f(t)u(t)\chi _{K_j}(t)𝑑t`$
$`=`$ $`v(x)\chi _{K_j}(x){\displaystyle _a^x}f(t)u(t)𝑑t.`$
Then
$$(TT_{K_j})f(x)=v(x)\chi _{\mathrm{\Gamma }K_j}(x)_a^xf(t)u(t)𝑑t$$
and, by Lemma 6.2, for some $`k^{}>j`$
$`TT_{K_j}`$ $``$ $`2^{2/p}c\left\{{\displaystyle \underset{iS}{}}2^k^{}\mu (Z_{k^{},i})\right\}^{1/p}`$
$``$ $`2^{2/p}c{\displaystyle \underset{iS}{}}\sigma _{k^{},i}2^{2/p}c{\displaystyle \underset{iS}{}}B_{k^{},i}^{1/p^{}}\sigma _{k^{},i}.`$
Thus $`TT_{K_j}0`$ as $`j\mathrm{}`$, and $`T`$ is compact since the $`T_{K_j}`$ are compact.$`\mathrm{}`$
###### Theorem 6.5
Let $`1<qp`$. Then, for some positive constant $`c`$,
$$\{a_n(T)\}_{l^q(𝐍)}cB_{k,i}^{1/p^{}}\sigma _{k,i}_{l^q(𝐙\times 𝐍)}\text{ if }1<qp$$
(6. 31)
Proof: Let $`\mathrm{\Gamma }_l,l=1,2,\mathrm{},M(\mathrm{\Gamma },\epsilon )`$ be a maximal set of subtrees of $`\mathrm{\Gamma }`$ from the definition of $`M(\mathrm{\Gamma },\epsilon )`$, so that $`A(\mathrm{\Gamma }_l)>\epsilon `$. Then, from (6. 21) and Corollary 3.3, for $`1qp`$,
$$\epsilon ^qM(\mathrm{\Gamma },\epsilon )\underset{l}{}T_l^qcB_{k,i}^{1/p^{}}\sigma _{k,i}_{l^q(𝐙\times 𝐍)}^q.$$
Since $`a_{3M(\mathrm{\Gamma },\epsilon )+4}(T)2\epsilon `$ by Lemma 3.9 and Theorem 3.18, it follows that
$`\mathrm{\#}\{m:a_m(T)>t\}`$ $``$ $`3M(\mathrm{\Gamma },t/2)+4`$
$``$ $`cM(\mathrm{\Gamma },t/2)`$
for $`c7`$. Thus
$$\mathrm{\#}\{m:a_m(T)>t\}ct^qB_{k,i}^{1/p^{}}\sigma _{k,i}_{l^q(𝐙\times 𝐍)}^q.$$
(6. 32)
We now proceed as in the proof of the Marcinkiewicz Interpolation Theorem (see ); we give the proof for completeness.
Define
$$v_1:=\{\begin{array}{cc}v\hfill & \text{ on }Z_{k,i}\text{ if }B_{k,i}^{1/p^{}}\sigma _{k,i}t/2,\hfill \\ 0\hfill & \text{ otherwise,}\hfill \end{array}$$
and set $`v_2=vv_1.`$ Denote $`T,\sigma _{k,i},`$ etc by $`T(v),\sigma _{k,i}(v)`$ to indicate the dependence on $`v`$, and set $`T^j=T(v_j),j=1,2.`$ Then, by \[3, Proposition II.2.2\],
$$a_{2n1}(T)a_n(T_1)+a_n(T_2)$$
and so
$$\{n:a_{2n1}(T)>t\}\{n:a_n(T_1)>t/2\}\{n:a_n(T_2)>t/2\}$$
and
$$\mathrm{\#}\{n:a_{2n1}(T)>t\}\mathrm{\#}\{n:a_n(T_1)>t/2\}+\mathrm{\#}\{n:a_n(T_2)>t/2\}.$$
(6. 33)
Set $`S_{k,i}=B_{k,i}^{1/p^{}}\sigma _{k,i}`$, and let $`1<q<q_1`$. Then, on using (6. 32) and (6. 33),
$`\{a_{2n1}(T)\}_{l^q(𝐍)}`$ $`=`$ $`q{\displaystyle _0^{\mathrm{}}}t^{q1}\mathrm{\#}\{n:a_{2n1}(T)>t\}𝑑t`$
$``$ $`cq{\displaystyle _0^{\mathrm{}}}t^{q1}\left\{t^{q_1}{\displaystyle \underset{S_{k,i}t/2}{}}S_{k,i}^{q_1}+t^1{\displaystyle \underset{S_{k,i}>t/2}{}}S_{k,i}\right\}𝑑t`$
$``$ $`c{\displaystyle \underset{k,i}{}}S_{k,i}^q,`$
whence (6. 31), since $`a_n(T)`$ decreases with $`n`$. $`\mathrm{}`$
###### Theorem 6.6
Let $`q(p,\mathrm{})`$. Then, for some positive constant $`c`$,
$$\{a_n(T)\}_{l^q(𝐍)}c\{\sigma _k\}_{l^q(𝐙)}.$$
(6. 34)
Proof. Let $`\{\mathrm{\Gamma }_l\}_1^{M(\mathrm{\Gamma },\epsilon )}`$ be as in the proof of Theorem 6.5 and define
$$F_j=\{\mathrm{\Gamma }_l:k^{}(\mathrm{\Gamma }_l)=k_j\}$$
in the notation of Lemma 6.2. Then, from the proof of Lemma 6.3, for $`q(p,\mathrm{}),`$
$`\epsilon ^p\mathrm{\#}F_j`$ $``$ $`{\displaystyle \underset{\mathrm{\Gamma }_lF_j}{}}T_l^p`$
$``$ $`c^p\sigma _{k_j}^p.`$
Thus, with $`m_j=[c^p\sigma _{k_j}^p/\epsilon ^p],`$
$$M(\mathrm{\Gamma },\epsilon )=\underset{j=1}{\overset{\mathrm{}}{}}\underset{m=1}{\overset{m_j}{}}1\underset{m=1}{\overset{\mathrm{}}{}}\mathrm{\#}\{j:c\sigma _{k_j}>m^{1/p}\epsilon \}.$$
Hence
$`\{a_n(T)\}_{l^q(𝐍)}`$ $`=`$ $`q{\displaystyle _0^{\mathrm{}}}t^{q1}\mathrm{\#}\{n:a_n(T)>t\}𝑑t`$
$``$ $`c{\displaystyle _0^{\mathrm{}}}t^{q1}M(\mathrm{\Gamma },t/2)𝑑t`$
$``$ $`c{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}t^{q1}\mathrm{\#}\{j:\sigma _j>m^{1/p}t\}dt`$
$``$ $`c{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}m^{q/p}\mathrm{\#}\{j:\sigma _j>t\}t^{q1}dt`$
$``$ $`c\{\sigma _k\}_{l^q(𝐙)}.`$
$`\mathrm{}`$
In the next theorem $`l_\omega ^q`$ denotes weak-$`l^q`$, that is, the space of sequences $`\{x_k\}`$ such that
$$\{x_k\}_{l_\omega ^q}:=\underset{t>0}{sup}\{t(\mathrm{\#}\{k:|x_k|>t\})^{1/q}\}<\mathrm{}.$$
###### Theorem 6.7
For some positive constant $`c`$,
$$\{a_n(T)\}_{l_\omega ^q(𝐍)}c\underset{i=1}{\overset{n_k}{}}B_{k,i}^{1/p^{}}\sigma _{k,i}_{l_\omega ^q(𝐙)}\text{ if }1<qp,$$
(6. 35)
and
$$\{a_n(T)\}_{l_\omega ^q(𝐍)}c\{\sigma _k\}_{l_\omega ^q(𝐙)}\text{ if }p<q<\mathrm{}.$$
(6. 36)
Proof. Let $`\{\mathrm{\Gamma }_l\}_1^{M(\mathrm{\Gamma },\epsilon )},F_j`$ be as in the proof of Theorem 6.6. Then, from the proof of Lemma 6.3,
$`\epsilon \mathrm{\#}F_j`$ $``$ $`{\displaystyle \underset{\mathrm{\Gamma }_lF_j}{}}T_l`$
$``$ $`c{\displaystyle \underset{i=1}{\overset{n_{k_j}}{}}}B_{k_j,i}^{1/p^{}}\sigma _{k_j,i}=:N_j`$
say. Thus
$`\epsilon ^q\mathrm{\#}\{n:a_n(T)>\epsilon \}`$ $``$ $`c\epsilon ^qM(\mathrm{\Gamma },\epsilon /2)`$
$``$ $`c\epsilon ^q{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=1}{\overset{[N_j/\epsilon ]}{}}}1`$
$``$ $`c\epsilon ^q{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}\mathrm{\#}\{j:N_j>m\epsilon \}`$
$``$ $`c{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}m^qt^q\mathrm{\#}\{j:N_j>t\}`$
and hence (6. 35). The proof of (6. 36) is similar, starting from
$$\epsilon ^p\mathrm{\#}F_jc\sigma _{k_j}^p.$$
$`\mathrm{}`$
Let us now suppose that the tree $`\mathrm{\Gamma }`$ satisfies the following condition:
$$B_{k,i}<B<\mathrm{}\text{for each admissible }k\text{ and }i.$$
(6. 37)
Then with this condition we can get lower estimates in Theorems 6.5 and 6.6.
We need the following result which is similar to \[2,Lemma 20\].
###### Lemma 6.8
Suppose that (6.10) is satisfied. Let $`S(\epsilon ):=\{(k,i):\sigma _{k,i}>\epsilon \}`$. If $`M+1/2\mathrm{\#}S(\epsilon )/4B`$, then $`a_M(T)>c\epsilon `$, where $`c`$ is an absolute constant.
Proof. It is sufficient to prove the result for $`S(\epsilon )`$ finite, for this will imply the result when $`\mathrm{\#}S(\epsilon )=\mathrm{}`$. The elements of $`S(\epsilon )`$ fall into two subsets according as $`k`$ is odd or even. At least one of them, say $`S_1(\epsilon )`$, has cardinal at least half that of $`S(\epsilon )`$. Thus, we may suppose that $`\mathrm{\#}S_1(\epsilon )>B`$.
Denote by $`\zeta _{k,j}`$ the point of $`Z_{k,j}`$ nearest to $`a`$, and define $`n(x):=\mathrm{\#}\{(k,j):\zeta _{k,j}_ax,(k,j)S_1(\epsilon )\}`$. Let $`l`$ be a path in $`\mathrm{\Gamma }`$ starting at $`a`$ and consisting of edges $`(x,y)`$ of $`\mathrm{\Gamma }`$, $`(x,y)`$ at each stage chosen so that $`n(y)`$ is as large as possible. Terminate the path at the point $`x=\zeta _{r,s}`$ at which $`n(x)=0`$. Define $`\xi l`$ by
$$\xi :=inf\{xl:n(x)=n(\zeta _{r1,j}),\zeta _{r1,j}l\},$$
( the infimum, which is being taken with respect to the total ordering on $`l`$ induced by $`_a`$, exists since $`n(a)>B`$ and $`n(\zeta _{r1,j})B`$ ). There are two possibilities : (i) $`\xi `$ may be a point $`\zeta _{k,j}`$, in which case define $`\mathrm{\Gamma }_1:=\{x:x_a\xi \}`$, or (ii) $`\xi `$ may be a vertex of $`\mathrm{\Gamma }`$ joined by a path $`l_1`$ to a point $`\zeta _{m,n}_a\xi `$, where $`(m,n)S_1(\epsilon )`$. In the latter case we define $`\mathrm{\Gamma }_1:=\{x:x_ay,yll_1\}`$. Then, in both cases, the closure of $`\mathrm{\Gamma }_1`$ is a subtree and so is its complement. Moreover, $`A(\overline{\mathrm{\Gamma }_1})>c\epsilon `$, where $`c`$ is an absolute constant. For, in case (i), if $`b`$ is a point of $`l`$ with $`U(b)=2^{(r1/2)(p^{}/p)}`$ and $`T^1`$ is the restriction of $`T`$ to $`\mathrm{\Gamma }_1`$, then, in the notation of the discussion preceding Lemma 3.5,
$$T_{b,1}^1,T_{b,2}^1(2^{r(p^{}/p)}2^{(r1/2)(p^{}/p)})^{1/p^{}}(2^{r/p}\epsilon );$$
this follows from \[7, Proposition 5.1\] where it is shown that
$$T_a\underset{x\mathrm{\Gamma }}{sup}u\chi _(a,x)_p^{}v\chi _{(a,x)^c}_p.$$
In case (ii) a similar result holds if $`b`$ is a point of $`ll_1`$ with $`U(b)=2^{(t1/2)(p^{}/p)}`$ and $`t`$ the greater of $`r,m`$. The lower bound for $`A(\overline{\mathrm{\Gamma }_1})`$ is then a consequence of Lemma 3.5. Note also that $`\mathrm{\Gamma }_1`$ contains at most $`2B`$ elements of $`S_1(\epsilon )`$.
The result now follows by induction on $`\mathrm{\#}S(\epsilon )`$ and Lemma 3.17. $`\mathrm{}`$
###### Lemma 6.9
Suppose that (6.10) is satisfied. Then, for all $`t>0`$,
$$\mathrm{\#}\{(k,i):\sigma _{k,i}>t\}4B\mathrm{\#}\{k𝐍:a_k(T)>ct\}+6B.$$
Proof. From Lemma 6.8,
$`\mathrm{\#}\{k𝐍:a_k(T)>ct\}`$ $``$ $`\left[\mathrm{\#}S(t)/4B1/2\right]`$
$``$ $`\mathrm{\#}S(t)/4B3/2,`$
whence the result. $`\mathrm{}`$
###### Lemma 6.10
Suppose that (6.10) is satisfied. Then, for all $`q>0`$,
$$\{\sigma _{k,i}\}_{l^q(𝐙\times 𝐍)}^qc_1\{a_k(T)\}_{l^q(𝐍)}^q+c_2\{\sigma _{k,i}\}_{l^{\mathrm{}}(𝐙\times 𝐍)}^q$$
Proof. Let $`\lambda =\{\sigma _{k,i}\}_{l^{\mathrm{}}(𝐙\times 𝐍)}`$. Then, by Lemma 6.9,
$`\{\sigma _{k,i}\}_{l^q(𝐙\times 𝐍)}^q`$ $``$ $`q{\displaystyle _0^\lambda }t^{q1}\mathrm{\#}\{(k,i)𝐙\times 𝐍:\sigma _{k,i}>t\}𝑑t`$
$``$ $`4Bq{\displaystyle _0^\lambda }t^{q1}\mathrm{\#}\{k𝐍;a_k(T)>c\epsilon \}𝑑t+6B\lambda ^q`$
$``$ $`c_1\{a_k(T)\}_{l^q(𝐍)}^q+c_2\lambda ^q.`$
$`\mathrm{}`$
###### Theorem 6.11
Let $`1<p<\mathrm{}`$ and suppose that (6.10) is satisfied. Then, for any $`q>0`$, there exists a constant $`c>0`$ such that
$$\{\sigma _{k,i}\}_{l^q(𝐙\times 𝐍)}c\{a_k(T)\}_{l^q(𝐍)}.$$
Proof. By Lemma 6.1,
$$\{\sigma _{k,i}\}_{l^{\mathrm{}}(𝐙\times 𝐍)}T=a_1(T)\{a_k(T)\}_{l^q(𝐍)}.$$
The result then follows from Lemma 6.10. $`\mathrm{}`$
###### Remark 6.12
(i) It follows from Theorem 6.5 and 6.11 that if (6. 37)is satisfied and $`1<qp`$, then
$$\{a_n(T)\}_{l^q(𝐍)}\{\sigma _{k,i}\}_{l^q(𝐙\times 𝐍)}.$$
For $`q>p`$, we have from Theorems 6.11 and 6.6
$`c_1\sigma _{k_i}_{l^q(𝐙\times 𝐍)}`$ $``$ $`\{a_n(T)\}_{l^q(𝐍)}c_2\sigma _k_{l^q(𝐙)}`$ (6. 38)
$``$ $`\sigma _{k,i}_{l^p(𝐙\times 𝐍)}.`$
(ii) Naimark and Solomyak take $`u=1`$, and in \[8,(4.8)\] they make the assumption that, for every edge $`y,zE(\mathrm{\Gamma }),`$
$$\mu _1|z|/|y|\mu _2,1<\mu _1\mu _2,$$
(6. 39)
where $`|y|,|z|`$ denote the lengths of the paths from the root of $`\mathrm{\Gamma }`$ to $`y,z`$ respectively. Let $`y_jV(t),`$ $`j=0,1,\mathrm{},`$ and suppose that $`|y_0|2^k`$ and $`|y_1|2^k`$. Then (6. 39) implies that
$$|y_n|\mu _1^{n1}|y_1|\mu _1^{n1}2^k2^{k+1}$$
if $`n1+\mathrm{log}2/\mathrm{log}\mu _1`$. Hence, if each vertex has constant branching number $`b`$ (ie. degree $`b+1`$), then
$$B_{k,i}b^{[\mathrm{log}2/log\mu _1+1]}$$
and hence (6. 37) is satisfied.
(iii) Theorem 4.1 in is valid under assumptions made on a sequence $`\{\eta _j\}`$ which is defined as follows : for any partition $`\mathrm{\Xi }`$ of $`\mathrm{\Gamma }`$ into a countable union of non-overlapping segments $`I_j=y_j,z_j`$,
$$\eta _j:=|z_j|_{I_j}v^2𝑑t.$$
Note that in our notation, the case $`p=2,u=1`$ is what is considered in . It is proved in \[9, Theorem 4.1 (i)\] that (6.15) for $`p=2`$ holds if, for some $`\mathrm{\Xi }`$, $`\{\eta _j\}l_{1/2}.`$
Choose $`\mathrm{\Xi }=_{k𝐍}\mathrm{\Xi }_k`$, where $`\mathrm{\Xi }_k`$ is a partition of $`Z_k`$. Then,
$`\sigma _{k,j}^2`$ $`=`$ $`2^k{\displaystyle \underset{I_sZ_{k,j}}{}}{\displaystyle _{I_s}}v^2𝑑t{\displaystyle \underset{I_sZ_{k,j}}{}}\eta _s`$
$``$ $`\left({\displaystyle \underset{I_sZ_{k,j}}{}}\eta _s^{1/2}\right)^2`$
and
$$\underset{k,j}{}\sigma _{k,j}\underset{s}{}\eta _s^{1/2}.$$
Thus, if (6. 37) is satisfied,
$$\underset{k,j}{}B_{k,j}^{1/2}\sigma _{k,j}B^{1/2}\underset{s}{}\eta _s^{1/2}.$$
(6. 40)
In the reverse directions we have
$$\sigma _{k,j}^21/2\underset{I_sZ_{k,j}}{}\eta _s$$
and so
$$\sigma _{k,j}2^{1/2}(\underset{I_sZ_{k,j}}{}\eta _s^{1/2})(\underset{I_sZ_{k,j}}{}1)^{1/2}.$$
Therefore
$$B_{k,j}^{1/2}\sigma _{k,j}2^{1/2}\left(\frac{B_{k,j}}{_{I_sZ_{k,j}}1}\right)^{1/2}\underset{I_sZ_{k,j}}{}\eta _s^{1/2}.$$
If
$$\underset{k,j}{inf}\left(\frac{B_{k,j}}{_{I_sZ_{k,j}}1}\right)=:c>0$$
(6. 41)
then
$$\underset{k,j}{}B_{k,j}^{1/2}\sigma _{k,j}(c/2)^{1/2}\underset{s}{}\eta _s^{1/2}.$$
(6. 42)
The condition (6. 41) is satisfied if the tree $`\mathrm{\Gamma }`$ is, in the terminology of , $`b`$ regular of type $`(b,2)`$ and $`\mathrm{\Xi }`$ consists of edges of $`\mathrm{\Gamma }`$. This means that every vertex of $`\mathrm{\Gamma }`$ has fixed branching number $`b`$, and any edge $`y,z`$ of the k-th generation is such that $`|y|=2^k,|z|=2^{k+1}.`$ Hence, in this case, (6. 41) is satisfied with $`c=1`$.
Acknowledgment. J. Lang wishes to record his gratitude to the Royal Society and NATO for support to visit the School of Mathematics at Cardiff during 1997/8, under their Postdoctoral Fellowship programmes. He also thanks the Grant Agency of the Czech Republic for partial support under grants 201/96/0431 and 201/98/P017.
REFERENCES
1. D.E.Edmunds, W.D.Evans and D.J.Harris. Approximation numbers of certain Volterra integral operators. J. London Math. Soc. (2) 37 (1988), 471–489.
2. D.E.Edmunds, W.D.Evans and D.J.Harris. Two–sided estimates of the approximation numbers of certain Volterra integral operators. Studia Math. 124 (1) (1997), 59–80.
3. D.E.Edmunds and W.D.Evans, Spectral Theory and Differential Operators, Oxford Univ. Press, Oxford, 1987.
4. D.E.Edmunds, P.Gurka and L.Pick. Compactness of Hardy–type integral operators in weighted Banach function spaces. Studia Math. 109 (1) (1994), 73–90.
5. W.D.Evans and D.J.Harris, Fractals, trees and the Neumann Laplacian. Math. Ann. 296 (1993), 493-527.
6. W.D.Evans, D.J.Harris and J.Lang, Two–sided estimates for the approximation numbers of Hardy–type operators in $`L^{\mathrm{}}`$ and $`L^1`$. Studia Math. 130 (2) (1998), 171-192.
7. Evans, W.D.; Harris, D.J.; Pick, L. Weighted Hardy and Poincar inequalities on trees. J. Lond. Math. Soc. 52 (2)(1995), 121-136.
8. K.Naimark and M.Solomyak, Eigenvalue Estimates for the weighted Laplacian on metric trees. Preprint.
9. J.Newman and M.Solomyak, Two–sided estimates of singular values for a class of integral operators on the semi–axis, Integral Equations Operator Theory, 20 (1994), 335–349.
10. B.Opic and A.Kufner, Hardy–type Inequalities, Pitman Res. Notes Math. Ser. 219, Longman Sci. & Tech., Harlow, 1990.
11. E.M.Stein, Singular Integrals and Differentiability Properties of Functions, Princeton Univ. Press, Princeton, 1970
W. D. Evans, D. J. Harris
School of Mathematics
Cardiff University
Senghennydd Road
Cardiff CF24 4YH
Wales, UK
E-mail : EvansWD@cardiff.ac.uk
J. Lang
Mathematics Department
202 Mathematical Sciences Bldg
University of Missouri
Columbia, MO 65211 USA
(permanent address: Math. Inst. AV CR, Zitna 25, Prague 1, Czech Republic)
E-mail : langjan@math.missouri.edu
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# 1 Introduction
## 1 Introduction
The investigation of the $`O(N)`$ vector model at large $`N`$ has a long-standing history in quantum field theory . One of the main aspects was the question of symmetry restoration at high temperature that for some time was controversial. The dynamical exploration of a special class of nonequilibrium properties has been developed only recently .
The out-of equilibrium configuration that has been studied mainly is characterized by an initial state in which one of the components has a spatially homogeneous classical expectation value $`\varphi (t)`$. This implies that the other $`N1`$ components $`\psi _i(𝐱,t),i=1\mathrm{}N1`$ have a mass that is different from the mass in the ground state. This means that their initial state is related to the Fock space vacuum state by a Bogoliubov transformation. The evolution of the system is governed by the classical equation of motion for the field $`\varphi (t)`$ and by the mode equations for the quantum fields $`\psi (𝐱,t)`$. The expectation value $`\psi (𝐱,t)\psi (𝐱,t)`$ appears in both equations of motion, this constitutes the quantum back reaction. In the one-loop approximation, in contrast to the large-$`N`$ approximation, this quantum back reaction only appears in the classical equation of motion. This leads to decisive differences in the late time behavior.
We have previously carried out such dynamical computations for the $`O(N)`$ vector model in the limit of large $`N`$ at finite temperature for the case of unbroken symmetry, i.e., with a positive mass term. Here we will consider the case of spontaneously broken symmetry. In this case, at low temperatures the fields $`\psi _i(𝐱,t)`$ will be the Goldstone modes. This is the case for the ground state at $`T=0`$ and at finite temperature; for nonequilibrium initial states these modes become massless when the system settles to a stationary state at late times. Symmetry restoration happens at high temperature and at large values of the initial field $`\varphi (0)`$; then at late times these modes stay massive while the classical field vanishes, and thereby the spontaneous symmetry breaking disappears.
Our investigation, as well as the analogous ones at $`T=0`$, are limited to fields, masses (as solutions of the gap equation) and temperatures much smaller than the scale of the Landau ghost $`m_x=m_1\mathrm{exp}(8\pi ^2/\lambda )`$, where $`m_1`$ is a renormalization scale, taken of order $`\sqrt{\lambda }v`$. So the question of symmetry non-restoration at “really” high temperatures will not be addressed here.
The plan of the paper is as follows: in section 2 we introduce the model and set up the equations governing the nonequilibrium evolution. In section 3 we discuss the renormalization of the equations of motion and of the energy-momentum tensor, some details are referred to Appendix A. In section 4 we discuss the phase structure of the system as a function of temperature and initial conditions. In section 5 we present the results of the numerical computations. Some conclusions are drawn in section 6.
## 2 Formulation of the model
We consider the $`O(N)`$ vector model with the Lagrangian
$$=\frac{1}{2}_\mu \varphi ^i^\mu \varphi ^i\frac{\lambda }{4N}(\varphi ^i\varphi ^iNv^2)^2$$
(2.1)
where $`\varphi ^i,i=1,..,N`$ are $`N`$ real scalar fields. The nonequilibrium state of the system is characterized by a classical expectation value which we take in the direction of $`\varphi _N`$. We split the field into its expectation value $`\varphi `$ and the quantum fluctuations $`\psi `$ via
$$\varphi ^i(𝐱,t)=\delta _N^i\sqrt{N}\varphi (t)+\psi ^i(𝐱,t).$$
(2.2)
In the large-$`N`$ limit one neglects, in the Lagrangian, all terms which are not of order $`N`$. In particular terms containing the fluctuation $`\psi _N`$ of the component $`\varphi _N`$ are at most of order $`\sqrt{N}`$ and are dropped, therefore. The fluctuations of the other components are identical, their summation produces factors $`N1=N(1+O(1/N))`$. In the broken symmetry case these are the Goldstone modes. Identifying all the fields $`\psi _1,..\psi _{N1}`$ as $`\psi `$ the leading order term in the Lagrangian then takes the form
$$=N\left(_\varphi +_\psi +_I\right),$$
(2.3)
with
$`_\varphi `$ $`=`$ $`{\displaystyle \frac{1}{2}}_\mu \varphi ^\mu \varphi {\displaystyle \frac{\lambda }{4}}\left(\varphi ^2v^2\right)^2,`$ (2.4)
$`_\psi `$ $`=`$ $`{\displaystyle \frac{1}{2}}_\mu \psi ^\mu \psi +{\displaystyle \frac{\lambda }{2}}v^2\psi ^2+{\displaystyle \frac{\lambda }{4}}(\psi ^2)^2,`$ (2.5)
$`_\mathrm{I}`$ $`=`$ $`{\displaystyle \frac{\lambda }{2}}\psi ^2\varphi ^2,`$ (2.6)
where $`\psi ^2`$ is to be identified with $`\psi ^i\psi ^i/N`$.
We decompose the fluctuating field into momentum eigenfunctions via
$$\psi (𝐱,t)=\frac{d^3k}{(2\pi )^32\omega _{k0}}\left[a_𝐤U_k(t)e^{i\mathrm{𝐤𝐱}}+a_𝐤^{}U_k^{}(t)e^{i\mathrm{𝐤𝐱}}\right],$$
(2.7)
with $`\omega _{k0}=\sqrt{m_0^2+k^2}`$. The mass $`m_0`$ will be specified below. This field decomposition defines a vacuum state as being annihilated by the operators $`a_𝐤`$.
The equations of motion for the field $`\varphi (t)`$ and of the fluctuations $`U_k(t)`$ have been derived in this formalism by various authors .
We include in the following the counterterms that we will need later in order to write the renormalized equations. The equation of motion for the field $`\varphi `$ becomes
$$\ddot{\varphi }(t)+\delta m^2\varphi (t)\lambda v^2\varphi (t)+(\lambda +\delta \lambda )\varphi (t)\left[\varphi ^2(t)+(t,T)\right]=0.$$
(2.8)
Here $`(t,T)`$ is the divergent fluctuation integral; it is given by the average of the fluctuation fields defined by the initial density matrix. For a thermal initial state of quanta with energy $`\omega _{k0}=\sqrt{k^2+m_0^2}`$ it is given by
$$(t,T)=\psi ^2(𝐱,t)=\frac{d^3k}{(2\pi )^32\omega _{k0}}\mathrm{coth}\frac{\beta \omega _{k0}}{2}|U_k(t)|^2.$$
(2.9)
The mode functions satisfy the equation:
$$\left[\frac{d^2}{dt^2}+\omega _k^2(t)\right]U_k(t)=0,$$
(2.10)
and the initial conditions
$$U_k(0)=1;\dot{U}_k(0)=i\omega _{k0}.$$
(2.11)
The time dependent frequency $`\omega _k(t)`$ is given by
$$\omega _k^2(t)=k^2+^2(t)$$
(2.12)
with the time dependent mass
$$^2(t)=\lambda v^2+\delta m^2+(\lambda +\delta \lambda )\left[\varphi ^2(t)+(t)\right].$$
(2.13)
Using this definition the classical equation of motion can be rewritten as
$$\ddot{\varphi }(t)+^2(t)\varphi (t)=0$$
(2.14)
which is the same equation as the one for $`U_k(t)`$ with $`k=0`$ (zero mode). Of course the initial conditions are different and $`\varphi (t)`$ is real.
As in our previous work we rewrite the mode equation in the form
$$\left[\frac{d^2}{dt^2}+\omega _{k0}^2\right]U_k(t)=𝒱(t)U_k(t),$$
(2.15)
whereby we have defined the time-dependent potential $`𝒱(t)=^2(t)^2(0)`$; we further identify $`m_0=(0)`$ as the “initial mass”.
The average of energy with respect to the initial density matrix is given by <sup>3</sup><sup>3</sup>3Note that twice the last term, with positive sign, is included in the fluctuation energy, since $`\omega _k^2(t)`$ contains $`(t,T)`$.
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\varphi }^2(t)+{\displaystyle \frac{1}{2}}(\lambda v^2+\delta m^2)\varphi ^2(t)+{\displaystyle \frac{\lambda +\delta \lambda }{4}}\varphi ^4(t)+\delta \mathrm{\Lambda }`$
$`+{\displaystyle \frac{d^3k}{(2\pi )^32\omega _{k0}}\mathrm{coth}\frac{\beta \omega _{k0}}{2}\left\{\frac{1}{2}|\dot{U}_k(t)|^2+\frac{1}{2}\omega _k^2(t)|U_k(t)|^2\right\}}`$
$`{\displaystyle \frac{\lambda +\delta \lambda }{4}}^2(t,T).`$
It is easy to check, using the equations of motion (2.14) and (2.10), that the energy is conserved. The energy density is the $`00`$ component of the energy-momentum tensor. The average of the energy momentum tensor for our system is diagonal, its space-space components define the pressure which is given by
$`p`$ $`=`$ $`\varphi ^2(t)+\delta \xi {\displaystyle \frac{d^2}{dt^2}}\left[\varphi ^2(t)+(t,T)\right]`$
$`+{\displaystyle \frac{d^3k}{(2\pi )^32\omega _{k0}}\mathrm{coth}\frac{\beta \omega _{k0}}{2}\left(\omega _{k0}^2+\frac{k^2}{3}\right)|U_k(t)|^2}.`$
$`\delta \xi `$ is the renormalization of the conformal coupling term $`\xi (g_{\mu \nu }^2_\mu _\nu )\varphi ^2`$, which has been used for the improved energy momentum tensor .
## 3 The renormalized equation of motion
The expressions for the time-dependent mass $`^2(t)`$, the energy density $`(t)`$ and the pressure are still undefined as they involve divergent integrals over the fluctuations. Our approach to regularization and renormalization has been presented previously . It is based on expanding the fluctuations $`U_k(t)`$ and subsequently the various integrals involving these fluctuations with respect to the time-dependent potential $`𝒱(t)`$. As this procedure has been presented elsewhere in detail we just give the outline, here.
The expansion of the fluctuations with respect to $`𝒱(t)`$ is given in Appendix A. We use this perturbative expansion in order to single out the divergent contributions in the fluctuation integral. One finds
$$(t)=I_1(m_0,T)I_3(m_0,T)\left[^2(t)^2(0)\right]+_{\mathrm{fin}}(t,T),$$
(3.1)
where the finite part of $`(t,T)`$ can be written as
$`_{\mathrm{fin}}(t,T)`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{4\omega _{k0}^3}\underset{0}{\overset{t}{}}𝑑t^{}\mathrm{cos}\left[2\omega _{k0}(tt^{})\right]\dot{𝒱}(t^{})\mathrm{coth}\frac{\beta \omega _{k0}}{2}}`$ (3.2)
$`+`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{2\omega _{k0}}\left[2\mathrm{R}\mathrm{e}f_k^{\overline{(2)}}(t)+|f_k^{\overline{(1)}}(t)|^2\right]\mathrm{coth}\frac{\beta \omega _{k0}}{2}},`$
and where the divergent integrals are defined as
$`I_1(m_0,T)`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{2\omega _{k0}}\left(1+\frac{2}{e^{\beta \omega _0}1}\right)}=I_1(m_0)+\mathrm{\Sigma }_1(m_0,T),`$ (3.3)
$`I_3(m_0,T)`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{4\omega _{k0}^3}\left(1+\frac{2}{e^{\beta \omega _0}1}\right)}=I_3(m_0)+\mathrm{\Sigma }_3(m_0,T).`$ (3.4)
The integrals $`I_k(m_0)`$ are the genuine divergences which appear in the renormalization at $`T=0`$. Their dimensionally regularized form is given by
$`I_3(m_0)`$ $`=`$ $`\left\{{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{4\omega _{k0}^3}}\right\}_{\mathrm{reg}}={\displaystyle \frac{1}{16\pi ^2}}\left\{{\displaystyle \frac{2}{ϵ}}+\mathrm{ln}{\displaystyle \frac{4\pi \mu ^2}{m_0^2}}\gamma \right\},`$ (3.5)
$`I_1(m_0)`$ $`=`$ $`\left\{{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{2\omega _{k0}}}\right\}_{\mathrm{reg}}={\displaystyle \frac{m_0^2}{16\pi ^2}}\left\{{\displaystyle \frac{2}{ϵ}}+\mathrm{ln}{\displaystyle \frac{4\pi \mu ^2}{m_0^2}}\gamma +1\right\}`$ (3.6)
$`=`$ $`m_0^2I_3(m_0){\displaystyle \frac{m_0^2}{16\pi ^2}}.`$
The additional temperature dependent terms $`\mathrm{\Sigma }_k(m_0,T)`$ are finite. They are defined as
$`\mathrm{\Sigma }_1(m_0,T)`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{\omega _{k0}\left(e^{\beta \omega _{k0}}1\right)}},`$ (3.7)
$`\mathrm{\Sigma }_3(m_0,T)`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{2\omega _{k0}^3\left(e^{\beta \omega _{k0}}1\right)}}.`$ (3.8)
It is convenient to include these finite terms into the definition of $`_{\mathrm{fin}}(t,T)`$. Then the time dependent mass takes the form
$$^2(t)=\lambda (\varphi ^2v^2)+\delta \lambda \varphi ^2+\delta m^2+(\lambda +\delta \lambda )\left[I_1(m_0)I_3(m_0)𝒱(t)+\stackrel{~}{}_{\mathrm{fin}}(t,T)\right],$$
(3.9)
with
$$\stackrel{~}{}_{\mathrm{fin}}(t,T)=\mathrm{\Sigma }_1(m_0,T)𝒱(t)\mathrm{\Sigma }_3(m_0,T)+_{\mathrm{fin}}(t,T).$$
(3.10)
The time dependent mass (3.9) contains both renormalization constants $`\delta m`$ and $`\delta \lambda `$. Furthermore, its definition by this equation is implicit, $`^2(t)`$ also appears on the right hand side of (3.9) in $`𝒱(t)`$.
We now have to fix the renormalization counterterms in such a way that the relation between the time dependent mass and $`\varphi (t)`$ becomes finite. An additional constraint derives from the requirement that the renormalization counterterms should not depend on the initial condition, but only on the parameters appearing in the Lagrangian, i.e., $`\lambda `$ and $`v`$ and renormalization conventions.
We first determine $`\delta \lambda `$ by considering the difference
$`𝒱(t)`$ $`=`$ $`^2(t)^2(0)`$
$`=`$ $`(\lambda +\delta \lambda )\left[\varphi ^2(t)\varphi ^2(0)I_3(m_0)𝒱(t)+\stackrel{~}{}_{\mathrm{fin}}(t,T)\stackrel{~}{}_{\mathrm{fin}}(0,T)\right].`$
The divergent parts depend on the initial mass $`m_0`$. We have to replace this by a renormalization scale independent of the initial conditions. In Ref. we had chosen the scale $`m`$, where $`m`$ was the mass parameter appearing in the Lagrangian. Here the analogous mass squared would be $`m^2=\lambda v^2`$ and so $`m`$ would be imaginary. We therefore choose another scale $`m_1`$ which we do not specify here. In the numerical computations we have used the physical Higgs mass $`m_1^2=m_H^2=2\lambda v^2`$.
We rewrite the implicit equation for $`𝒱(t)`$ as
$`𝒱(t)\left[1+(\lambda +\delta \lambda )I_3(m_1)\right]`$ $`=`$ $`(\lambda +\delta \lambda )\{\varphi ^2(t)\varphi ^2(0)[I_3(m_0)I_3(m_1)]𝒱(t)`$ (3.12)
$`+\stackrel{~}{}_{\mathrm{fin}}(t,T)\stackrel{~}{}_{\mathrm{fin}}(0,T)\}`$
and require
$$\frac{\lambda +\delta \lambda }{1+(\lambda +\delta \lambda )I_3(m_1)}=\lambda .$$
(3.13)
Solving with respect to $`\delta \lambda `$ we find
$$\delta \lambda =\frac{\lambda ^2I_3(m_1)}{1\lambda I_3(m_1)}.$$
(3.14)
Inserting this relation into (3.12) we find
$$𝒱(t)=\lambda 𝒞\left[\varphi ^2(t)\varphi ^2(0)+\stackrel{~}{}_{\mathrm{fin}}(t,T)\stackrel{~}{}_{\mathrm{fin}}(0,T)\right].$$
(3.15)
with
$$𝒞=\frac{1}{1+\lambda \left[I_3(m_0)I_3(m_1)\right]}=\frac{1}{1+{\displaystyle \frac{\lambda }{16\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{m_1^2}{m_0^2}}\right)}.$$
(3.16)
Eq. (3.15) is a finite relation for the potential $`𝒱(t)`$ since the difference $`[I_3(m_0)I_3(m_1)]`$ is finite. Going back to Eq. (3.10) we realize that $`\stackrel{~}{}_{\mathrm{fin}}`$ on the right hand side contains itself a term proportional to $`𝒱(t)`$. Taking account of this term we rewrite $`𝒱(t)`$ in terms of $`_{\mathrm{fin}}`$ as
$$𝒱(t)=\lambda 𝒞_T\left[\varphi ^2(t)\varphi ^2(0)+_{\mathrm{fin}}(t,T)\right]$$
(3.17)
with
$$𝒞_T=\frac{1}{1+{\displaystyle \frac{\lambda }{16\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{m_1^2}{m_0^2}}\right)+\lambda \mathrm{\Sigma }_3(m_0,T)}.$$
(3.18)
Recall that $`_{\mathrm{fin}}(t)`$ is the mode integral of second order in $`𝒱(t)`$ and vanishes at $`t=0`$.
We now go back to equation (3.9) which we take at the initial time $`t=0`$:
$$m_0^2^2(0)=\lambda [\varphi ^2(0)v^2]+\delta \lambda \varphi ^2(0)+\delta m^2+(\lambda +\delta \lambda )\left[I_1(m_0)+\stackrel{~}{}_{\mathrm{fin}}(0,T)\right].$$
(3.19)
This is an implicit relation between $`m_0`$ and $`\varphi (0)`$ which, however, contains still the infinite quantities $`\delta \lambda `$, $`\delta m`$ and $`I_1(m_0)`$. Using Eq. (3.6) we can rewrite Eq. (3.19) as
$$m_0^2=\left(\lambda v^2+\delta m^2\right)+(\lambda +\delta \lambda )\left[\varphi ^2(0)m_0^2I_3(m_0)\frac{m_0^2}{16\pi ^2}+\stackrel{~}{}_{\mathrm{fin}}(0,T)\right].$$
(3.20)
As renormalization condition we require $`m_0`$ to vanish, for temperature $`T=0`$, at the minimum of the potential $`\varphi =v`$, as it is the case on the tree level. We note that $`m_0^2=0`$ is not the curvature of the tree level potential at $`\varphi =v`$ which is $`m_H^2=2\lambda v^2`$. It is the mass of the fluctuations at $`\varphi =v`$ in the large-$`N`$ approximation. For $`T=0`$ we have $`\stackrel{~}{}_{\mathrm{fin}}(t=0,T=0)=\mathrm{\Sigma }_1(m_0,T=0)=0`$. Setting $`m_0=0`$, $`\varphi (0)=v`$ in the gap equation (3.20) we get immediately
$$\delta m^2=\delta \lambda v^2=\frac{\lambda ^2v^2I_3(m_1)}{1\lambda I_3(m_1)}.$$
(3.21)
Inserting this into Eq. (3.20) we obtain the renormalized gap equation
$$m_0^2=\lambda 𝒞\left[\varphi ^2(0)v^2\frac{m_0^2}{16\pi ^2}+\mathrm{\Sigma }_1(m_0,T)\right].$$
(3.22)
For the numerical computation it is easier to choose some $`m_0^20`$ and to use the gap equation solved for $`\varphi ^2(0)`$:
$$\varphi ^2(0)=\frac{m_0^2}{\lambda }+v^2+\frac{m_0^2}{16\pi ^2}\left(1+\mathrm{ln}\frac{m_1^2}{m_0^2}\right)\mathrm{\Sigma }_1(m_0,T).$$
(3.23)
For $`t>0`$ the renormalized relation for the mass squared $`^2(t)`$ we find, using Eqns. (3.15) and (3.22), is
$$^2(t)=m_0^2+𝒱(t)=\lambda 𝒞\left[\varphi ^2(t)v^2\frac{m_0^2}{16\pi ^2}+\stackrel{~}{}_{\mathrm{fin}}(t,T)\right].$$
(3.24)
Having thus obtained a finite relation between $`\varphi (t)`$ and $`(t)`$ the equations of motion for the classical field $`\varphi (t)`$ and for the modes $`U_k(t)`$ are well-defined and finite.
The way in which we have renormalized has made the cutoff disappear. This was possible only to the extent that we could safely neglect corrections of order $`ϵ`$ in the evaluation of the divergent integrals. One way of achieving this is to take the limit $`ϵ0`$. This implies for the bare coupling $`\lambda _0`$
$$\lambda _0=\underset{ϵ0}{lim}\frac{\lambda }{1{\displaystyle \frac{\lambda }{16\pi ^2}}{\displaystyle \frac{2}{ϵ}}}=0^{},$$
(3.25)
so this is the case of “negative bare coupling” as discussed in . One can leave the cutoff finite, however, as long as the masses and momenta are much smaller than the scale of the Landau ghost, $`m_x=m_1^2\mathrm{exp}(8\pi ^2/\lambda )`$. This will be case here. This is not related to a pragmatic momentum cutoff that we apply to the convergent integrals of the finite part.
While we have found here the gap equation as a self-consistency condition, it can also be derived from a potential (free energy) which here takes the form
$`V(m_0^2,\mathrm{\Phi }^2,T)`$ $`=`$ $`{\displaystyle \frac{m_0^2}{2}}\left\{\varphi ^2v^2{\displaystyle \frac{m_0^2}{2\lambda }}+{\displaystyle \frac{m_0^2}{32\pi ^2}}\left[\mathrm{ln}\left({\displaystyle \frac{m_0^2}{m_1^2}}\right){\displaystyle \frac{3}{2}}\right]\right\}`$ (3.26)
$`+{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{\beta }\mathrm{ln}\left[1\mathrm{exp}(\beta \omega _0)\right]}.`$
The gap equation then follows from the condition
$$\frac{V(m_0^2,\varphi ^2,T)}{m_0^2}=0.$$
(3.27)
It should be mentioned here that the gap equation has two solutions, one of which lies above the scale of the Landau ghost, $`m_x=m_1\mathrm{exp}(8\pi ^2/\lambda ^2)`$. In the sense that we consider here the model as giving rise to a low energy effective theory we discard this high mass solution, and its discussion. The solution we consider is the low energy one which is of order $`\sqrt{\lambda }v`$ or $`m_1`$.
The energy density is given by
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\varphi }^2(t)+{\displaystyle \frac{1}{4}}\left(\lambda +\delta \lambda \right)\left(\varphi ^2v^2\right)^2+\delta \mathrm{\Lambda }`$ (3.28)
$`+_{\mathrm{fl}}(t,T){\displaystyle \frac{\lambda +\delta \lambda }{4}}^2(t,T).`$
Here we have used already that $`\delta m^2=\delta \lambda v^2`$, and part of the “cosmological constant” counterterm $`\delta \mathrm{\Lambda }`$ is included in $`\delta \lambda v^4/4`$. The fluctuation energy is given by
$$_{\mathrm{fl}}(t,T)=\frac{d^3k}{(2\pi )^32\omega _{k0}}\mathrm{coth}\frac{\beta \omega _{k0}}{2}\left\{\frac{1}{2}|\dot{U}_k(t)|^2+\frac{1}{2}\omega _k^2(t)|U_k(t)|^2\right\}.$$
(3.29)
We again split off the temperature-dependent contribution via
$$_{\mathrm{fl}}(t,T)=_{\mathrm{fl}}(t,0)+\mathrm{\Delta }_{\mathrm{fl}}(t,T),$$
(3.30)
where the second term on the right hand side
$$\mathrm{\Delta }_{\mathrm{fl}}(t,T)=\frac{d^3k}{(2\pi )^32\omega _{k0}}\frac{2}{e^{\beta \omega _{k0}}1}\left\{\frac{1}{2}|\dot{U}_k(t)|^2+\frac{1}{2}\omega _k^2(t)|U_k(t)|^2\right\},$$
(3.31)
is finite. The divergences of the first term are given by the decomposition
$$_{\mathrm{fl}}(t,0)=I_1(m_0)+\frac{1}{2}𝒱(t)I_1(m_0)\frac{1}{4}𝒱^2(t)I_3(m_0)+_{\mathrm{fl},\mathrm{fin}}(t,0)$$
(3.32)
with
$$_{\mathrm{fl},\mathrm{fin}}(t,0)=\frac{1}{2}\frac{d^3k}{(2\pi )^32\omega _{k0}}\left\{\frac{1}{2}|\dot{f}_k^{\overline{(1)}}|^2+\frac{𝒱(t)}{2}\left[2\mathrm{R}\mathrm{e}f_k^{\overline{(1)}}+|f_k^{\overline{(1)}}|^2\right]+\frac{𝒱^2(t)}{8\omega _{k0}^2}\right\}.$$
(3.33)
We denote the sum of $`_{\mathrm{fl},\mathrm{fin}}(t,0)`$ and $`\mathrm{\Delta }_{\mathrm{fl}}(t,T)`$ finite contributions as $`_{\mathrm{fl},\mathrm{fin}}(t,T)`$. The expression for the energy then takes the form
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\varphi }^2+{\displaystyle \frac{\lambda +\delta \lambda }{4}}\left(\varphi ^2v^2\right)^2+_{\mathrm{fl},\mathrm{fin}}(t,T)+I_1(m_0)+{\displaystyle \frac{1}{2}}𝒱(t)I_1(m_0){\displaystyle \frac{1}{4}}𝒱^2(t)I_3(m_0)`$ (3.34)
$`{\displaystyle \frac{\lambda +\delta \lambda }{4}}^2(t,T)+\delta \mathrm{\Lambda }.`$
In addition to the divergences arising from $`_{\mathrm{fl}}(t,T)`$ we have to take into consideration those of $`^2(t,T)`$ which we have analyzed above. If all divergences and the renormalization constant $`\delta \lambda `$ are inserted, the expression turns out to be finite, i.e., the remaining counterterm $`\delta \mathrm{\Lambda }`$ is needed only for a finite renormalization. We require the energy to vanish at $`T=0`$ for $`\varphi (t)v`$, which implies $`m_0=0`$. Then $`\delta \mathrm{\Lambda }=0`$. There remains a finite constant dependent on the initial condition
$$\mathrm{\Delta }\mathrm{\Lambda }=\frac{m_0^4}{128\pi ^2}\left(1+\frac{2\lambda 𝒞}{16\pi ^2}\right)$$
(3.35)
and the energy is given by
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\varphi }^2+{\displaystyle \frac{\lambda }{4}}𝒞(\varphi ^2v^2)^2+{\displaystyle \frac{1}{2}}\mathrm{\Delta }m^2(\varphi ^2v^2)`$
$`+_{\mathrm{fl},\mathrm{fin}}(t,T){\displaystyle \frac{\lambda }{4}}𝒞\stackrel{~}{}_{\mathrm{fin}}^2(t,T)+\mathrm{\Delta }\mathrm{\Lambda }.`$
Here $`\mathrm{\Delta }m^2`$ is given by
$$\mathrm{\Delta }m^2=\lambda 𝒞\frac{m_0^2}{16\pi ^2}.$$
(3.37)
We write the pressure in the form
$$p=\dot{\varphi }^2(t)+p_{\mathrm{fl}}(t,T)+\delta \xi \frac{d^2}{dt^2}\left[\varphi ^2(t)+(t,T)\right].$$
(3.38)
The renormalization does not differ form the case of unbroken symmetry discussed in Ref. and is not presented again. We find
$$\delta \xi =\frac{\lambda x}{6(1\lambda x)}=\frac{\lambda I_3(m)}{6(1\lambda I_3(m))}.$$
(3.39)
The final result for the renormalized pressure reads
$$p=\dot{\varphi }^2(t)+p_{\mathrm{fl},\mathrm{fin}}(t,T)\frac{m_0^4}{96\pi ^2}\frac{m_0^2}{48\pi ^2}𝒱(t)\frac{1}{96\pi ^2}\left[\mathrm{ln}\left(\frac{m_1^2}{m_0^2}\right)+2\right]\ddot{𝒱}(t)$$
(3.40)
with
$`p_{\mathrm{fl},\mathrm{fin}}(t,0)`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d^3k}{(2\pi )^32\omega _{k0}}}\{(\omega _{k0}^2+{\displaystyle \frac{\stackrel{}{k}^2}{3}})[2\mathrm{R}\mathrm{e}f_k^{\overline{(2)}}(t)+|f_k^{\overline{(1)}}(t)|^2]`$ (3.42)
$`+\left({\displaystyle \frac{1}{6\omega _{k0}^2}}{\displaystyle \frac{m_0^2}{24\omega _{k0}^4}}\right){\displaystyle \underset{0}{\overset{t}{}}}𝑑t^{}\mathrm{cos}2\omega _k^0(tt^{})\stackrel{\mathrm{}}{𝒱}(t^{})`$
$`+\left({\displaystyle \frac{1}{12\omega _{k0}^2}}+{\displaystyle \frac{m_0^2}{24\omega _{k0}^4}}\right)\mathrm{cos}(2\omega _{k0}t)\ddot{𝒱}(0)`$
$`+|\dot{f}_k^{\overline{(1)}}(t)|^22\mathrm{R}\mathrm{e}[i\omega _{k0}\dot{f}_k^{\overline{(1)}}(t)+i\omega _{k0}f_k^{\overline{(1)}}(t)f_k^{\overline{(1)}}(t)]\}.`$
## 4 Analysis of the gap equation and of the phase structure
The dynamical evolution of the nonequilibrium system depends on two parameters, the temperature $`T`$ and the initial amplitude of the classical field $`\varphi (0)=\varphi _0`$ which in analogy with thermal equilibrium systems can be considered as an external parameter. There are two regions from which we can start the system, which we will call regions II and III. There is, in addition, one region into which the system can evolve when one considers $`\varphi (t=\mathrm{})=\varphi _{\mathrm{}}`$ and not $`\varphi _0`$ as the external parameter. We call it region I. In this section we will characterize these regions and describe the dynamical evolution as to be expected from the analysis at $`T=0`$. This analysis is based on certain empirical results that, though unproven, seem to be at least almost exact. We will generalize these results to finite temperature in a plausible heuristic way, to be confirmed by the numerical computations. We think that the way in which we generalize these results will give a further clue to understanding them.
### 4.1 Region I, $`m_0^2<0`$
The gap equation requires $`m_0^2`$ to be positive. The point where $`m_0^2=0`$ marks an initial condition that leads to a solution $`\varphi =\mathrm{const}.`$, if $`\dot{\varphi }_0=0`$ as we will assume in the following. For $`T=0`$ this stationary amplitude is $`\varphi =v`$. For $`T>0`$ we can easily find this amplitude as well. Indeed for $`m_0=0`$ the integral $`\mathrm{\Sigma }_1(m_0,T)`$ is given by its value for massless quanta, i.e.,
$$\mathrm{\Sigma }_1(0,T)=\frac{T^2}{12},$$
(4.1)
therefore
$$\varphi _1^2(T)=\varphi _0^2|_{m_0=0}=v^2\frac{T^2}{12}.$$
(4.2)
For $`\varphi _0<\varphi _1(T)`$ the gap equation has no real solution $`m_0`$. The region below the boundary (4.2) is region I.
If nevertheless one wants to start the system with $`\varphi _0`$ in region I one faces the problem that in this region the gap equation requires $`m_0^2`$ to be negative. Then the low-momentum modes with $`k^2<m_0^2`$ have imaginary frequencies. So from an orthodox point of view (to which we adhere here) the system cannot be quantized properly. One may avoid this problem by redefining the dispersion relation for the initial frequencies via $`\omega _{k0}^2=k^2+|m_0^2|`$ in this region. Of course at $`t>0`$ $`^2(t)`$ will be negative so the “mass squared” changes sign at $`T=0`$, a situation called “quench” in analogy by a similar transition form a stable to an unstable state by a sudden drop of temperature or inversion of a magnet field. On the other hand the amplitude $`\varphi (t)`$ can reach this region at late times, but then it is in a quantum state different form the ones we use as initial states.
### 4.2 Region II: $`m_0^2>0,m_{\mathrm{}}^2=0`$
We now assume $`\varphi _0`$ is started above the boundary value (4.2). If $`\varphi _0`$ is not too large the system may, at $`t>0`$, enter a region where $`^2(t)<0`$, i.e., region I. Then the quantum fluctuations with momenta $`k^2<^2(t)`$ will increase exponentially, signalling instability. This causes $`^2(t)`$ to increase so that it is driven back to a value $`^2(t)>0`$. If the initial amplitude $`\varphi _0`$ is sufficiently small this forth-and-back reaction will lead $`^2(t)`$ to stabilize at $`_{\mathrm{}}^2=0`$. So at late times $`\stackrel{~}{}_{\mathrm{fin}}(t,T)`$ is determined by quantum modes $`U_k(t)`$ that oscillate with time-independent frequencies $`\omega _{\mathrm{}}=k`$, it becomes stationary as well and will be positive. Therefore $`\varphi (t)`$ stabilizes at some value
$$\varphi _{\mathrm{}}^2=v^2\stackrel{~}{}_{\mathrm{fin}}(\mathrm{},T)<v^2\stackrel{~}{}_{\mathrm{fin}}(0,T).$$
(4.3)
This is entirely analogous to the behavior found at $`T=0`$ . We call the region of initial values $`\varphi _0`$ leading to this late-time behavior region II.
The stabilization by back-reaction onto the fluctuations obtained in the large-$`N`$ approximation is not present in the one-loop approximation. In this approximation, once $`\varphi (t)`$ dips into the unstable region $`\varphi (t)<v/\sqrt{3}`$, the mass squared of the fluctuations becomes negative and the low momentum modes evolve exponentially. The effective mass of the classical field increases exponentially as well, and continues to do so, but the mass squared of the fluctuations stays negative. The amplitude $`\varphi (t)`$ is driven towards zero. Nevertheless the classical energy continues to increase, as the field oscillates faster and faster, this energy being extracted from the energy of the quantum fluctuations. Obviously this signals the instability of the quantum vacuum, as already apparent from the fact that the effective potential is complex in this region. We will illustrate this by a numerical example, to be presented in the next section.
At $`T=0`$ the final value $`\varphi _{\mathrm{}}`$ was found to be related to the initial value $`\varphi _0`$ by an empirical relation
$$\varphi _{\mathrm{}}^2=\sqrt{\varphi _0^2(2v^2\varphi _0^2)}T=0.$$
(4.4)
It is not obvious how to generalize this relation to finite temperature. It was remarked in Ref. that the relation only depends on the initial, purely classical, energy, which is given by $`E=\lambda (\varphi ^2v^2)^2/4`$. Obviously it satisfies the constraints that $`\varphi _{\mathrm{}}^2=v^2`$ if $`\varphi _0=v^2`$, and that $`\varphi _{\mathrm{}}=0`$ if classically the system can reach the maximum of the potential; this happens at $`\varphi _0^2=\varphi _2^2(T=0)=2v^2`$. So Eq. (4.4) seems to be related to energy considerations. We further observe that the classical turning point is at $`\overline{\varphi }_0^2=2v^2\varphi _0^2`$ so that one may write Eq. (4.4) as the geometric mean
$$\varphi _{\mathrm{}}^2=\sqrt{\varphi _0^2\overline{\varphi }_0^2}.$$
(4.5)
This form turns indeed out to lead to the correct generalization for finite temperature.
Obviously the relation is characterized by the motion at early times when the quantum fluctuations have not yet evolved. When discussing renormalization we have made an expansion with respect to the “potential” $`𝒱(t)`$ which vanishes at $`t=0`$. So the same expansion can be used to study the early time behavior. In the energy the coefficients of the terms of first and second order in $`𝒱`$ have been absorbed into renormalization constants. However, the thermal fluctuations are not absorbed in this way and will add to the classical terms in an early time expansion. These appear in the energy, see Eq.(3.34), via
$$\mathrm{\Delta }_{\mathrm{fl}}(t,T)=\mathrm{\Sigma }_1(m_0)+\frac{1}{2}𝒱(t)\mathrm{\Sigma }_1(m_0)\frac{1}{4}𝒱^2(t)\mathrm{\Sigma }_3(m_0)+O(𝒱^3)$$
(4.6)
as a part of $`_{\mathrm{fl},\mathrm{fin}}(t,T)`$ and via Eq. (3.10) in $`\stackrel{~}{}_{\mathrm{fin}}(t,T)`$. Taking these expansions into account the energy can be written in the form
$$E\frac{\lambda }{4}𝒞\left[a\varphi ^4+\stackrel{~}{a}\varphi _0^4+b\varphi ^2+\stackrel{~}{b}\varphi _0^2+c\varphi ^2\varphi _0^2\right]+\mathrm{const}.$$
(4.7)
up to terms of order $`𝒱^3`$. We need the coefficients
$`a`$ $`=`$ $`1\lambda 𝒞_T\mathrm{\Sigma }_3(m_0,T),`$ (4.8)
$`b`$ $`=`$ $`2\left[v^2\mathrm{\Sigma }_1(m_0,T)\right],`$ (4.9)
$`c`$ $`=`$ $`\lambda 𝒞𝒞_T\mathrm{\Sigma }_3(m_0,T).`$ (4.10)
The classical turning point is given by
$$\overline{\varphi }_0^2=\frac{b+(a+c)\varphi _0^2}{a}=\frac{1}{1\lambda 𝒞_T\mathrm{\Sigma }_3}\left[2v^2\mathrm{\Sigma }_1(1+\lambda 𝒞_T\mathrm{\Sigma }_3)\varphi _0^2\right],$$
(4.11)
so that we are led to suppose
$$\varphi _{\mathrm{}}^2(T)=\sqrt{\frac{1}{1\lambda 𝒞_T\mathrm{\Sigma }_3}}\sqrt{\varphi _0^2\left[2v^22\mathrm{\Sigma }_1(1+\lambda 𝒞_T\mathrm{\Sigma }_3)\varphi _0^2\right]}.$$
(4.12)
We find indeed (see below) that this relation is very well fulfilled numerically. According to this formula the region II is limited by the requirement that the expression in the square root be positive, so the boundary between region II and the new region III is given by
$$\varphi _2^2=2\frac{v^2\mathrm{\Sigma }_1(m_0,T)}{1+\lambda 𝒞_T\mathrm{\Sigma }_3(m_0,T)}.$$
(4.13)
We note that the relation is implicit, the value of $`m_0`$ that appears on the right hand side is related to $`\varphi _2^2`$ on the left hand side by the gap equation.
### 4.3 Region III, $`\varphi _{\mathrm{}}=0`$ and $`_{\mathrm{}}^2>0`$
If the value $`\varphi _0`$ becomes larger than $`\varphi _2`$ the stationary state with constant $`\varphi `$ and vanishing mass $`^2(t)`$ is no longer attained, and the system reaches another asymptotic regime where $`^2(\mathrm{})0`$ whereas $`\varphi (t)0`$. This regime is similar to the one that describes the late time behavior for the unbroken symmetry case. We call the region of initial values $`\varphi _0`$ that leads to such a behavior region III.
There are two phenomena that characterize the transition to this region. On the one hand the stabilization of the system is taken over by the phenomenon of parametric resonance. On the other hand the system has enough energy so that $`\varphi (t)`$ can move over the maximum of the potential at $`\varphi =0`$, and indeed will oscillate around $`\varphi =0`$. Accordingly the threshold value of $`\varphi _0`$ at which these two phenomena set in can be characterized by two - a priori unrelated - criteria. Both rely on plausible assumptions, which at $`T=0`$ lead to the same prediction for the critical value of $`\varphi _0`$.
The criterion based on the energy consideration has been presented in the previous subsection, we now describe the criterion supplied by the phenomenon of parametric resonance. For the case of unbroken symmetry it was found at zero and finite temperature , that the late time behavior is described by an empirical sum rule which relates $`_{\mathrm{}}^2`$ to the initial amplitude. For $`T=0`$ an analogous sum rule was found to hold for the case of spontaneously broken symmetry as well . It is given by
$$\mu _{\mathrm{}}^2=1+\frac{\eta _0^2}{2}.$$
(4.14)
Here $`\mu `$ and $`\eta `$ are normalized in such a way that the classical equation of motion at early times, i.e., in the parametric resonance regime without back reaction, reads
$$\eta ^{\prime \prime }\eta +\eta ^3=0,$$
(4.15)
where the prime denotes a derivative with respect to $`\tau =\alpha t`$ and where $`\eta =\beta \varphi `$, also $`\mu =/\alpha `$. With $`\eta (\tau )`$ a solution of Eq. (4.15) the mode equation becomes a Lamé equation. The sum rule implies , that the frequencies $`\omega ^2(t)=^2(t)+k^2`$ are shifted outside the parametric resonance band of the Lamé equation. Though there is no rigorous derivation for the sum rule, it accordingly seems related to the parametric resonance phenomenon.
As the shift of the frequencies outside the parametric resonance region must have happened at the end of the phase where the evolution of the system is described by parametric resonance, we will again consider the initial classical evolution. Again, in addition to the classical terms we have to take into account the terms due to the thermal fluctuations. In terms of the parameters introduced in the previous section the equation of motion is given by
$`\ddot{\varphi }+\lambda 𝒞a\varphi ^3+{\displaystyle \frac{\lambda }{2}}𝒞(b+c\varphi _0^2)\varphi =0.`$ (4.16)
Comparing to the normalized equation (4.15) we determine the factors $`\alpha `$ and $`\beta `$ to be
$`\alpha `$ $`=`$ $`\sqrt{{\displaystyle \frac{\lambda 𝒞}{2}}}\sqrt{b+c\varphi _0^2},`$ (4.17)
$`\beta `$ $`=`$ $`\sqrt{{\displaystyle \frac{2a}{b+c\varphi _0^2}}},`$ (4.18)
so that the asymptotic mass is given by
$`_{\mathrm{}}^2`$ $`=`$ $`\alpha ^2(1+{\displaystyle \frac{1}{2}}\beta ^2\varphi _0^2)`$
$`=`$ $`\lambda 𝒞\left\{v^2+\mathrm{\Sigma }_1(m_0,T)+{\displaystyle \frac{1}{2}}\left[1+\lambda 𝒞_T\mathrm{\Sigma }_3(m_0,T)\right]\varphi _0^2\right\}.`$
Again $`\varphi _0`$ and $`m_0`$ are related by the gap equation. At the transition from region II to region III the asymptotic mass vanishes. It is easily seen that this criterion leads to an identical equation for the boundary, i.e., Eq. (4.13).
The field amplitude decreases to zero at late times, in this regime. So the symmetry is restored dynamically at high excitation characterized by a high value of $`\varphi _0`$.
At the critical temperature $`T=\sqrt{12}v`$ both boundaries $`\varphi _1(T)`$ and $`\varphi _2(T)`$ become zero. Above $`T_C`$ the behavior of the system is the same as for region III, for all initial values of $`\varphi _0`$. While at the border between region I and II there was a lowest value for $`\varphi _0`$ for obtaining real solutions of the gap equation, now there is a lowest value of $`m_0`$, the one for which $`\varphi _0=0`$. It is obtained by solving the gap equation for $`\varphi _0=0`$ and agrees with the thermodynamical equilibrium value $`m_\beta `$ at that temperature, as defined, e.g., in Eq. (3.38) of Ref. . Of course with $`\varphi _0=0`$ the system remains static.
Having defined the three regions by the two boundaries (4.2) and (4.13) we present, in Fig. 1, a phase diagram in the $`\varphi _0^2T`$ plane. Fig. 2 shows the phase diagram in the $`m_0^2T`$ plane, displaying, above $`T_C`$, the region $`m_0<m_\beta `$ which is excluded as an initial condition. We have to stress that the boundary between regions II and III relies on an empirical relation.
The symmetry restoration above a critical temperature is expected naively. However, if the temperature becomes nonperturbatively large, $`T\sqrt{12}m_1\mathrm{exp}(8\pi ^2/\lambda )`$, the gap equation does not have solutions any longer. Then the free energy attains its maximum at the boundary $`m_0=0`$ and the $`O(N)`$ symmetry is again broken . This phenomenon of “symmetry non-restoration”, as well as the existence of the second solution of the gap equation above $`m_x=m_1\mathrm{exp}(8\pi ^2/\lambda )`$ will not be discussed here, as it is not part of the low energy effective theory.
## 5 Numerical Results
We have discussed already in the previous section the type of nonequilibrium behavior to be expected in the different regions of phase space. The numerical results follow these expectations. We have chosen generally the parameters $`v=1`$ and $`\lambda =1`$. We present results for the various regions in the $`T,\varphi _0`$ plane. The critical temperature is $`2\sqrt{3}=3.464`$. We choose the temperatures between $`T=1`$ and $`4`$, the latter one being above the phase transition. The numerical method has been described in . We just recall that all the integrals computed numerically are finite, so cutting off the momentum integration at some reasonable value is unrelated to cutoffs used for renormalization.
We first consider initial conditions in region II. The expectation value of $`\varphi `$, shown in Fig. 3 becomes constant and different from zero as $`t\mathrm{}`$. This signals spontaneous breakdown of the $`O(N)`$ symmetry. As displayed in Fig. 4 the mass $`^2(t)`$ vanishes as $`t\mathrm{}`$, as expected form the Goldstone theorem. The momentum distribution of the quantum fluctuations peaks at $`k=0`$ as $`|U_k(t)|^2k^2`$, leading to long-range correlations, a phenomenon called “dynamical Bose-Einstein condensation” in Ref. and investigated further, for finite volume, in Ref.. We show an example of the momentum distribution in Fig. 5, but we have not studied the phenomenon in detail.
The relation between the asymptotic value as $`t\mathrm{}`$ for $`\varphi (t)`$ and the initial amplitude $`\varphi _0`$ is displayed in Figs. 6 to 8, for $`T=1,2.5`$ and $`3`$. We compare the data with our generalization (4.12) of the empirical formula (4.4) given in Ref. . The data are obtained by averaging over the second half of the time interval. The agreement is excellent, except at the phase boundary where the averaging converges slowly.
As an illustration of the behavior of the system in the unstable region in the one-loop approximation we show, in Fig. 9, the evolution of the field amplitude, and, in Fig. 10, the exponential behavior of the fluctuation integral and of the effective mass squared $`^2(t)`$ of the classical field.
The behavior of the system in region III is displayed in Figs.11 and 12. The amplitude $`\varphi (t)`$ is seen to decrease to zero. The decrease is powerlike, not exponential, a phenomenon called anomalous relaxation in Ref. . Fig. 12 shows the squared mass $`^2(t)`$ which is seen to converge to an asymptotic value $`_{\mathrm{}}^2`$. The sum rule for this asymptotic value, Eq. (4.12), is compared to the data in Fig. 13 for $`T=1.5,2.5`$ and $`4`$. The agreement is again excellent.
We have not presented the results for the pressure and the ratio of pressure and energy which varies between $`0`$ for a nonrelativisticœ\[Dœ\[Dœ\[Dœ\[Distic and $`1/3`$ for an ultrarelativistic ensemble. Here these are dominated, already at $`T=1`$, by the purely thermal contributions, so that the fluctuations generated by the motion of the field $`\varphi (t)`$ are relatively unimportant.
## 6 Conclusions and Outlook
The dynamical exploration of the quantum states of the $`O(N)`$ $`\lambda \mathrm{\Phi }^4`$ theory in the limit $`N\mathrm{}`$ has been extended here to finite temperature. We have performed numerical simulations with various initial fields $`\varphi _0=\varphi (0)`$ and initial masses $`m_0`$ related by the gap equation, and for various temperatures $`T`$. Depending on the initial conditions we find, in analogy to computations at zero temperature , final states with restored $`O(N)`$ symmetry and final states for which the symmetry is spontaneously broken. The resulting phase diagrams resemble typical phase diagrams of thermodynamical systems, with the temperature and an external variable as parameters. Instead of, e.g., the magnetic field or the pressure we have here the initial value $`\varphi _0`$ as external parameter. While the initial states are thermal states, the final states are not.
We have generalized two empirical formulae, the relation between the initial and asymptotic field amplitudes in region II, and the formula for the asymptotic value of $`^2(t)`$ in region III to finite temperature, extending the plausibility arguments given in . While we have not been able, either, to derive these formulae, the way of generalizing them may give some clue for such a derivation. Both relations are linked as they give the same formula for the boundary between regions II and III, though the arguments for their heuristic derivation are seemingly different. Furthermore, it is clear that both of them are based on the early time behavior. Obviously the fluctuations have to be included up to order $`𝒱^2(t)`$ in a perturbative expansion. At $`T=0`$ these terms are essentially absorbed into renormalization constants, so that the purely classical behavior prevails. One may also formulate the modifications at finite temperature in terms of temperature dependent masses and couplings. It is the rôle of the large-$`N`$ quantum back reaction to transmit the early time behavior into the late time one.
Unfortunately there are many interesting models for which the large-$`N`$ approach is not possible or not adequate. The one-loop approximation, on the other hand, can be applied in general. However, it shows features that seem to make it obsolete for describing nonequilibrium phenomena, especially for theories with spontaneous symmetry breaking. As an illustration we have shown the typical behavior of a spontaneously broken $`\lambda \mathrm{\Phi }^4`$ model in the one-loop approximation. The system does not reach a stationary state at late times: the effective mass of the classical field diverges exponentially, while the effective mass of the quantum fluctuations is and stays negative. This is due to the lack of the quantum back reaction onto the fluctuations. The fact that one finds such a pathological behavior may, however, indicate the correct physics and is not necessarily a consequence of an inadequate approximation. It is known that the system is indeed unstable for spatially constant static fields, it is an instability with respect to formation of domains . For space dependent fields like minimal bubble configurations the one-loop approximation to the effective action does not display any unplausible features , though the effective potential is complex in the unstable region. So it is not clear whether the “taming” of the instability introduced by the large-$`N`$ approximation necessarily improves the understanding of the physics.
In this situation it is certainly very important to develop new approaches to the evolution of quantum systems for theories with spontaneously broken symmetry . There are indications in a large-$`N`$ quantum mechanical system that the large-$`N`$ limit may be misleading, as the next-to-leading corrections become large especially at late times. It is not clear, however, what the impact of these results on quantum field theory will be. One of the problems is that, in contrast to the large-$`N`$ and one loop approximations, alternative wave functionals pose problems with renormalization . This is not only a technical problem. It is connected (trivially) to the fact that the higher the dimension of space, the more the ultraviolet behavior of the system will be important.
We think nevertheless, that a good understanding of the leading order approximation may improve the understanding of the corrections. That these become large at late times is not too surprising, it is therefore even more important to realize (once more) that the late-time behavior is related to the early-time behavior, which will therefore set the initial conditions for other approximations as well. This should apply to the phase structure as well.
## Acknowledgments
The authors thank D. Cormier, H. de Vega, and J. Salgado for useful discussions. J. B. thanks the Deutsche Forschungsgemeinschaft for partial support under grant No Ba 703/6-1. K.H. thanks the Graduiertenkolleg “Erzeugung und Zerfälle von Elementarteilchen” for partial support.
## Appendix A Perturbative expansion
The mode functions $`U_k(t)`$ with the initial conditions introduced in section 2 satisfy the integral equation
$$U_k(t)=e^{i\omega _{k0}t}+\underset{0}{\overset{\mathrm{}}{}}𝑑t^{}\mathrm{\Delta }_{k,\mathrm{ret}}(tt^{})𝒱(t^{})U_k(t^{}),$$
(A.1)
with
$$\mathrm{\Delta }_{k,\mathrm{ret}}(tt^{})=\frac{1}{\omega _{k0}}\mathrm{\Theta }(tt^{})\mathrm{sin}\left(\omega _{k0}(tt^{})\right).$$
(A.2)
We separate $`U_k(t)`$ into the trivial part corresponding to the case $`𝒱(t)=0`$ and a function $`f_k(t)`$ which represents the reaction to the potential by making the ansatz
$$U_k(t)=e^{i\omega _{k0}t}[1+f_k(t)].$$
(A.3)
$`f_k(t)`$ satisfies then the integral equation
$$f_k(t)=\underset{0}{\overset{t}{}}𝑑t^{}\mathrm{\Delta }_{k,\mathrm{ret}}(tt^{})𝒱(t^{})[1+f_k(t^{})]e^{i\omega _{k0}(tt^{})},$$
(A.4)
and an equivalent differential equation
$$\ddot{f}_k(t)2i\omega _{k0}\dot{f}_k(t)=𝒱(t)[1+f_k(t)],$$
(A.5)
with the initial conditions $`f_k(0)=\dot{f}_k(0)=0`$. We expand now $`f_k(t)`$ with respect to orders in $`𝒱(t)`$ by writing
$`f_k(t)`$ $`=`$ $`f_k^{(1)}(t)+f_k^{(2)}(t)+f_k^{(3)}(t)+\mathrm{}`$ (A.6)
$`=`$ $`f_k^{(1)}(t)+f_k^{\overline{(2)}}(t),`$ (A.7)
where $`f_k^{(n)}(t)`$ is of n’th order in $`𝒱(t)`$ and $`f_k^{\overline{(n)}}(t)`$ is the sum over all orders beginning with the n’th one:
$$f_k^{\overline{(n)}}(t)=\underset{l=n}{\overset{\mathrm{}}{}}f_k^{(l)}(t).$$
(A.8)
The $`f_k^{(n)}`$ are obtained by iterating the integral equation (A.4) or the differential equation (A.5). The function $`f_k^{\overline{(1)}}(t)`$ is identical to the function $`f_k(t)`$ itself which is obtained by solving (A.5). The function $`f_k^{\overline{(2)}}(t)`$ can again be obtained by iteration via
$$f_k^{\overline{(2)}}(t)=\underset{0}{\overset{t}{}}𝑑t^{}\mathrm{\Delta }_{k,\mathrm{ret}}(tt^{})𝒱(t^{})f_k^{\overline{(1)}}(t^{})e^{i\omega _{k0}(tt^{})}.$$
(A.9)
The integral equations can be used in order to derive the asymptotic behavior as $`\omega _{k0}\mathrm{}`$ and to separate divergent and finite contributions. This has been described previously in extenso . We illustrate the procedure by calculating the relevant leading terms for $`f_k^{(1)}(t)`$. We have
$$f_k^{(1)}(t)=\frac{i}{2\omega _{k0}}\underset{0}{\overset{t}{}}𝑑t^{}(\mathrm{exp}(2i\omega _{k0}(tt^{}))1)𝒱(t^{}).$$
(A.10)
Integrating by parts we obtain
$$f_k^{(1)}(\tau )=\frac{i}{2\omega _{k0}}\underset{0}{\overset{t}{}}𝑑t^{}𝒱(t^{})\frac{1}{4\omega _{k0}^2}𝒱(t)+\frac{1}{4\omega _{k0}^2}\underset{0}{\overset{t}{}}𝑑t^{}\mathrm{exp}(2i\omega _{k0}(tt^{}))\dot{𝒱}(t^{}),$$
(A.11)
For the expansion of the fluctuation integral $`(t)`$ we need the real part of $`f_k^{(1)}`$ for which we find
$$\mathrm{Re}h_k^{(1)}(t)=\frac{1}{4\omega _{k0}^2}𝒱(t)+\frac{1}{4\omega _{k0}^2}\underset{0}{\overset{t}{}}𝑑t^{}\mathrm{cos}(2\omega _{k0}(tt^{}))\dot{𝒱}(t^{}).$$
(A.12)
The second term decreases at least as $`\omega _{k0}^3`$. In terms of the perturbative expansion for the functions $`f_k`$ we can the mode functions appearing in the fluctuation integral as
$$|U_k|^2=1+2\mathrm{R}\mathrm{e}f_k^{\overline{(1)}}+|f_k^{\overline{(1)}}|^2.$$
(A.13)
Using Eq. (A.12) the leading behavior of this expression is
$`1+2\mathrm{R}\mathrm{e}f_k^{\overline{(1)}}+|f_k^{\overline{(1)}}|^21{\displaystyle \frac{1}{2\omega _{k0}^2}}𝒱(t).`$ (A.14)
Similarly the integrand of the energy density and pressure can be expanded . As these are more divergent, the calculations require more integrations by parts in order to single out the leading powers in $`\omega _{k0}`$ and they become more involved.
## Figure Captions
Fig. 1: Phase diagram in the $`\varphi _0^2T`$ plane.
Fig. 2: Phase diagram in the $`m_0^2T`$ plane.
Fig. 3: Evolution of classical field in region II.
Fig. 4: Evolution of $`^2(t)`$ in region II.
Fig. 5: The momentum spectrum for $`T=1`$ at $`t=75`$ displaying “dynamical Bose-Einstein condensation”, with a fit $`sin^2(kt)/k^2`$.
Fig. 6: Late time amplitude $`\varphi (\mathrm{})`$ vs. initial amplitude $`\varphi _0`$ for $`T=1`$ (asteriscs), compared with Eq. (4.12) (solid line).
Fig. 7: The same as Fig. 3 for $`T=2`$.
Fig. 8: The same as Fig. 3 for $`T=3`$.
Fig. 9: Evolution of the classical field in the one-loop approximation.
Fig. 10: The fluctuation integral (solid line) and $`^2(t)`$ (dashed line) in the one-loop approximation.
Fig. 11: Evolution of the classical field in region III.
Fig. 12: Evolution of $`^2(t)`$ in region III.
Fig. 13: The asymptotic sum rule for $`^2(t)`$. The data for $`T=1.5`$ (diamonds), $`T=2.5`$ (asteriscs) and $`T=4`$ (triangles) are compared to Eq. (4.3) (solid lines).
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# 1 Introduction
## 1 Introduction
The computation of correlation functions in field theoretical models is a difficult problem. In this article we present a novel approach which applies to models where a two point function can be written as
$$S(x,y)=[P+Q]_{x,y}^1𝑑\mu (Q).$$
$`(1.1)`$
Here $`P`$ is some operator diagonal in momentum space, typically determined by the unperturbed Hamiltonian, and $`Q`$ is diagonal in coordinate space. The functional integral is taken with respect to some probability measure $`d\mu (Q)`$ and goes over the matrix elements of $`Q`$. $`[]_{x,y}^1`$ denotes the $`x,y`$-entry of the matrix $`[P+Q]^1`$. Our starting point is always a model in finite volume and with positive lattice spacing in which case the operator $`P+Q`$ and the functional integral in (1.1) becomes huge- but finite-dimensional. In the end we take the infinite volume limit and, if wanted, the continuum limit.
Our treatment is based on the following identity which is obtained by repeated application of the Feshbach formula (Lemma 3.2 below). It is proven in Theorem 3.3. Let $`B=(B_{kp})_{k,p}^{N\times N}`$, $``$ some index set, $`||=N`$ and let
$$G(k):=\left[B^1\right]_{kk}$$
$`(1.2)`$
Then one has
$$G(k)=\frac{1}{B_{kk}+{\displaystyle \underset{r=2}{\overset{N}{}}}(1)^{r+1}{\displaystyle \underset{\genfrac{}{}{0pt}{}{p_2\mathrm{}p_r\{k\}}{p_ip_j}}{}}B_{kp_2}G_k(p_2)B_{p_2p_3}\mathrm{}B_{p_{r1}p_r}G_{kp_2\mathrm{}p_{r1}}(p_r)B_{p_rk}}$$
$`(1.3)`$
where $`G_{k_1\mathrm{}k_j}(p)=\left[(B_{st})_{s,t\{k_1\mathrm{}k_j\}}\right]_{pp}^1`$ is the $`p,p`$ entry of the inverse of the matrix which is obtained from $`B`$ by deleting the rows and columns labelled by $`k_1,\mathrm{},k_j`$. In Section 2 we apply this formula to a matrix of the form $`B=`$ self adjoint + $`i\epsilon Id`$, which, for $`\epsilon 0`$, has the property that all submatrices $`(B_{st})_{s,t\{k_1\mathrm{}k_j\}}`$ are invertible.
There is also a formula for the off-diagonal inverse matrix elements. It reads
$$\left[B^1\right]_{kp}=G(k)B_{kp}G_k(p)+\underset{r=3}{\overset{N}{}}(1)^{r+1}\underset{\genfrac{}{}{0pt}{}{t_3\mathrm{}t_r\{k,p\}}{t_it_j}}{}G(k)B_{kt_3}G_k(t_3)B_{t_3t_4}\mathrm{}B_{t_rp}G_{kt_3\mathrm{}t_r}(p)$$
$`(1.4)`$
These formulae also hold in the case where the matrix $`B`$ has a block structure $`B_{kp}=(B_{k\sigma ,p\tau })`$ where, say, $`\sigma ,\tau \{,\}`$ are some spin variables. In that case the $`B_{kp}`$ are small matrices, the $`G_{k_1\mathrm{}k_j}(p)`$ are matrices of the same size and the 1/$``$ in (1.3) means inversion of matrices, see Theorem 3.3 below.
The paper is organized as follows. In the next section we demonstrate the method by applying it to the averaged Green function of the Anderson model. The Schwinger-Dyson equation for that model reads $`G^1=G_0^1+\mathrm{\Sigma }(G_0)`$ where $`\mathrm{\Sigma }(G_0)`$ is the sum of all two-legged one-particle irreducible diagrams. Application of (1.3) leads to an integral equation $`G^1=G_0^1+\sigma (G)`$ where $`\sigma (G)`$ is the sum of all two-legged graphs without two-legged subgraphs. The latter equation has two advantages. First, $`\mathrm{\Sigma }`$ is the sum of one-particle irreducible diagrams, but these diagrams may very well have two-legged subdiagrams and usually these are the diagrams which produce anomalously large contributions. And second, the propagator for $`\sigma (G)`$ is the interacting two point function $`G`$, which, for the Anderson model, is more regular than the free two point function $`G_0`$ which is the propagator for the diagrams contributing to $`\mathrm{\Sigma }(G_0)`$. More precisely, the series for $`\sigma (G)`$ can be expected to be asymptotic, that is, its lowest order contributions are a good approximation if the coupling is small, but, usually, the series for $`\mathrm{\Sigma }(G_0)`$ is not asymptotic.
For the many-electron system and for the $`\phi ^4`$ model repeated application of (1.3,4) amounts to a resummation of two- and four-legged subgraphs. This is discussed in section 4. In section 5 we discuss how our method is related to the integral equations which can be found in the literature. The proof of the inversion formula is given in section 3.
## 2 Application to the Anderson Model
Let coordinate space be a lattice of finite volume with periodic boundary conditions, lattice spacing $`1/M`$ and volume $`[0,L]^d`$:
$$\mathrm{\Gamma }=\left\{x=\frac{1}{M}(n_1,\mathrm{},n_d)\right|\mathrm{\hspace{0.33em}0}n_iML1\}=\left(\frac{1}{M}\right)^d/(L)^d$$
$`(2.1)`$
Momentum space is given by
$$:=\mathrm{\Gamma }^{\mathrm{}}=\left\{k=\frac{2\pi }{L}(m_1,\mathrm{},m_d)\right|\mathrm{\hspace{0.33em}0}m_iML1\}=\left(\frac{2\pi }{L}\right)^d/(2\pi M)^d$$
$`(2.2)`$
We consider the averaged Green function of the Anderson model given by
$$G(x,x^{}):=\left[\mathrm{\Delta }z+\lambda V\right]_{x,x^{}}^1𝑑P(V)$$
$`(2.3)`$
where the random potential is Gaussian,
$$dP(V)=\underset{x\mathrm{\Gamma }}{\Pi }e^{\frac{V_x^2}{2}}\frac{dV_x}{\sqrt{2\pi }}.$$
$`(2.4)`$
Here $`z=E+i\epsilon `$ and $`\mathrm{\Delta }`$ is the discrete Laplacian,
$$\left[\mathrm{\Delta }z+\lambda V\right]_{x,x^{}}=M^2\underset{i=1}{\overset{d}{}}\left(\delta _{x^{},x+e_i/M}+\delta _{x^{},xe_i/M}2\delta _{x^{},x}\right)z\delta _{x,x^{}}+\lambda V_x\delta _{x,x^{}}$$
$`(2.5)`$
By taking the Fourier transform, one has
$`G(x,x^{})`$ $`=`$ $`\frac{1}{M^dL^d}{\displaystyle \underset{k}{}}e^{ik(x^{}x)}G(k)`$ (2.6)
$`G(k)`$ $`=`$ $`{\displaystyle _{^{N^d}}}\left[a_k\delta _{k,p}+\frac{\lambda }{\sqrt{N^d}}v_{kp}\right]_{k,k}^1𝑑P(v)`$ (2.7)
where $`N^d=(ML)^d=|\mathrm{\Gamma }|=||`$, $`dP(v)`$ is given by (2.10) or (2.11) below, depending on whether $`N^d`$ is even or odd, and
$$a_k=4M^2\underset{i=1}{\overset{d}{}}\mathrm{sin}^2\left[\frac{k_i}{2M}\right]Ei\epsilon $$
$`(2.8)`$
The rigorous control of $`G(k)`$ for small disorder $`\lambda `$ and energies inside the spectrum of the unperturbed Laplacian, $`E[0,4M^2]`$, in which case $`a_k`$ has a root if $`\epsilon 0`$, is still an open problem \[AG,K,MPR,P,W\]. It is expected that $`lim_{\epsilon 0}lim_L\mathrm{}G(k)=1/(a_k\sigma _k)`$ where Im$`\sigma =O(\lambda ^2)`$.
The integration variables $`v_q`$ in (2.7) are given by the discrete Fourier transform of the $`V_x`$. In particular, observe that, if $`F`$ denotes the unitary matrix of discrete Fourier transform, the variables
$$v_q(FV)_q=\frac{1}{\sqrt{N^d}}\underset{x\mathrm{\Gamma }}{}e^{iqx}V_x=\left(\frac{M}{L}\right)^{\frac{d}{2}}\frac{1}{M^d}\underset{x\mathrm{\Gamma }}{}e^{iqx}V_x\left(\frac{M}{L}\right)^{\frac{d}{2}}\widehat{V}_q$$
$`(2.9)`$
would not have a limit if $`V_x`$ would be deterministic and cutoffs are removed, since the $`\widehat{V}_q`$ are the quantities which have a limit in that case. But since the $`V_x`$ are integration variables, we choose a unitary transform to keep the integration measure invariant. Observe also that $`v_q`$ is complex, $`v_q=u_q+iw_q`$. Since $`V_x`$ is real, $`u_q=u_q`$ and $`w_q=w_q`$. In order to transform $`dP(V)`$ to momentum space, we have to choose a set $`^+`$ such that either $`q^+`$ or $`q^+`$. If $`N`$ is odd, the only momentum with $`q=q`$ or $`w_q=0`$ is $`q=0`$. In that case $`dP(V)`$ becomes
$$dP(v)=e^{\frac{u_0^2}{2}}\frac{du_0}{\sqrt{2\pi }}\underset{q^+}{\Pi }e^{(u_q^2+w_q^2)}\frac{du_qdw_q}{\pi }$$
$`(2.10)`$
For even $`N`$ we get
$$dP(v)=e^{\frac{1}{2}(u_0^2+u_{q_0}^2)}\frac{du_0du_{q_0}}{2\pi }\underset{q^+}{\Pi }e^{(u_q^2+w_q^2)}\frac{du_qdw_q}{\pi }$$
$`(2.11)`$
where $`q_0=\frac{2\pi m}{L}`$ is the unique nonzero momentum for which $`\frac{2\pi }{L}m=2\pi M(1,\mathrm{},1)\frac{2\pi }{L}m`$.
Now we apply the inversion formula (1.3) to the inverse matrix element in
$$G(k)=_{^{N^d}}\left[a_k\delta _{k,p}+\frac{\lambda }{\sqrt{N^d}}v_{kp}\right]_{k,k}^1𝑑P(v)$$
$`(2.7)`$
We start with the ‘two loop approximation’, which we define by retaining only the $`r=2`$ term in the denominator of the right hand side of (1.3),
$$G(k)\frac{1}{B_{kk}_{p\{k\}}B_{kp}G_k(p)B_{pk}}$$
$`(2.12)`$
Thus, let
$$G(k):=\left[a_k\delta _{k,p}+\frac{\lambda }{\sqrt{N^d}}v_{kp}\right]_{k,k}^1=G(k;v,z)$$
$`(2.13)`$
In the infinite volume limit the spacing $`2\pi /L`$ of the momenta becomes arbitrary small. Hence, in computing an inverse matrix element, it should not matter whether a single column and row labelled by some momentum $`t`$ is absent or not. In other words, in the infinite volume limit one should have
$$G_t(p)=G(p)\mathrm{for}L\mathrm{}$$
$`(2.14)`$
and similarly $`G_{t_1\mathrm{}t_j}(p)=G(p)`$ as long as $`j`$ is independent of the volume. We remark however that if the matrix has a block structure, say $`B=(B_{k\sigma ,p\tau })`$ with $`\sigma ,\tau \{,\}`$ some spin variables, this structure has to be respected. That is, for a given momentum $`k`$ all rows and columns labelled by $`k,k`$ have to be eliminated, since otherwise (2.14) may not be true.
Thus the two loop approximation gives
$$G(k)=\frac{1}{a_k+\frac{\lambda }{\sqrt{N^d}}v_0\frac{\lambda ^2}{N^d}_{pk}v_{kp}G(p)v_{pk}}$$
$`(2.15)`$
For large $`L`$, we can disregard the $`\frac{\lambda }{\sqrt{N^d}}v_0`$ term. Introducing $`\sigma _k=\sigma _k(v,z)`$ according to
$$G(k)=:\frac{1}{a_k\sigma _k},$$
$`(2.16)`$
we get
$$\sigma _k=\frac{\lambda ^2}{N^d}\underset{pk}{}\frac{|v_{kp}|^2}{a_p\sigma _p}\frac{\lambda ^2}{N^d}\underset{p}{}\frac{|v_{kp}|^2}{a_p\sigma _p}$$
$`(2.17)`$
and arrive at
$$G(k)=\frac{1}{a_k\frac{\lambda ^2}{N^d}{\displaystyle \underset{p}{\overset{}{}}}{\displaystyle \frac{|v_{kp}|^2}{a_p\frac{\lambda ^2}{N^d}_t\frac{|v_{pt}|^2}{a_t\frac{\lambda ^2}{N^d}\mathrm{\Sigma }\mathrm{}}}}}𝑑P(v)$$
$`(2.18)`$
Now consider the infinite volume limit $`L\mathrm{}`$ or $`N=ML\mathrm{}`$. By the central limit theorem of probability $`\frac{1}{\sqrt{N^d}}_q\left(|v_q|^2|v_q|^2\right)`$ is, as a sum of independent random variables, normal distributed. Note that only a prefactor of $`1/\sqrt{N^d}`$ is required for that property. In particular, if $`F`$ is some bounded function independent of $`N`$, sums which come with a prefactor of $`1/N^d`$ like $`\frac{1}{N^d}_qc_q|v_q|^2`$ can be substituted by their expectation value,
$$\underset{N\mathrm{}}{lim}F\left(\frac{1}{N^d}\underset{k}{}c_k|v_k|^2\right)𝑑P(v)=F\left(\underset{N\mathrm{}}{lim}\frac{1}{N^d}\underset{k}{}c_k|v_k|^2\right)$$
$`(2.19)`$
Therefore, in the two loop approximation, one obtains in the infinite volume limit
$$G(k)=\frac{1}{a_k\frac{\lambda ^2}{N^d}{\displaystyle \underset{p}{\overset{}{}}}{\displaystyle \frac{|v_{kp}|^2}{a_p\frac{\lambda ^2}{N^d}_t\frac{|v_{pt}|^2}{a_t\frac{\lambda ^2}{N^d}\mathrm{\Sigma }\mathrm{}}}}}=:\frac{1}{a_k\sigma _k}$$
$`(2.20)`$
where the quantity $`\sigma _k`$ satisfies the integral equation
$$\sigma _k=\frac{\lambda ^2}{N^d}\underset{p}{}\frac{|v_{kp}|^2}{a_p\sigma _p}\stackrel{L\mathrm{}}{}\frac{\lambda ^2}{M^d}_{[0,2\pi M]^d}\frac{d^dp}{(2\pi )^d}\frac{|v_{kp}|^2}{a_p\sigma _p}$$
$`(2.21)`$
For a Gaussian distribution $`|v_q|^2=1`$ for all $`q`$ such that $`\sigma _k=\sigma `$ becomes independent of $`k`$. Thus we end up with
$$G(k)=\frac{1}{4M^2_{i=1}^d\mathrm{sin}^2\left[\frac{k_i}{2M}\right]Ei\epsilon \sigma }$$
$`(2.22)`$
where $`\sigma `$ is a solution of
$`\sigma `$ $`=`$ $`\frac{\lambda ^2}{M^d}{\displaystyle _{[0,2\pi M]^d}}{\displaystyle \frac{d^dp}{(2\pi )^d}}{\displaystyle \frac{1}{4M^2_{i=1}^d\mathrm{sin}^2\left[\frac{p_i}{2M}\right]z\sigma }}`$ (2.23)
$`=`$ $`\lambda ^2{\displaystyle _{[0,2\pi ]^d}}{\displaystyle \frac{d^dp}{(2\pi )^d}}{\displaystyle \frac{1}{4M^2_{i=1}^d\mathrm{sin}^2\left[\frac{p_i}{2}\right]z\sigma }}.`$
This equation is of course well known and one deduces from it that it generates a small imaginary part Im$`\sigma =O(\lambda ^2)`$ if the energy $`E`$ is within the spectrum of $`\mathrm{\Delta }`$.
We now add the higher loop terms (the terms for $`r>2`$ in the denominator of (1.3)) to our discussion and give an interpretation in terms of Feynman diagrams. To make the volume factors more explicit, asume that the lattice spacing in coordinate space is $`1/M=1`$ such that $`N=L`$.
For the Anderson model, Feynman graphs may be obtained by brutally expanding
$`{\displaystyle \left[a_k\delta _{k,p}+\frac{\lambda }{\sqrt{L}^d}v_{kp}\right]_{k,k}^1𝑑P}={\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}{\displaystyle \left(C[VC]^r\right)_{kk}𝑑P}`$
$`={\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}\frac{(\lambda )^r}{\sqrt{L^d}^r}{\displaystyle \underset{p_2\mathrm{}p_r}{}}\frac{1}{a_ka_{p_2}\mathrm{}a_{p_r}a_k}{\displaystyle v_{kp_2}v_{p_2p_3}\mathrm{}v_{p_rk}𝑑P}`$
For a given $`r`$, this may be represented as in figure 1 ($`c_k:=1/a_k`$).
The integral over the $`v`$ gives a sum of $`(r1)!!`$ terms where each term is a product of $`r/2`$ Kroenecker-delta’s, the terms for odd $`r`$ vanish. If this is substituted in (2.24), the number of independent momentum sums is cut down to $`r/2`$ and each of the $`(r1)!!`$ terms may be represented by a diagram
where, say, the value of the third diagram is given by $`\frac{\lambda ^4}{L^{2d}}_{p_1,p_2}\frac{1}{a_ka_{k+p_1}a_{k+p_1+p_2}a_{k+p_2}a_k}`$. For short:
$$G(k)=\mathrm{sum}\mathrm{of}\mathrm{all}\mathrm{two}\mathrm{legged}\mathrm{diagrams}.$$
$`(2.25)`$
Since the value of a diagram depends on its subgraph structure, one distinguishes, in the easiest case, two types of diagrams. Diagrams with or without two-legged subdiagrams. Those diagrams with two-legged subgraphs usually produce anomalously large contributions. They are devided further into the one-particle irreducible ones and the reducible ones. Thereby a diagram is called one-particle reducible if it can be made disconnected by cutting one solid or ‘particle’ line (no squiggle or dashed line), see also figure 3.
The reason for introducing reducible and irreducible diagrams is that the reducible ones can be easily resummed by writing down the Schwinger-Dyson equation which states that if the self energy $`\mathrm{\Sigma }_k`$ is defined through
$$G(k)=\frac{1}{a_k\mathrm{\Sigma }_k(G_0)}$$
$`(2.26)`$
then $`\mathrm{\Sigma }_k(G_0)`$ is the sum of all amputated (no $`1/a_k`$’s at the ends) one particle irreducible diagrams. Here we wrote $`\mathrm{\Sigma }_k(G_0)`$ to indicate that the factors (‘propagators’) assigned to the solid lines of the diagrams contributing to $`\mathrm{\Sigma }_k`$ are given by the free two point function $`G_0(p)=\frac{1}{a_p}`$. However, the diagrams contributing to $`\mathrm{\Sigma }_k(G_0)`$ still contain anomalously large contributions, namely irreducible graphs which contain two-legged subgraphs like diagram (c) in figure 3.
In the following we show, using the inversion formula (1.3) including all higher loop terms, that all graphs with two-legged subgraphs can be eliminated or resummed by writing down the following integral equation for $`G`$:
$$G(k)=\frac{1}{a_k\sigma _k(G)}$$
$`(2.27)`$
where
> $`\sigma _k(G)`$ is the sum of all amputated two-legged diagrams which do not contain any two-legged subdiagrams, but now with propagators $`G(k)=\frac{1}{a_k\sigma _k}`$ instead of $`G_0=\frac{1}{a_k}`$
which may be formalized as in (2.35) below. The advantage of this formula is that the series for $`\sigma _k(G)`$ can be expected to be asymptotic, that is, its lowest order contributions are a good approximation if the coupling is small, but, usually, the series for $`\mathrm{\Sigma }_k(G_0)`$ is not asymptotic. Thus, in order to rigorously controll $`G(k)`$, one has to define a suitable space of propagators, to estimate the sum of all amputated two-legged graphs without two-legged subgraphs on that space and then finally to show that the equation (2.27) has a solution on this space. We intend to address this problem in another paper.
We now show (2.27) for the Anderson model. For fixed $`v`$ one has
$$G(k,v)=\frac{1}{a_k\sigma _k(v)}$$
$`(2.28)`$
where
$$\sigma _k(v)=\underset{r=2}{\overset{L^d}{}}(1)^r\underset{\genfrac{}{}{0pt}{}{p_2\mathrm{}p_r}{p_ip_j,p_ik}}{}\frac{\lambda ^r}{\sqrt{L^d}^r}G_k(p_2)\mathrm{}G_{kp_2\mathrm{}p_{r1}}(p_r)v_{kp_2}\mathrm{}v_{p_rk}$$
$`(2.29)`$
We cutoff the $`r`$-sum in (2.29) at some arbitrary but fixed order $`\mathrm{}<L^d`$ where $`\mathrm{}`$ is choosen to be independent of the volume. Furthermore we substitute $`G_{kp_2\mathrm{}p_j}(p)`$ by $`G(p)`$. Thus
$$G(k)=\frac{1}{a_k_{r=2}^{\mathrm{}}\sigma _k^r(v)}$$
$`(2.30)`$
where
$$\sigma _k^r(v)=(1)^r\frac{\lambda ^r}{\sqrt{L^d}^r}\underset{\genfrac{}{}{0pt}{}{p_2\mathrm{}p_r}{p_ip_j,p_ik}}{}G(p_2)\mathrm{}G(p_r)v_{kp_2}\mathrm{}v_{p_rk}$$
$`(2.31)`$
Consider first two strings $`s_k^{r_1}`$, $`s_k^{r_2}`$ where
$$s_k^r(v)=\frac{\lambda ^r}{\sqrt{L^d}^r}\underset{\genfrac{}{}{0pt}{}{p_2\mathrm{}p_r}{p_ip_j,p_ik}}{}c_{kp_2\mathrm{}p_r}^rv_{kp_2}\mathrm{}v_{p_rk}$$
$`(2.32)`$
and the $`c_{kp_2\mathrm{}p_r}^r`$ are some numbers. Then in the infinite volume limit
$$s_k^{r_1}s_k^{r_2}=s_k^{r_1}s_k^{r_2}$$
$`(2.33)`$
because all pairings which connect the two strings have an extra volume factor $`1/L^d`$. Namely, if the two strings are disconnected, there are $`(r_1+r_2)/2`$ loops and a volume factor of $`1/\sqrt{L^d}^{(r_1+r_2)}`$ giving $`(r_1+r_2)/2`$ Riemannian sums. If the two strings are connected, there are only $`(r_1+r_22)/2`$ loops leaving an extra factor of $`1/L^d`$. By the same argument one has in the infinite volume limit
$$(s_k^{r_1})^{n_1}\mathrm{}(s_k^{r_m})^{n_m}=s_k^{r_1}^{n_1}\mathrm{}s_k^{r_m}^{n_m}$$
$`(2.34)`$
which results in
$$G(k)=\frac{1}{a_k{\displaystyle \underset{r=2}{\overset{\mathrm{}}{}}}\frac{(\lambda )^r}{\sqrt{L^d}^r}{\displaystyle \underset{\genfrac{}{}{0pt}{}{p_2\mathrm{}p_r}{p_ip_j,p_ik}}{}}G(p_2)\mathrm{}G(p_r)v_{kp_2}\mathrm{}v_{p_rk}}$$
$`(2.35)`$
The condition $`p_2,\mathrm{},p_rk`$ and $`p_ip_j`$ means exactly that two-legged subgraphs are forbidden. Namely, for a two-legged subdiagram as in (c) in figure 3, the incomming and outgoing momenta $`p`$, $`p^{}`$ (to which are assigned propagators $`G(p)`$, $`G(p^{})`$) must be equal which is forbidden by the condidtion $`p_ip_j`$ in (2.35).
However, we cannot take the limit $`\mathrm{}\mathrm{}`$ in (2.35) since the series in the denominator of (2.35) is only an asymptotic one. To see this a bit more clearly suppose for the moment that there were no restrictions on the momentum sums. Then, if $`V=(\frac{\lambda }{\sqrt{L^d}}v_{kp})_{k,p}`$ and $`G=(G(k)\delta _{k,p})_{k,p}`$,
$$\frac{\lambda ^r}{\sqrt{L^d}^r}\underset{p_2\mathrm{}p_r}{}G(p_2)\mathrm{}G(p_r)v_{kp_2}\mathrm{}v_{p_rk}=(V[GV]^{r1})_{kk}$$
$`(2.36)`$
and for $`\mathrm{}\mathrm{}`$ we would get
$$G(k)=\frac{1}{a_k(V\frac{GV}{Id+GV})_{kk}}=\frac{1}{a_k(V\frac{1}{G^1+V}V)_{kk}}$$
$`(2.37)`$
That is, the factorials produced by the number of diagrams in the denominator of (2.35) are basically the same as those in the expansion
$$_{}\frac{x^2}{z+\lambda x}e^{\frac{x^2}{2}}\frac{dx}{\sqrt{2\pi }}=\underset{r=0}{\overset{\mathrm{}}{}}\frac{\lambda ^{2r}}{z^{2r+1}}(2r+1)!!+R_{\mathrm{}+1}(\lambda )$$
$`(2.38)`$
where the remainder satisfies the bound $`|R_{\mathrm{}+1}(\lambda )|\mathrm{}!const_z^{\mathrm{}}\lambda ^{\mathrm{}}`$.
We close this section with two further remarks. So far the computations were done in momentum space. One may wonder what one gets if the inversion formula (1.3) is applied to $`[\mathrm{\Delta }+z+\lambda V]^1`$ in coordinate space. Whereas a geometric series expansion of $`[\mathrm{\Delta }+z+\lambda V]^1`$ gives a representation in terms of the simple random walk, application of (1.3) results in a representation in terms of the self avoiding walk:
$$[\mathrm{\Delta }+z+\lambda V]_{0,x}^1=\underset{\genfrac{}{}{0pt}{}{\gamma :0x}{\gamma \mathrm{self}\mathrm{avoiding}}}{}\frac{det\left[(\mathrm{\Delta }+z+\lambda V)_{y,y^{}\mathrm{\Gamma }\gamma }\right]}{det\left[(\mathrm{\Delta }+z+\lambda V)_{y,y^{}\mathrm{\Gamma }}\right]}$$
$`(2.39)`$
where $`\mathrm{\Gamma }`$ is the lattice in coordinate space. Namely, if $`|x|>1`$, the inversion formula (1.4) for the off-diagonal elements gives
$`[\mathrm{\Delta }+\lambda V]_{0,x}^1={\displaystyle \underset{r=3}{\overset{L^d}{}}}(1)^{r+1}{\displaystyle \underset{\genfrac{}{}{0pt}{}{x_3\mathrm{}x_r\mathrm{\Gamma }\{0,x\}}{x_ix_j}}{}}G(0)G_0(x_3)\mathrm{}G_{0x_3\mathrm{}x_r}(x)(\mathrm{\Delta })_{0x_3}\mathrm{}(\mathrm{\Delta })_{x_rx}`$
$`={\displaystyle \underset{r=3}{\overset{L^d}{}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{x_2=0,x_3,\mathrm{},x_r,x_{r+1}=x\mathrm{\Gamma }}{|x_ix_{i+1}|=1i=2\mathrm{}r}}{}}{\displaystyle \frac{det\left[(\mathrm{\Delta }+\lambda V)_{y,y^{}\mathrm{\Gamma }\{0\}}\right]}{det\left[(\mathrm{\Delta }+\lambda V)_{y,y^{}\mathrm{\Gamma }}\right]}}\mathrm{}{\displaystyle \frac{det\left[(\mathrm{\Delta }+\lambda V)_{y,y^{}\mathrm{\Gamma }\{0,x_3\mathrm{}x_r,x\}}\right]}{det\left[(\mathrm{\Delta }+\lambda V)_{y,y^{}\mathrm{\Gamma }\{0,x_3\mathrm{}x_r\}}\right]}}`$
which coincides with (2.39).
Finally we remark that, while the argument following (2.32) leads to a factorization property for on-diagonal elements in momentum space, $`G(k)G(p)=G(k)G(p)`$, there is no such property for products of off-diagonal elements which appear in a quantity like
$$\mathrm{\Lambda }(q)=\frac{1}{L^d}\underset{k,p}{}\left[a_k\delta _{k,p}+\frac{\lambda }{\sqrt{L}^d}v_{kp}\right]_{k,p}^1\left[\overline{a}_k\delta _{k,p}+\frac{\lambda }{\sqrt{L}^d}\overline{v}_{kp}\right]_{kq,pq}^1$$
$`(2.40)`$
which is the Fourier transform of $`\left|[\mathrm{\Delta }+z+\lambda V]_{x,y}^1\right|^2`$. (Each off-diagonal inverse matrix element is proportional to $`1/\sqrt{L^d}`$, therefore the prefactor of $`1/L^d`$ in (2.40) is correct.)
## 3 Proof of the Inversion Formula
Lemma 3.1: Let $`B^{k\times n}`$, $`C^{n\times k}`$ and let $`Id_k`$ denote the identity in $`^{k\times k}`$. Then:
1. $`Id_kBC\mathrm{invertible}Id_nCB\mathrm{invertible}`$.
2. If the left or the right hand side of (i) fullfilled, then $`C\frac{1}{Id_kBC}=\frac{1}{Id_nCB}C`$.
Proof: Let
$$B=\left(\begin{array}{ccc}& \stackrel{}{b}_1& \\ & \mathrm{}& \\ & \stackrel{}{b}_k& \end{array}\right),C=\left(\begin{array}{ccc}|& & |\\ \stackrel{}{c}_1& \mathrm{}& \stackrel{}{c}_k\\ |& & |\end{array}\right)$$
where the $`\stackrel{}{b}_j`$ are $`n`$-component row vectors and the $`\stackrel{}{c}_j`$ are $`n`$-component column vectors. Let $`\stackrel{}{x}\mathrm{Kern}(IdCB)`$. Then $`\stackrel{}{x}=CB\stackrel{}{x}=_j\lambda _j\stackrel{}{c}_j`$ if we define $`\lambda _j:=(\stackrel{}{b}_j,\stackrel{}{x})`$. Let $`\stackrel{}{\lambda }=(\lambda _j)_{1jk}`$. Then $`[(IdBC)\stackrel{}{\lambda }]_i=\lambda _i_j(\stackrel{}{b}_i,\stackrel{}{c}_j)\lambda _j=(\stackrel{}{b}_i,\stackrel{}{x})_j(\stackrel{}{b}_i,\stackrel{}{c}_j)\lambda _j=0`$ since $`\stackrel{}{x}=_j\lambda _j\stackrel{}{c}_j`$, thus $`\stackrel{}{\lambda }\mathrm{Kern}(IdBC)`$. On the other hand, if some $`\stackrel{}{\lambda }\mathrm{Kern}(IdBC)`$, then $`\stackrel{}{x}:=_j\lambda _j\stackrel{}{c}_j\mathrm{Kern}(IdCB)`$ which proves (i). Part (ii) then follows from $`C=\frac{1}{Id_nCB}(Id_nCB)C=\frac{1}{Id_nCB}C(Id_kBC)`$ $`\mathrm{}`$
The inversion formula (1.3,4) is obtained by iterative application of the next lemma, which states the Feshbach formula for finite dimensional matrices. For a more general version one may look in \[BFS\], Theorem 2.1.
Lemma 3.2: Let $`h=\left(\genfrac{}{}{0pt}{}{A}{C}\genfrac{}{}{0pt}{}{B}{D}\right)^{n\times n}`$ where $`A^{k\times k}`$, $`D^{(nk)\times (nk)}`$ are invertible and $`B^{k\times (nk)}`$, $`C^{(nk)\times k}`$. Then
$$h\mathrm{invertible}ABD^1C\mathrm{invertible}DCA^1B\mathrm{invertible}$$
$`(3.1)`$
and if one of the conditions in (3.1) is fullfilled, one has $`h^1=\left(\genfrac{}{}{0pt}{}{E}{G}\genfrac{}{}{0pt}{}{F}{H}\right)`$ where
$`E=\frac{1}{ABD^1C},`$ $`H=\frac{1}{DCA^1B},`$ (3.2)
$`F=EBD^1=A^1BH,`$ $`G=HCA^1=D^1CE.`$ (3.3)
Proof: We have, using Lemma 3.1 in the second line,
$`ABD^1C\mathrm{inv}.`$ $``$ $`Id_kA^1BD^1C\mathrm{inv}.`$
$``$ $`Id_{nk}D^1CA^1B\mathrm{inv}.`$
$``$ $`DCA^1B\mathrm{inv}.`$
Furthermore, again by Lemma 3.1,
$$D^1C\frac{1}{IdA^1BD^1C}=\frac{1}{IdD^1CA^1B}D^1C=\frac{1}{DCA^1B}C=HC$$
and
$$A^1B\frac{1}{IdD^1CA^1B}=\frac{1}{IdA^1BD^1C}A^1B=\frac{1}{ABD^1C}B=EB$$
which proves the last equalities in (3.3), $`HCA^1=D^1CE`$ and $`EBD^1=A^1BH`$. Using these equations and the definition of $`E,F,G`$ and $`H`$ one computes
$$\left(\genfrac{}{}{0pt}{}{A}{C}\genfrac{}{}{0pt}{}{B}{D}\right)\left(\genfrac{}{}{0pt}{}{E}{G}\genfrac{}{}{0pt}{}{F}{H}\right)=\left(\genfrac{}{}{0pt}{}{E}{G}\genfrac{}{}{0pt}{}{F}{H}\right)\left(\genfrac{}{}{0pt}{}{A}{C}\genfrac{}{}{0pt}{}{B}{D}\right)=\left(\genfrac{}{}{0pt}{}{Id}{0}\genfrac{}{}{0pt}{}{0}{Id}\right)$$
It remains to show that the invertibility of $`h`$ implies the invertibility of $`ABD^1C`$. To this end let $`P=\left(\genfrac{}{}{0pt}{}{Id}{}\genfrac{}{}{0pt}{}{}{0}\right)`$, $`\overline{P}=\left(\genfrac{}{}{0pt}{}{0}{}\genfrac{}{}{0pt}{}{}{Id}\right)`$ such that $`ABD^1C=PhPPh\overline{P}(\overline{P}h\overline{P})^1\overline{P}hP`$. Then
$`(ABD^1C)Ph^1P`$ $`=`$ $`PhPh^1PPh\overline{P}(\overline{P}h\overline{P})^1\overline{P}hPh^1P`$
$`=`$ $`Ph(1\overline{P})h^1PPh\overline{P}(\overline{P}h\overline{P})^1\overline{P}h(1\overline{P})h^1P`$
$`=`$ $`PPh\overline{P}h^1P+Ph\overline{P}h^1P=P`$
and similarly $`Ph^1P(ABD^1C)=P`$ which proves the invertibility of $`ABD^1C`$ $`\mathrm{}`$
Theorem 3.3: Let $`B^{nN\times nN}`$ be given by $`B=(B_{kp})_{k,p}`$, $``$ some index set, $`||=N`$, and $`B_{kp}=(B_{k\sigma ,p\tau })_{\sigma ,\tau I}^{n\times n}`$ where $`I`$ is another index set, $`|I|=n`$. Suppose that $`B`$ and, for any $`𝒩`$, the submatrix $`(B_{kp})_{k,p𝒩}`$ is invertible. For $`k`$ let
$$G(k):=\left[B^1\right]_{kk}^{n\times n}$$
$`(3.4)`$
and, if $`𝒩`$, $`k𝒩`$,
$$G_𝒩(k):=\left[\{(B_{st})_{s,t𝒩}\}^1\right]_{kk}^{n\times n}$$
$`(3.5)`$
Then one has
1. The on-diagonal block matrices of $`B^1`$ are given by
$$G(k)=\frac{1}{B_{kk}{\displaystyle \underset{r=2}{\overset{N}{}}}(1)^r{\displaystyle \underset{\genfrac{}{}{0pt}{}{p_2\mathrm{}p_r\{k\}}{p_ip_j}}{}}B_{kp_2}G_k(p_2)B_{p_2p_3}\mathrm{}B_{p_{r1}p_r}G_{kp_2\mathrm{}p_{r1}}(p_r)B_{p_rk}}$$
$`(3.6)`$
where $`1/`$ is inversion of $`n\times n`$ matrices.
2. Let $`k,p`$, $`kp`$. Then the off-diagonal block matrices of $`B^1`$ can be expressed in terms of the $`G_𝒩(s)`$ and the $`B_{st}`$,
$$[B^1]_{kp}=G(k)B_{kp}G_k(p)\underset{r=3}{\overset{N}{}}(1)^r\underset{\genfrac{}{}{0pt}{}{t_3\mathrm{}t_r\{k,p\}}{t_it_j}}{}G(k)B_{kt_3}G_k(t_3)B_{t_3t_4}\mathrm{}B_{t_rp}G_{kt_3\mathrm{}t_r}(p)$$
$`(3.7)`$
Proof: Let $`k`$ be fixed and let $`p,p^{}\{k\}`$ below label columns and rows. By Lemma 3.2 we have
$$\left(\begin{array}{cccc}G(k)& & & \\ & & & \\ & & & \end{array}\right)=\left(\begin{array}{cccc}B_{kk}& & B_{kp}& \\ |& & & \\ B_{p^{}k}& & B_{p^{}p}& \\ |& & & \end{array}\right)^1=\left(\begin{array}{cccc}E& & F& \\ |& & & \\ G& & H& \\ |& & & \end{array}\right)$$
where
$`G(k)=E`$ $`=`$ $`{\displaystyle \frac{1}{B_{kk}{\displaystyle \underset{p,p^{}k}{}}B_{kp}\left[\{(B_{p^{}p})_{p^{},p\{k\}}\}^1\right]_{pp^{}}B_{p^{}k}}}`$
$`=`$ $`{\displaystyle \frac{1}{B_{kk}{\displaystyle \underset{pk}{}}B_{kp}G_k(p)B_{pk}{\displaystyle \underset{\genfrac{}{}{0pt}{}{p,p^{}k}{pp^{}}}{}}B_{kp}\left[\{(B_{p^{}p})_{p^{},p\{k\}}\}^1\right]_{pp^{}}B_{p^{}k}}}`$
and
$`F_{kp}=\left[B^1\right]_{kp}`$ $`=`$ $`G(k){\displaystyle \underset{tk}{}}B_{kt}\left[\{(B_{p^{}p})_{p^{},p\{k\}}\}^1\right]_{tp}`$
$`=`$ $`G(k)B_{kp}G_k(p)G(k){\displaystyle \underset{tk,p}{}}B_{kt}\left[\{(B_{p^{}p})_{p^{},p\{k\}}\}^1\right]_{tp}`$
Apply Lemma 3.2 now to the matrix $`\{(B_{p^{}p})_{p^{},p\{k\}}\}^1`$ and proceed by induction to obtain after $`\mathrm{}`$ steps
$`G(k)`$ $`=`$ $`{\displaystyle \frac{1}{B_{kk}{\displaystyle \underset{r=2}{\overset{\mathrm{}}{}}}(1)^r{\displaystyle \underset{\genfrac{}{}{0pt}{}{p_2\mathrm{}p_r\{k\}}{p_ip_j}}{}}B_{kp_2}G_k(p_2)B_{p_2p_3}\mathrm{}B_{p_{r1}p_r}G_{kp_2\mathrm{}p_{r1}}(p_r)B_{p_rk}R_{\mathrm{}+1}}}`$ (3.10)
$`F_{kp}`$ $`=`$ $`G(k)B_{kp}G_k(p){\displaystyle \underset{r=3}{\overset{\mathrm{}}{}}}(1)^r{\displaystyle \underset{\genfrac{}{}{0pt}{}{t_3\mathrm{}t_r\{k,p\}}{t_it_j}}{}}G(k)B_{kt_3}G_k(t_3)B_{t_3t_4}\mathrm{}B_{t_rp}G_{kt_3\mathrm{}t_r}(p)\stackrel{~}{R}_{\mathrm{}+1}`$
where
$`R_{\mathrm{}+1}`$ $`=`$ $`(1)^{\mathrm{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{p_2\mathrm{}p_{\mathrm{}+1}\{k\}}{p_ip_j}}{}}B_{kp_2}G_k(p_2)\mathrm{}B_{p_\mathrm{}1p_{\mathrm{}}}\left[\{(B_{p^{}p})_{p^{},p\{kp_2\mathrm{}p_{\mathrm{}}\}}\}^1\right]_{p_{\mathrm{}}p_{\mathrm{}+1}}B_{p_{\mathrm{}+1}k}`$
$`\stackrel{~}{R}_{\mathrm{}+1}`$ $`=`$ $`(1)^{\mathrm{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{t_3\mathrm{}t_{\mathrm{}+1}\{k,p\}}{t_it_j}}{}}G(k)B_{kt_3}\mathrm{}G_{kt_3\mathrm{}t_\mathrm{}1}(t_{\mathrm{}})B_{t_{\mathrm{}}t_{\mathrm{}+1}}\left[\{(B_{p^{}p})_{p^{},p\{kt_3\mathrm{}t_{\mathrm{}}\}}\}^1\right]_{t_{\mathrm{}+1}p}`$
Since $`R_{N+1}=\stackrel{~}{R}_{N+1}=0`$ the theorem follows $`\mathrm{}`$
## 4 Application to the Many-Electron System and to the $`\phi ^4`$-Model
### 4.1 The Many-Electron System
We consider the many-electron system in the grand canonical ensemble in finite volume $`[0,L]^d`$ and at some small but positive temperature $`T=1/\beta >0`$ with attractive delta-interaction given by the Hamiltonian
$$H=H_0\lambda H_{\mathrm{int}}=\frac{1}{L^d}\underset{𝐤\sigma }{}(\frac{𝐤^2}{2m}\mu )a_{𝐤\sigma }^+a_{𝐤\sigma }\frac{\lambda }{L^{3d}}\underset{\mathrm{𝐤𝐩𝐪}}{}a_𝐤^+a_{𝐪𝐤}^+a_{𝐪𝐩}a_𝐩$$
$`(4.1)`$
Our normalization conventions concerning the volume factors are such that the canonical anticommutation relations read $`\{a_{𝐤\sigma },a_{𝐩\tau }^+\}=L^d\delta _{𝐤,𝐩}\delta _{\sigma ,\tau }`$. The momentum sums range over some subset of $`\left(\frac{2\pi }{L}\right)^d`$, say $`=\{𝐤\left(\frac{2\pi }{L}\right)^d||e_𝐤|1\}`$, $`e_𝐤=𝐤^2/2m\mu `$, and $`𝐪\{𝐤𝐩|𝐤,𝐩\}`$.
We are interested in the momentum distribution
$$a_{𝐤\sigma }^+a_{𝐤\sigma }=Tr[e^{\beta H}a_{𝐤\sigma }^+a_{𝐤\sigma }]/Tre^{\beta H}$$
$`(4.2)`$
and in the expectation value of the energy
$$H_{\mathrm{int}}=\underset{𝐪}{}\mathrm{\Lambda }(𝐪)$$
$`(4.3)`$
where
$$\mathrm{\Lambda }(𝐪)=\frac{\lambda }{L^{3d}}\underset{𝐤,𝐩}{}Tr[e^{\beta H}a_𝐤^+a_{𝐪𝐤}^+a_{𝐪𝐩}a_𝐩]/Tre^{\beta H}$$
$`(4.4)`$
By writing down the perturbation series for the partition function, rewriting it as a Grassmann integral
$`\frac{Tre^{\beta (H_0\lambda H_{\mathrm{int}})}}{Tre^{\beta H_0}}`$ $`=`$ $`{\displaystyle e^{\frac{\lambda }{(\beta L^d)^3}_{kpq}\overline{\psi }_k\overline{\psi }_{qk}\psi _{qp}\psi _p}𝑑\mu _C(\psi ,\overline{\psi })}`$ (4.5)
$`d\mu _C`$ $`=`$ $`\underset{k\sigma }{\Pi }\frac{\beta L^d}{ik_0e_𝐤}e^{\frac{1}{\beta L^d}_{k\sigma }(ik_0e_𝐤)\overline{\psi }_{k\sigma }\psi _{k\sigma }}\underset{k\sigma }{\Pi }d\psi _{k\sigma }d\overline{\psi }_{k\sigma },`$
performing a Hubbard-Stratonovich transformation ($`\varphi _q=u_q+iv_q`$, $`d\varphi _qd\overline{\varphi }_q:=du_qdv_q`$)
$$e^{_qa_qb_q}=e^{i_q(a_q\varphi _q+b_q\overline{\varphi }_q)}e^{_q|\varphi _q|^2}\underset{q}{\Pi }\frac{d\varphi _qd\overline{\varphi }_q}{\pi }$$
$`(4.6)`$
with
$$a_q=\frac{\lambda ^{\frac{1}{2}}}{(\beta L^d)^{\frac{3}{2}}}\underset{k}{}\overline{\psi }_k\overline{\psi }_{qk},b_q=\frac{\lambda ^{\frac{1}{2}}}{(\beta L^d)^{\frac{3}{2}}}\underset{p}{}\psi _p\psi _{qp}$$
$`(4.7)`$
and then integrating out the $`\psi ,\overline{\psi }`$ variables, one arrives at the following representation which is the starting point for our analysis (for more details, see \[FKT\] or \[L1\]):
$$\frac{1}{L^d}a_{𝐤\sigma }^+a_{𝐤\sigma }=\frac{1}{\beta L^d}\frac{1}{\beta }\underset{k_0\frac{\pi }{\beta }(2+1)}{}\psi _{𝐤k_0\sigma }^+\psi _{𝐤k_0\sigma }$$
$`(4.8)`$
where, abbreviating $`k=(𝐤,k_0)`$, $`\kappa =\beta L^d`$, $`a_k=ik_0e_𝐤`$, $`g=\lambda ^{\frac{1}{2}}`$,
$$\frac{1}{\kappa }\overline{\psi }_{t\sigma }\psi _{t\sigma }=\left[\begin{array}{cc}a_k\delta _{k,p}& \frac{ig}{\sqrt{\kappa }}\overline{\varphi }_{pk}\\ \frac{ig}{\sqrt{\kappa }}\varphi _{kp}& a_k\delta _{k,p}\end{array}\right]_{t\sigma ,t\sigma }^1𝑑P(\varphi )$$
$`(4.9)`$
and $`dP(\varphi )`$ is the normalized measure
$$dP(\varphi )=\frac{1}{Z}det\left[\begin{array}{cc}a_k\delta _{k,p}& \frac{ig}{\sqrt{\kappa }}\overline{\varphi }_{pk}\\ \frac{ig}{\sqrt{\kappa }}\varphi _{kp}& a_k\delta _{k,p}\end{array}\right]e^{_q|\varphi _q|^2}\underset{q}{\Pi }d\varphi _qd\overline{\varphi }_q$$
$`(4.10)`$
Furthermore
$$\mathrm{\Lambda }(𝐪)=\frac{1}{\beta }\underset{q_0\frac{2\pi }{\beta }}{}\mathrm{\Lambda }(𝐪,q_0)$$
$`(4.11)`$
where
$`\mathrm{\Lambda }(q)`$ $`=`$ $`\frac{\lambda }{(\beta L^d)^3}{\displaystyle \underset{k,p}{}}\overline{\psi }_k\overline{\psi }_{qk}\psi _{qp}\psi _p`$ (4.12)
$`=`$ $`|\varphi _q|^21`$
and the expectation in the last line is integration with respect to $`dP(\varphi )`$. The expectation on the $`\psi `$ variables $`\overline{\psi }_{k\sigma }\psi _{k\sigma }=\frac{1}{𝒵}\overline{\psi }_{k\sigma }\psi _{k\sigma }e^{\frac{\lambda }{\kappa ^3}_{k,p,q}\overline{\psi }_k\overline{\psi }_{qk}\psi _{qp}\psi _p}𝑑\mu _C`$ is Grassmann integration, but these representations are not used in the following. The matrix and the integral in (4.9) become finite dimensional if we choose some cutoff on the $`k_0`$ variables which is removed in the end. The set $``$ for the spatial momenta is already finite since we have chosen a fixed UV-cuttoff $`|e_𝐤|=|𝐤^2/2m\mu |1`$ which will not be removed in the end since we are interested in the infrared properties at $`𝐤^2/2m=\mu `$.
Our goal is to apply the inversion formula to the inverse matrix element in (4.9). Instead of writing the matrix in terms of four $`N\times N`$ blocks $`(a_k\delta _{k,p})_{k,p}`$, $`(\overline{\varphi }_{pk})_{k,p}`$, $`(\varphi _{kp})_{k,p}`$ and $`(a_k\delta _{k,p})_{k,p}`$ where $`N`$ is the number of the $`d+1`$-dimensional momenta $`k,p`$, we interchange rows and columns to rewrite it in terms of $`N`$ blocks of size $`2\times 2`$ (the matrix $`U`$ in the next line interchanges the rows and columns):
$$U\left[\begin{array}{cc}a_k\delta _{k,p}& \frac{ig}{\sqrt{\kappa }}\overline{\varphi }_{pk}\\ \frac{ig}{\sqrt{\kappa }}\varphi _{kp}& a_k\delta _{k,p}\end{array}\right]U^1=B=(B_{kp})_{k,p}$$
where the $`2\times 2`$ blocks $`B_{kp}`$ are given by
$$B_{kk}=\left(\begin{array}{cc}a_k& \frac{ig}{\sqrt{\kappa }}\overline{\varphi }_0\\ \frac{ig}{\sqrt{\kappa }}\varphi _0& a_k\end{array}\right),B_{kp}=\frac{ig}{\sqrt{\kappa }}\left(\begin{array}{cc}0& \overline{\varphi }_{pk}\\ \varphi _{kp}& 0\end{array}\right)\mathrm{if}kp.$$
$`(4.13)`$
We want to compute the $`2\times 2`$ matrix
$$G(k)=G(k)𝑑P(\varphi )$$
$`(4.14)`$
where
$$G(k)=[B^1]_{kk}$$
$`(4.15)`$
We start again with the two loop approximation which retains only the $`r=2`$ term in the denominator of (1.3). The result will be equation (4.20) below where the quantities $`\sigma _k`$ and $`|\varphi _0|^2`$ appearing in (4.20) have to satisfy the equations (4.21) and (4.24) which have to be solved in conjunction with (4.29). The solution to these equations is discussed below (4.30).
We first derive (4.20). In the two loop approximation,
$`G(k)`$ $``$ $`\left[B_{kk}{\displaystyle \underset{pk}{}}B_{kp}G_k(p)B_{pk}\right]^1`$ (4.22)
$`=`$ $`\left[\left(\begin{array}{cc}a_k& \frac{ig}{\sqrt{\kappa }}\overline{\varphi }_0\\ \frac{ig}{\sqrt{\kappa }}\varphi _0& a_k\end{array}\right)+\frac{\lambda }{\kappa }{\displaystyle \underset{pk}{}}\left(\begin{array}{cc}& \overline{\varphi }_{pk}\\ \varphi _{kp}& \end{array}\right)G_k(p)\left(\begin{array}{cc}& \overline{\varphi }_{kp}\\ \varphi _{pk}& \end{array}\right)\right]^1`$
$`=:`$ $`\left[\left(\begin{array}{cc}a_k& \frac{ig}{\sqrt{\kappa }}\overline{\varphi }_0\\ \frac{ig}{\sqrt{\kappa }}\varphi _0& \overline{a}_k\end{array}\right)+\mathrm{\Sigma }(k)\right]^1`$ (4.25)
where, substituting again $`G_k(p)`$ by $`G(p)`$ in the infinite volume limit,
$$\mathrm{\Sigma }(k)=\frac{\lambda }{\kappa }\underset{pk}{}\left(\begin{array}{cc}& \overline{\varphi }_{pk}\\ \varphi _{kp}& \end{array}\right)\left[\left(\begin{array}{cc}a_p& \frac{ig}{\sqrt{\kappa }}\overline{\varphi }_0\\ \frac{ig}{\sqrt{\kappa }}\varphi _0& \overline{a}_p\end{array}\right)+\mathrm{\Sigma }(p)\right]^1\left(\begin{array}{cc}& \overline{\varphi }_{kp}\\ \varphi _{pk}& \end{array}\right)$$
$`(4.17)`$
Anticipating the fact that the off-diagonal elements of $`\mathrm{\Sigma }(k)`$ will be zero (for ‘zero external field’), we make the Ansatz
$$\mathrm{\Sigma }(k)=\left(\begin{array}{cc}\sigma _k& \\ & \overline{\sigma }_k\end{array}\right)$$
$`(4.18)`$
and obtain
$$\left(\begin{array}{cc}\sigma _k& \\ & \overline{\sigma }_k\end{array}\right)=\frac{\lambda }{\kappa }\underset{pk}{}\frac{1}{(a_p+\sigma _p)(\overline{a}_p+\overline{\sigma }_p)+\frac{\lambda }{\kappa }|\varphi _0|^2}\left(\begin{array}{cc}(a_k+\sigma _k)|\varphi _{pk}|^2& \frac{ig}{\sqrt{\kappa }}\varphi _0\overline{\varphi }_{kp}\overline{\varphi }_{pk}\\ \frac{ig}{\sqrt{\kappa }}\overline{\varphi }_0\varphi _{kp}\varphi _{pk}& (\overline{a}_k+\overline{\sigma }_k)|\varphi _{kp}|^2\end{array}\right)$$
$`(4.19)`$
As for the Anderson model, we perform the functional integral by substituting the quantities $`|\varphi _q|^2`$ by their expectation values $`|\varphi _q|^2`$. Apparently this is less obvious in this case since $`dP(\varphi )`$ is no longer Gaussian and the $`|\varphi _q|^2`$ are no longer identically, independently distributed. We will comment on this after (4.37) below and at the end of the next section by reinterpreting this procedure as a resummation of diagrams. For now, we simply continue in this way. Then
$$G(k)=\frac{1}{|a_k+\sigma _k|^2+\frac{\lambda }{\kappa }|\varphi _0|^2}\left(\begin{array}{cc}\overline{a}_k+\overline{\sigma }_k& \frac{ig}{\sqrt{\kappa }}\overline{\varphi }_0\\ \frac{ig}{\sqrt{\kappa }}\varphi _0& a_k+\sigma _k\end{array}\right)$$
$`(4.20)`$
where the quantity $`\sigma _k`$ has to satisfy the equation
$$\sigma _k=\frac{\lambda }{\kappa }\underset{pk}{}\frac{\overline{a}_p+\overline{\sigma }_p}{|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2}|\varphi _{pk}|^2$$
$`(4.21)`$
Since $`dP(\varphi )`$ is not Gaussian, we do not know the expectations $`|\varphi _q|^2`$. However, by partial integration, we obtain
$$|\varphi _q|^2=1+\frac{ig}{\sqrt{\kappa }}\underset{p}{}\varphi _q[B^1(\varphi )]_{p,p+q}𝑑P(\varphi )$$
$`(4.22)`$
Namely,
$`|\varphi _q|^2`$ $`=`$ $`\frac{1}{Z}{\displaystyle \varphi _q\overline{\varphi }_qdet\left[\{B_{kp}(\varphi )\}_{k,p}\right]e^{_q|\varphi _q|^2}d\varphi _qd\overline{\varphi }_q}`$
$`=`$ $`1+\frac{1}{Z}{\displaystyle \varphi _q\left(\frac{}{\varphi _q}det\left[\{B_{kp}(\varphi )\}_{k,p}\right]\right)e^{_q|\varphi _q|^2}𝑑\varphi _q𝑑\overline{\varphi }_q}`$
$`=`$ $`1+\frac{1}{Z}{\displaystyle \varphi _q\underset{p,\tau }{}det\left[\begin{array}{ccc}|& |& |\\ B_{k\sigma ,p^{}\tau ^{}}& \frac{B_{k\sigma ,p\tau }}{\varphi _q}& B_{k\sigma ,p^{\prime \prime }\tau ^{\prime \prime }}\\ |& |& |\end{array}\right]e^{_q|\varphi _q|^2}d\varphi _qd\overline{\varphi }_q}`$
Since
$$\frac{}{\varphi _q}B_{kp}=\frac{ig}{\sqrt{\kappa }}\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)\delta _{kp,q}$$
we have
$$det\left[\begin{array}{ccc}|& |& |\\ B_{k\sigma ,p^{}\tau ^{}}& \frac{B_{k\sigma ,p\tau }}{\varphi _q}& B_{k\sigma ,p^{\prime \prime }\tau ^{\prime \prime }}\\ |& |& |\end{array}\right]/det\left[\{B_{kp}\}_{k,p}\right]=\{\begin{array}{cc}0\hfill & \mathrm{if}\tau =\hfill \\ & \\ \frac{ig}{\sqrt{\kappa }}[B^1]_{p,p+q}\hfill & \mathrm{if}\tau =\hfill \end{array}$$
which results in (4.22).
The inverse matrix element in (4.22) we compute again with (1.3,4) in the two loop approximation. Consider first the case $`q=0`$. Then one gets
$`|\varphi _0|^2`$ $`=`$ $`1+\frac{ig}{\sqrt{\kappa }}{\displaystyle \underset{p}{}}{\displaystyle \varphi _0G(p)_{}𝑑P(\varphi )}`$ (4.25)
$`=`$ $`1+\frac{ig}{\sqrt{\kappa }}{\displaystyle \underset{p}{}}{\displaystyle \varphi _0\frac{1}{|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2}\left(\begin{array}{cc}\overline{a}_p+\overline{\sigma }_p& \frac{ig}{\sqrt{\kappa }}\overline{\varphi }_0\\ \frac{ig}{\sqrt{\kappa }}\varphi _0& a_p+\sigma _p\end{array}\right)_{}𝑑P(\varphi )}`$
$`=`$ $`1+\frac{\lambda }{\kappa }{\displaystyle \underset{p}{}}{\displaystyle \varphi _0\frac{\overline{\varphi }_0}{|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2}𝑑P(\varphi )}`$ (4.26)
Performing the functional integral by substitution of expectation values gives
$$|\varphi _0|^2=1+\frac{\lambda }{\kappa }\underset{p}{}|\varphi _0|^2\frac{1}{|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2}$$
or
$$|\varphi _0|^2=\frac{1}{1\frac{\lambda }{\kappa }{\displaystyle \underset{p}{\overset{}{}}}\frac{1}{|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2}}$$
$`(4.24)`$
Before we discuss (4.24), we write down the equation for $`q0`$. In that case we use (1.4) to compute $`[B^1(\varphi )]_{p,p+q}`$ in the two loop approximation. We get
$`[B^1(\varphi )]_{p,p+q}\left[G(p)B_{p,p+q}G(p+q)\right]_{}`$
$`=\frac{1}{|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2}\frac{1}{|a_{p+q}+\sigma _{p+q}|^2+\frac{\lambda }{\kappa }|\varphi _0|^2}\frac{ig}{\sqrt{\kappa }}\times `$ (4.27)
$`\left(\begin{array}{cc}\frac{ig}{\sqrt{\kappa }}[(\overline{a}+\overline{\sigma })_{p+q}\overline{\varphi }_0\varphi _q+(\overline{a}+\overline{\sigma })_p\varphi _0\overline{\varphi }_q]& (\overline{a}+\overline{\sigma })_p(a+\sigma )_{p+q}\overline{\varphi }_q\frac{\lambda }{\kappa }\overline{\varphi }_0^2\varphi _q\\ (a+\sigma )_p(\overline{a}+\overline{\sigma })_{p+q}\varphi _q\frac{\lambda }{\kappa }\varphi _0^2\overline{\varphi }_q& \frac{ig}{\sqrt{\kappa }}[(a+\sigma )_{p+q}\varphi _0\overline{\varphi }_q+(a+\sigma )_p\overline{\varphi }_0\varphi _q]\end{array}\right)_{}`$
$`=\frac{ig}{\sqrt{\kappa }}\frac{(\overline{a}+\overline{\sigma })_p(a+\sigma )_{p+q}\overline{\varphi }_q\frac{\lambda }{\kappa }\overline{\varphi }_0^2\varphi _q}{\left(|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2\right)\left(|a_{p+q}+\sigma _{p+q}|^2+\frac{\lambda }{\kappa }|\varphi _0|^2\right)}`$ (4.28)
which gives
$`|\varphi _q|^2`$ $`=`$ $`1+\frac{\lambda }{\kappa }{\displaystyle \underset{p}{}}{\displaystyle \varphi _q\frac{(\overline{a}+\overline{\sigma })_p(a+\sigma )_{p+q}\overline{\varphi }_q\frac{\lambda }{\kappa }\overline{\varphi }_0^2\varphi _q}{\left(|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2\right)\left(|a_{p+q}+\sigma _{p+q}|^2+\frac{\lambda }{\kappa }|\varphi _0|^2\right)}𝑑P(\varphi )}`$ (4.26)
$`=`$ $`1+\frac{\lambda }{\kappa }{\displaystyle \underset{p}{}}\frac{(\overline{a}_p+\overline{\sigma }_p)(a_{p+q}+\sigma _{p+q})|\varphi _q|^2\frac{\lambda }{\kappa }\overline{\varphi }_0^2\varphi _q\varphi _q}{\left(|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2\right)\left(|a_{p+q}+\sigma _{p+q}|^2+\frac{\lambda }{\kappa }|\varphi _0|^2\right)}`$
Although one may think that the expectation $`\overline{\varphi }_0^2\varphi _q\varphi _q`$ vanishes for zero external field, this is not so. This can be seen again by partial integration:
$`\overline{\varphi }_0^2\varphi _q\varphi _q`$ $`=`$ $`\frac{1}{Z}{\displaystyle \overline{\varphi }_0^2\varphi _q\varphi _qdet\left[\{B_{kp}(\varphi )\}_{k,p}\right]e^{_q|\varphi _q|^2}d\varphi _qd\overline{\varphi }_q}`$
$`=`$ $`\frac{1}{Z}{\displaystyle \overline{\varphi }_0^2\varphi _q\left(\frac{}{\overline{\varphi }_q}det\left[\{B_{kp}(\varphi )\}_{k,p}\right]\right)e^{_q|\varphi _q|^2}𝑑\varphi _q𝑑\overline{\varphi }_q}`$
$`=`$ $`\frac{1}{Z}{\displaystyle \overline{\varphi }_0^2\varphi _q\underset{p,\tau }{}det\left[\begin{array}{ccc}|& |& |\\ B_{k\sigma ,p^{}\tau ^{}}& \frac{B_{k\sigma ,p\tau }}{\overline{\varphi }_q}& B_{k\sigma ,p^{\prime \prime }\tau ^{\prime \prime }}\\ |& |& |\end{array}\right]e^{_q|\varphi _q|^2}d\varphi _qd\overline{\varphi }_q}`$
The above determinant is multiplied and devided by $`det\left[\{B_{kp}\}_{k,p}\right]`$ to give
$$det\left[\begin{array}{ccc}|& |& |\\ B_{k\sigma ,p^{}\tau ^{}}& \frac{B_{k\sigma ,p\tau }}{\overline{\varphi }_q}& B_{k\sigma ,p^{\prime \prime }\tau ^{\prime \prime }}\\ |& |& |\end{array}\right]/det\left[\{B_{kp}\}_{k,p}\right]=\{\begin{array}{cc}0\hfill & \mathrm{if}\tau =\hfill \\ & \\ \frac{ig}{\sqrt{\kappa }}[B^1]_{p,p+q}\hfill & \mathrm{if}\tau =\hfill \end{array}$$
Computing the inverse matrix element again in the two loop approximation (4.25), we arrive at
$$\overline{\varphi }_0^2\varphi _q\varphi _q=\frac{\lambda }{\kappa }\underset{p}{}\frac{(a_p+\sigma _p)(\overline{a}_{p+q}+\sigma _{p+q})\overline{\varphi }_0^2\varphi _q\varphi _q\frac{\lambda }{\kappa }\overline{\varphi }_0^2\varphi _0^2\varphi _q\overline{\varphi }_q}{\left(|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2\right)\left(|a_{p+q}+\sigma _{p+q}|^2+\frac{\lambda }{\kappa }|\varphi _0|^2\right)}$$
Abbreviating
$$g_p=\frac{a_p+\sigma _p}{|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2},f_p=\frac{\sqrt{\frac{\lambda }{\kappa }|\varphi _0|^2}}{|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2}$$
$`(4.27)`$
this gives
$$\frac{\lambda }{\kappa }\overline{\varphi }_0^2\varphi _q\varphi _q=\frac{\lambda }{\kappa }\underset{p}{}g_p\overline{g}_{p+q}\frac{\lambda }{\kappa }\overline{\varphi }_0^2\varphi _q\varphi _q\frac{\lambda }{\kappa }\underset{p}{}f_pf_{p+q}\frac{\lambda }{\kappa }|\varphi _0|^2|\varphi _q|^2$$
or
$$\frac{\lambda }{\kappa }\overline{\varphi }_0^2\varphi _q\varphi _q=\frac{\frac{\lambda }{\kappa }\underset{p}{}f_pf_{p+q}\frac{\lambda }{\kappa }|\varphi _0|^2}{1\frac{\lambda }{\kappa }_pg_p\overline{g}_{p+q}}|\varphi _q|^2$$
$`(4.28)`$
Substituting this in (4.26), we finally arrive at
$$|\varphi _q|^2=\frac{1\frac{\lambda }{\kappa }\underset{p}{}g_p\overline{g}_{p+q}}{\left|1\frac{\lambda }{\kappa }_pg_p\overline{g}_{p+q}\right|^2\left(\frac{\lambda }{\kappa }_pf_pf_{p+q}\right)^2}$$
$`(4.29)`$
where $`g_p,f_p`$ are given by (4.27). Observe that, since $`dP(\varphi )`$ is complex, also $`|\varphi _q|^2`$ is in general complex. Only after summation over the $`q_0`$ variables we obtain necessarily a real quantity which is given by (4.4,11).
We now discuss the solutions to (4.24) and (4.29). We assume that the solution $`\sigma _k`$ of (4.21) is sufficiently small such that the BCS equation
$$\frac{\lambda }{\kappa }\underset{p}{\overset{}{}}\frac{1}{|a_p+\sigma _p|^2+|\mathrm{\Delta }|^2}=1$$
$`(4.30)`$
has a nonzero solution $`\mathrm{\Delta }0`$ (in particular this excludes large corrections like $`\sigma _pp_0^\alpha `$, $`\alpha 1/2`$, which one may expect in the case of Luttinger liquid behaviour, for $`d=1`$ one should make a seperate analysis), and make the Ansatz
$$\lambda |\varphi _0|^2=\beta L^d|\mathrm{\Delta }|^2+\eta $$
$`(4.31)`$
where $`\eta `$ is independent of the volume. Then
$`\frac{\lambda }{\kappa }{\displaystyle \underset{p}{\overset{}{}}}\frac{1}{|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2}`$ $`=`$ $`\frac{\lambda }{\kappa }{\displaystyle \underset{p}{\overset{}{}}}\frac{1}{|a_p+\sigma _p|^2+|\mathrm{\Delta }|^2+\frac{\eta }{\kappa }}`$ (4.32)
$`=`$ $`\frac{\lambda }{\kappa }{\displaystyle \underset{p}{\overset{}{}}}\frac{1}{|a_p+\sigma _p|^2+|\mathrm{\Delta }|^2}\frac{\lambda }{\kappa }{\displaystyle \underset{p}{\overset{}{}}}\frac{\eta /\kappa }{(|a_p+\sigma _p|^2+|\mathrm{\Delta }|^2)^2}+O\left((\frac{\eta }{\kappa })^2\right)`$
$`=`$ $`1c_\mathrm{\Delta }\frac{\eta }{\kappa }+O\left((\frac{\eta }{\kappa })^2\right)`$
where we put $`c_\mathrm{\Delta }=\frac{\lambda }{\kappa }_p\frac{1}{(|a_p+\sigma _p|^2+|\mathrm{\Delta }|^2)^2}`$ and used the BCS equation (4.30) in the last line. Equation (4.24) becomes
$`\kappa |\mathrm{\Delta }|^2+\eta `$ $`=`$ $`{\displaystyle \frac{\lambda }{c_\mathrm{\Delta }\frac{\eta }{\kappa }+O\left((\frac{\eta }{\kappa })^2\right)}}=\kappa {\displaystyle \frac{\lambda }{c_\mathrm{\Delta }\eta }}+O(1)`$
and has a solution $`\eta =\lambda /(c_\mathrm{\Delta }|\mathrm{\Delta }|^2)`$.
Now consider $`|\varphi _q|^2`$ for small but nonzero $`q`$. In the limit $`q0`$ the denominator in (4.29) vanishes, or more precisely, is of order $`O(1/\kappa )`$ since
$$1\frac{\lambda }{\kappa }\underset{p}{}g_p\overline{g}_p\frac{\lambda }{\kappa }\underset{p}{}f_pf_p=\mathrm{\hspace{0.33em}1}\frac{\lambda }{\kappa }\underset{p}{}\frac{1}{|a_p+\sigma _p|^2+\frac{\lambda }{\kappa }|\varphi _0|^2}=O(1/\kappa )$$
because of (4.32). If we assume the second derivatives of $`\sigma _k`$ to be integrable (which should be the case for $`d=3`$ and $`|\varphi _q|^21/q^2`$ by virtue of (4.21)), then, since the denominator in (4.29) is an even function of $`q`$, the small $`q`$ behaviour of $`|\varphi _q|^2`$ is $`1/q^2`$. This agrees with the common expectations \[FMRT,CFS,B\]. Usually the behaviour of $`|\varphi _q|^2`$ is infered from the second order Taylor expansion of the effective potential
$$V_{\mathrm{eff}}(\{\varphi _q\})=\underset{q}{}|\varphi _q|^2\mathrm{log}det\left[\begin{array}{cc}\delta _{k,p}& \frac{ig}{\sqrt{\kappa }}\frac{\overline{\varphi }_{pk}}{a_k}\\ \frac{ig}{\sqrt{\kappa }}\frac{\varphi _{kp}}{a_k}& \delta _{k,p}\end{array}\right]$$
$`(4.33)`$
around its global minimum \[L2\]
$$\varphi _q^{\mathrm{min}}=\sqrt{\beta L^d}\frac{|\mathrm{\Delta }|}{\sqrt{\lambda }}\delta _{q,0}e^{i\theta _0}$$
$`(4.34)`$
where the phase $`\theta _0`$ of $`\varphi _0`$ is arbitrary. If one expands $`V_{\mathrm{eff}}`$ up to second order in
$$\xi _q=\varphi _q\delta _{q,0}\sqrt{\beta L^d}\frac{|\mathrm{\Delta }|}{\sqrt{\lambda }}e^{i\theta _0}=\{\begin{array}{cc}\left(\rho _0\sqrt{\beta L^d}\frac{|\mathrm{\Delta }|}{\sqrt{\lambda }}\right)e^{i\theta _0}\hfill & \text{for }q=0\text{ }\hfill \\ \rho _qe^{i\theta _q}\hfill & \text{for }q0\hfill \end{array}$$
$`(4.35)`$
one obtains \[L2\]
$`V_{\mathrm{eff}}(\{\varphi _q\})`$ $`=`$ $`V_{\mathrm{min}}+2\beta _0(\rho _0\sqrt{\beta L^d}\frac{|\mathrm{\Delta }|}{\sqrt{\lambda }})^2+{\displaystyle \underset{q0}{}}(\alpha _q+i\gamma _q)\rho _q^2`$ (4.36)
$`+\frac{1}{2}{\displaystyle \underset{q0}{}}\beta _q|e^{i\theta _0}\varphi _q+e^{i\theta _0}\overline{\varphi }_q|^2+O(\xi ^3)`$
where for small $`q`$ one has $`\alpha _q,\gamma _qq^2`$. Hence, if $`V_{\mathrm{eff}}`$ is substituted by the right hand side of (4.36) one obtains $`|\varphi _q|^21/q^2`$.
For $`d=3`$, this seems to be the right answer, but in lower dimensions one would expect an integrable singularity due to (4.21) and (4.3,4,11). In particular, we think it would be a very interesting problem to solve the integral equations (4.21,24,29) for $`d=1`$ and to check the result for Luttinger liquid behaviour. A good warm up excercise would be to consider the $`0+1`$ dimensional problem, that is, we only have the $`k_0,p_0,q_0`$-variables. In that case the ‘bare BCS equation’
$$\frac{\lambda }{\beta }\underset{p_0\frac{\pi }{\beta }(2+1)}{}\frac{1}{p_0^2+|\mathrm{\Delta }|^2}=1$$
still has a nonzero solution $`\mathrm{\Delta }`$ for sufficiently small $`T=1/\beta `$ and the question would be whether the correction $`\sigma _{p_0}`$ is sufficiently big to destroy the gap. That is, does the ‘renormalized BCS equation’
$$\frac{\lambda }{\beta }\underset{p_0\frac{\pi }{\beta }(2+1)}{}\frac{1}{|p_0+\sigma _{p_0}|^2+|\mathrm{\Delta }|^2}=1$$
$`\sigma _{p_0}`$ being the solution to (4.21,24,29), still have a nonzero solution? We remark that, if the gap vanishes (for arbitrary dimension), then also the singularity of $`|\varphi _q|^2`$ disappears. Namely, if the gap equation has no solution, that is, if $`\frac{1}{\kappa }_p\frac{1}{|a_p+\sigma _p|^2}<\mathrm{}`$, then $`|\varphi _0|^2`$ given by (4.24) is no longer macroscopic (for sufficiently small coupling) and $`\frac{\lambda }{\kappa }|\varphi _0|^2`$ vanishes in the infinite volume limit. And the denominator in (4.29) becomes for $`q0`$
$$1\frac{\lambda }{\kappa }\underset{p}{}\frac{1}{|a_p+\sigma _p|^2}$$
which would be nonzero (for sufficiently small coupling).
Finally we argue why it is reasonable to substitute $`|\varphi _0|^2`$ by its expectation value while performing the functional integral. We may write the effective potential (4.33) as
$$V_{\mathrm{eff}}(\{\varphi _q\})=V_1(\varphi _0)+V_2(\{\varphi _q\})$$
$`(4.37)`$
where
$`V_1(\varphi _0)`$ $`=`$ $`|\varphi _0|^2{\displaystyle \underset{k}{}}\mathrm{log}\left[1+\frac{\lambda }{\kappa }\frac{|\varphi _0|^2}{k_0^2+e_𝐤^2}\right]`$ (4.38)
$`=`$ $`\kappa \left(\frac{|\varphi _0|^2}{\kappa }\frac{1}{\kappa }{\displaystyle \underset{k}{}}\mathrm{log}\left[1+\frac{\lambda \frac{|\varphi _0|^2}{\kappa }}{k_0^2+e_𝐤^2}\right]\right)\kappa V_{\mathrm{BCS}}\left(\frac{|\varphi _0|}{\sqrt{\kappa }}\right)`$
and
$$V_2(\{\varphi _q\})=\underset{q0}{}|\varphi _q|^2\mathrm{log}det\left[\left(\begin{array}{cc}\delta _{k,p}& \frac{ig}{\sqrt{\kappa }}\frac{\overline{\varphi }_0}{a_k}\delta _{k,p}\\ \frac{ig}{\sqrt{\kappa }}\frac{\varphi _0}{a_k}\delta _{k,p}& \delta _{k,p}\end{array}\right)^1\left(\begin{array}{cc}\delta _{k,p}& \frac{ig}{\sqrt{\kappa }}\frac{\overline{\varphi }_{pk}}{a_k}\\ \frac{ig}{\sqrt{\kappa }}\frac{\varphi _{kp}}{a_k}& \delta _{k,p}\end{array}\right)\right]$$
$`(4.39)`$
If we ignore the $`\varphi _0`$-dependence of $`V_2`$, then the $`\varphi _0`$-integral
$`{\displaystyle \frac{F\left(\frac{1}{\kappa }|\varphi _0|^2\right)e^{V_1(\varphi _0)}𝑑\varphi _0𝑑\overline{\varphi }_0}{e^{V_1(\varphi _0)}𝑑\varphi _0𝑑\overline{\varphi }_0}}`$ $`=`$ $`{\displaystyle \frac{F\left(\rho ^2\right)e^{\kappa V_{\mathrm{BCS}}(\rho )}\rho 𝑑\rho }{e^{\kappa V_{\mathrm{BCS}}(\rho )}\rho 𝑑\rho }}\stackrel{\kappa \mathrm{}}{}F(\rho _{\mathrm{min}}^2)=F\left(\frac{1}{\kappa }|\varphi _0|^2\right)`$
simply puts $`|\varphi _0|^2`$ at the global minimum of the (BCS) effective potential.
### 4.2 The $`\phi ^4`$-Model
In this section we choose the $`\phi ^4`$-model as a typical bosonic model to demonstrate our method. As in section 2, we start in finite volume $`[0,L]^d`$ on a lattice with lattice spacing $`1/M`$. The two point function is given by
$$S(x,y)=\phi _x\phi _y:=\frac{_{^{N^d}}\phi _x\phi _ye^{\frac{g^2}{2}\frac{1}{M^d}_x\phi _x^4}e^{\frac{1}{M^{2d}}_{x,y}(\mathrm{\Delta }+m^2)_{x,y}\phi _x\phi _y}\underset{x}{\Pi }d\phi _x}{_{^{N^d}}e^{\frac{g^2}{2}\frac{1}{M^d}_x\phi _x^4}e^{\frac{1}{M^{2d}}_{x,y}(\mathrm{\Delta }+m^2)_{x,y}\phi _x\phi _y}\Pi _xd\phi _x}$$
$`(4.41)`$
where
$$(\mathrm{\Delta }+m^2)_{x,y}=M^d\left[M^2\underset{i=1}{\overset{d}{}}(\delta _{x,ye_i/M}+\delta _{x,y+e_i/M}2\delta _{x,y})+m^2\delta _{x,y}\right]$$
$`(4.42)`$
First we have to bring this into the form $`[P+Q]_{x,y}^1𝑑\mu `$, $`P`$ diagonal in momentum space, $`Q`$ diagonal in coordinate space. This is done again by making a Hubbard Stratonovich transformation which in this case reads
$$e^{\frac{1}{2}_xa_x^2}=e^{i_xa_xu_x}e^{\frac{1}{2}_xu_x^2}\underset{x}{\Pi }\frac{du_x}{\sqrt{2\pi }}$$
$`(4.43)`$
with
$$a_x=\frac{g}{\sqrt{M^d}}\phi _x^2$$
$`(4.44)`$
The result is Gaussian in the $`\phi _x`$-variables and the integral over these variables gives
$$S(x,y)=_{^{N^d}}\left[\frac{1}{M^{2d}}(\mathrm{\Delta }+m^2)_{x,y}\frac{ig}{\sqrt{M^d}}u_x\delta _{x,y}\right]_{x,y}^1𝑑P(u)$$
$`(4.45)`$
where
$$dP(u)=\frac{1}{Z}det\left[\frac{1}{M^{2d}}(\mathrm{\Delta }+m^2)_{x,y}\frac{ig}{\sqrt{M^d}}u_x\delta _{x,y}\right]^{\frac{1}{2}}e^{\frac{1}{2}_xu_x^2}\underset{x}{\Pi }du_x$$
$`(4.46)`$
Since we have bosons, the determinant comes with a power of $`1/2`$ which is the only difference compared to a fermionic system. In momentum space this reads (compare equations (2.7-11))
$$S(xy)=\frac{1}{L^d}\underset{k}{}e^{ik(xy)}G(k)$$
$`(4.47)`$
where ($`\gamma _q=v_q+iw_q`$, $`\gamma _q=\overline{\gamma }_q`$, $`d\gamma _qd\overline{\gamma }_q:=dv_qdw_q`$)
$$G(k)=_{^{N^d}}\left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]_{kk}^1𝑑P(\gamma )$$
$`(4.48)`$
and
$$dP(\gamma )=\frac{1}{Z}det\left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]^{\frac{1}{2}}e^{\frac{1}{2}v_0^2}dv_0\underset{q^+}{\Pi }e^{|\gamma _q|^2}d\gamma _qd\overline{\gamma }_q$$
$`(4.49)`$
and $`^+`$ again is a set such that either $`q^+`$ or $`q^+`$. Furthermore
$$a_k=4M^2\underset{i=1}{\overset{d}{}}\mathrm{sin}^2\left[\frac{k_i}{2M}\right]+m^2$$
$`(4.50)`$
Equation (4.48) is our starting point. We apply (1.3) to the inverse matrix element in (4.48). In the two loop approximation one obtains $`(\gamma _0=v_0)`$
$$[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}]_{kk}^1\frac{1}{a_k\frac{igv_0}{\sqrt{L^d}}+\frac{g^2}{L^d}_{pk}G_k(p)|\gamma _{kp}|^2}=:\frac{1}{a_k+\sigma _k}$$
$`(4.51)`$
where
$$\sigma _k=\frac{ig}{\sqrt{L^d}}v_0+\frac{g^2}{L^d}\underset{pk}{}\frac{|\gamma _{kp}|^2}{a_p\frac{igv_0}{\sqrt{L^d}}+\sigma _p}$$
$`(4.52)`$
which results in
$$G(k)=\frac{1}{a_k+\sigma _k}$$
$`(4.53)`$
where $`\sigma _k`$ has to satisfy the equation
$`\sigma _k`$ $`=`$ $`\frac{ig}{\sqrt{L^d}}v_0+\frac{g^2}{L^d}{\displaystyle \underset{pk}{}}{\displaystyle \frac{|\gamma _{kp}|^2}{a_p+\sigma _p}}`$ (4.54)
$`=`$ $`\frac{g^2}{2L^d}{\displaystyle \underset{p}{}}G(p)+\frac{g^2}{L^d}{\displaystyle \underset{pk}{}}{\displaystyle \frac{|\gamma _{kp}|^2}{a_p+\sigma _p}}=\frac{g^2}{L^d}{\displaystyle \underset{pk}{}}{\displaystyle \frac{|\gamma _{kp}|^2+\frac{1}{2}}{a_p+\sigma _p}}`$
where the last line is due to
$`v_0`$ $`=`$ $`\frac{1}{Z}{\displaystyle v_0det\left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]^{\frac{1}{2}}e^{\frac{1}{2}v_0^2}dv_0\underset{q^+}{\Pi }e^{|\gamma _q|^2}d\gamma _qd\overline{\gamma }_q}`$ (4.55)
$`=`$ $`\frac{1}{Z}{\displaystyle \left\{\frac{}{v_0}det\left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]^{\frac{1}{2}}\right\}e^{\frac{1}{2}v_0^2}𝑑v_0\underset{q^+}{\Pi }e^{|\gamma _q|^2}d\gamma _qd\overline{\gamma }_q}`$
$`=`$ $`\frac{1}{2}{\displaystyle \underset{p}{}}(\frac{ig}{\sqrt{L^d}}){\displaystyle \left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]_{pp}^1𝑑P(\gamma )}`$
As for the Many-Electron system, we can derive an equation for $`|\gamma _q|^2`$ by partial integration:
$`|\gamma _q|^2`$ $`=`$ $`\frac{1}{Z}{\displaystyle \gamma _q\overline{\gamma }_qdet\left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]^{\frac{1}{2}}e^{\frac{v_0^2}{2}}dv_0\underset{q}{\Pi }e^{|\gamma _q|^2}d\gamma _qd\overline{\gamma }_q}`$ (4.56)
$`=`$ $`\mathrm{\hspace{0.33em}1}+\frac{1}{Z}{\displaystyle \gamma _q\frac{}{\gamma _q}\left\{det\left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]^{\frac{1}{2}}\right\}e^{\frac{v_0^2}{2}}𝑑v_0\underset{q}{\Pi }e^{\frac{1}{2}|\gamma _q|^2}d\gamma _qd\overline{\gamma }_q}`$
$`=`$ $`\mathrm{\hspace{0.33em}1}\frac{1}{2}{\displaystyle \gamma _q\frac{\frac{}{\gamma _q}det\left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]}{det\left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]}𝑑P(\gamma )}`$
$`=`$ $`\mathrm{\hspace{0.33em}1}\frac{1}{2}{\displaystyle \underset{p}{}}\frac{ig}{\sqrt{L^d}}{\displaystyle \gamma _q\left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]_{p,p+q}^1𝑑P(\gamma )}`$
Computing the inverse matrix element in (4.56) again in the two loop approximation, one arrives at
$$|\gamma _q|^2=1|\gamma _q|^2\frac{g^2}{2L^d}\underset{p}{}\frac{1}{(a_p+\sigma _p)(a_{p+q}+\sigma _{p+q})}$$
or
$$|\gamma _q|^2=\frac{1}{1+\frac{g^2}{2}_{[0,2\pi M]^d}\frac{d^dp}{(2\pi )^d}\frac{1}{(a_p+\sigma _p)(a_{p+q}+\sigma _{p+q})}}$$
$`(4.57)`$
which has to be solved in conjunction with
$$\sigma _k=g^2_{[0,2\pi M]^d}\frac{d^dp}{(2\pi )^d}\frac{|\gamma _{kp}|^2+\frac{1}{2}}{a_p+\sigma _p}$$
$`(4.58)`$
Introducing the rescaled quantities
$$\sigma _k=M^2s_{\frac{p}{M}},|\gamma _q|^2=\lambda _{\frac{q}{M}},a_k=M^2\epsilon _{\frac{k}{M}},\epsilon _k=\underset{i=1}{\overset{d}{}}\mathrm{sin}^2\frac{k_i}{2}+\frac{m^2}{M^2}$$
$`(4.60)`$
(4.57,58) read
$`s_k`$ $`=`$ $`M^{d4}g^2_{[0,2\pi ]^d}\frac{d^dp}{(2\pi )^d}\frac{\lambda _{kp}+\frac{1}{2}}{\epsilon _p+s_p}`$ (4.61)
$`\lambda _q`$ $`=`$ $`{\displaystyle \frac{1}{1+M^{d4}\frac{g^2}{2}_{[0,2\pi ]^d}\frac{d^dp}{(2\pi )^d}\frac{1}{(\epsilon _p+s_p)(\epsilon _{p+q}+s_{p+q})}}}`$ (4.62)
Unfortunately we cannot check this result with the rigorously proven triviality theorem since $`\sigma _k`$ and $`|\gamma _q|^2`$ only give information on the 2-point function $`S(x,y)`$, (4.41), and on $`\frac{g^2}{M^d}_x\phi (x)^4=_q\mathrm{\Lambda }(q)`$ where $`\mathrm{\Lambda }(q)=|\gamma _q|^21`$. However, the triviality theorem \[F,FFS\] makes a statement on the connected 4-point function $`S_{4,c}(x_1,x_2,x_3,x_4)`$ at noncoinciding arguments, namely that this function vanishes in the continuum limit in dimension $`d>4`$.
Before we include the higher loop terms of (1.3,4) and give an interpretation in terms of diagrams, we would like to comment shortly on a problem which was suggested to us by A. Sokal after a preprint of this paper was published on the web. It refers to the $`\varphi ^2\psi ^2`$\- or $`\varphi _1^2\varphi _2^2`$-model. That is, we have two scalar bosonic fields on a lattice with unit lattice spacing with action
$$𝒮(\varphi _1,\varphi _2)=\underset{i=1}{\overset{2}{}}\underset{x}{}\varphi _i(x)(\mathrm{\Delta }+m^2)\varphi _i(x)+\lambda \underset{x}{}\varphi _1(x)^2\varphi _2(x)^2$$
The question is whether there is exponential decay (or a gap in momentum space) for the two point function $`G(x,y)=\varphi _1(x)\varphi _1(y)e^{𝒮(\varphi )}/e^{𝒮(\varphi )}`$ in the zero mass $`m0`$ limit. A computation with the above formalism in two loop approximation gives $`G(k)=\frac{1}{k^2+\sigma }`$ where the gap $`\sigma `$ has to satisfy the equation $`\sigma =\lambda _{[\pi ,\pi ]^d}\frac{d^dp}{(2\pi )^d}\frac{1}{p^2+\sigma }`$ which gives
$$\sigma =\{\begin{array}{cc}\hfill O(\lambda )& \hfill \mathrm{if}d3\\ \hfill O\left(\lambda \mathrm{log}[1/\lambda ]\right)& \hfill \mathrm{if}d=2\\ \hfill O(\lambda ^{\frac{2}{3}})& \hfill \mathrm{if}d=1\end{array}$$
We now include the higher loop terms of (1.3,4) and give an interpretation in terms of diagrams. The exact equations for $`G(k)`$ and $`|\gamma _q|^2`$ are
$`G(k)`$ $`=`$ $`{\displaystyle \left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]_{kk}^1𝑑P(\gamma )}=\frac{1}{a_k+\sigma _k}`$ (4.63)
$`\sigma _k`$ $`=`$ $`\frac{ig}{\sqrt{L^d}}v_0+{\displaystyle \underset{r=2}{\overset{N^d}{}}}\left(\frac{ig}{\sqrt{L^d}}\right)^r{\displaystyle \underset{\genfrac{}{}{0pt}{}{p_2\mathrm{}p_rk}{p_ip_j}}{}}G_k(p_2)\mathrm{}G_{kp_2\mathrm{}p_{r1}}(p_r)\gamma _{kp_2}\gamma _{p_2p_3}\mathrm{}\gamma _{p_rk}`$
and
$`|\gamma _q|^2`$ $`=`$ $`1+\frac{ig}{2\sqrt{L^d}}{\displaystyle \underset{p}{}}{\displaystyle \gamma _q\left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]_{p,p+q}^1𝑑P(\gamma )}`$
$`\stackrel{pp_2}{=}`$ $`1+\frac{1}{2}{\displaystyle \underset{r=2}{\overset{N^d}{}}}\left(\frac{ig}{\sqrt{L^d}}\right)^r{\displaystyle \underset{\genfrac{}{}{0pt}{}{p_2\mathrm{}p_rp_2+q}{p_ip_j}}{}}G(p_2)G_{p_2}(p_3)\mathrm{}G_{p_2\mathrm{}p_{r1}}(p_r)G_{p_2\mathrm{}p_r}(p_2+q)\times `$
$`\gamma _{p_2p_3}\mathrm{}\gamma _{p_{r1}p_r}\gamma _{p_rp_2q}\gamma _{p_2+qp_2}`$
For $`r>2`$, we obtain terms $`\gamma _{k_1}\mathrm{}\gamma _{k_r}`$ whose connected contributions
are, in terms of the electron or $`\phi ^4`$-lines, are at least six-legged. Since for the many-electron system and for the $`\phi ^4`$-model (for $`d=4`$) the relevant diagrams are two- and four-legged \[FT,R\], one may start with an approximation which ignores the connected $`r`$-loop contributions for $`r>2`$. This is obtained by writing
$$\gamma _{k_1}\mathrm{}\gamma _{k_n}\gamma _{k_1}\mathrm{}\gamma _{k_n}_2$$
$`(4.65)`$
where (the index ‘2’ for ‘retaining only two-loop contributions’)
$$\gamma _{k_1}\mathrm{}\gamma _{k_{2n}}_2:=\underset{\mathrm{pairings}\sigma }{}\gamma _{k_{\sigma 1}}\gamma _{k_{\sigma 2}}\mathrm{}\gamma _{k_{\sigma (2n1)}}\gamma _{k_{\sigma 2n}}=\gamma _{k_1}\mathrm{}\gamma _{k_{2n}}𝑑P_2(\gamma )$$
$`(4.66)`$
if we define
$$dP_2(\gamma ):=e^{_q\frac{|\gamma _q|^2}{|\gamma _q|^2}}\underset{q}{\Pi }\frac{d\gamma _qd\overline{\gamma }_q}{\pi |\gamma _q|^2}$$
$`(4.67)`$
Substituting $`dP`$ by $`dP_2`$ in (4.63,64), we obtain a model which differs from the original model only by irrelevant contributions and for which we are able to write down a closed set of equations for the two-legged particle correlation function $`G(k)`$ and the two-legged squiggle correlation function $`|\gamma _q|^2`$ by resumming all two-legged (particle and squiggle) subdiagrams. The exact equations of this model are
$`G(k)`$ $`=`$ $`{\displaystyle \left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]_{kk}^1𝑑P_2(\gamma )}`$ (4.68)
$`|\gamma _q|^2`$ $`=`$ $`1+\frac{ig}{2\sqrt{L^d}}{\displaystyle \underset{p}{}}{\displaystyle \gamma _q\left[a_k\delta _{k,p}\frac{ig}{\sqrt{L^d}}\gamma _{kp}\right]_{p,p+q}^1𝑑P_2(\gamma )}`$ (4.69)
and the resummation of the two-legged particle and squiggle subdiagrams is obtained by applying the inversion formula (1.3,4) to the inverse matrix elements in (4.68,69). A discussion similar to those of section 2 gives the following closed set of equations for the quantities $`G(k)`$ and $`|\gamma _q|^2`$:
$$G(k)=\frac{1}{a_k+\sigma _k},|\gamma _q|^2=\frac{1}{1+\pi _q}$$
$`(4.70)`$
where
$`\sigma _k`$ $`=`$ $`\frac{g^2}{2L^d}{\displaystyle \underset{p}{}}G(p)+{\displaystyle \underset{r=2}{\overset{\mathrm{}}{}}}\left(\frac{ig}{\sqrt{L^d}}\right)^r{\displaystyle \underset{\genfrac{}{}{0pt}{}{p_2\mathrm{}p_rk}{p_ip_j}}{}}G(p_2)\mathrm{}G(p_r)\gamma _{kp_2}\gamma _{p_2p_3}\mathrm{}\gamma _{p_rk}_2`$
$`\pi _q`$ $`=`$ $`\frac{1}{2}{\displaystyle \underset{r=2}{\overset{\mathrm{}}{}}}\left(\frac{ig}{\sqrt{L^d}}\right)^r{\displaystyle \underset{s=3}{\overset{r1}{}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{p_2\mathrm{}p_rp_2+q}{p_ip_j}}{}}(\delta _{q,p_{s+1}p_s}G(p_2)\mathrm{}G(p_r)G(p_2+q)\times `$
$`\gamma _{p_2p_3}\mathrm{}\widehat{\gamma }_{p_sp_{s+1}}\mathrm{}\gamma _{p_{r1}p_r}\gamma _{p_rp_2q}_2)`$
In the last line we used that $`\gamma _q`$ in (4.64) cannot contract to $`\gamma _{p_2p_3}`$ or to $`\gamma _{p_rp_2q}`$. If the expectations of the $`\gamma `$-fields on the right hand side of (4.71,72) are computed according to (4.66), one obtains the expansion into diagrams. The graphs contributing to $`\sigma _k`$ have exactly one string of particle lines, each line having $`G`$ as propagator, and no particle loops (up to the tadpole diagram). Each squiggle corresponds to a factor $`|\gamma |^2`$. The diagrams contributing to $`\pi `$ have exactly one particle loop, the propagators being again the interacting two point functions, $`G`$ for the particle lines and $`|\gamma |^2`$ for the squiggles. In both cases there are no two-legged subdiagrams. However, although the equation $`|\gamma _q|^2=\frac{1}{1+\pi _q}`$ resums ladder or bubble diagrams (which is apparent from (4.57) or (4.26)) and more general four-legged particle subdiagrams if the terms for $`r4`$ in (4.72) are taken into account, the right hand side of (4.71,72) still contains diagrams with four-legged particle subdiagrams. Thus, the resummation of four-legged particle subdiagrams is only partially through the complete resummation of two-legged squiggle diagrams. Also observe that, in going from (4.68,69) to (4.70-72), we cut off the $`r`$-sum at some fixed order $`\mathrm{}`$ independent of the volume since we can only expect that the expansions are asymptotic ones, compare the discussion in section 2.
## 5 Concluding Remarks
In the general case, without making the approximation (4.65), we expect the following picture for a generic quartic field theoretical model. Let $`G`$ and $`G_0`$ be the interacting and free particle Greens function (one solid line goes in, one solid line goes out), and let $`D`$ and $`D_0`$ be the interacting and free interaction Greens function (one wavy line goes in, one wavy line goes out). Then we expect the following closed set of integral equations for $`G`$ and $`D`$:
$$G=\frac{1}{G_0^1+\sigma (G,D)},D=\frac{1}{D_0^1+\pi (G,D)}$$
$`(5.1)`$
where $`\sigma `$ and $`\pi `$ are the sum of all two legged diagrams without two legged (particle and wavy line) subdiagrams with propagators $`G`$ and $`D`$ (instead of $`G_0`$, $`D_0`$). Thus (5.1) simply eliminates all two legged insertions by substituting them by the full propagators. For the Anderson model $`D=D_0=1`$ and (5.1) reduces to (2.27,35).
A variant of equations (5.1) has been derived on a more heuristic level in \[CJT\] and \[LW\]. Their integral equation (for example equation (40) of \[LW\]) reads
$$G=\frac{1}{G_0^1+\stackrel{~}{\sigma }(G,D_0)}$$
$`(5.2)`$
where $`\stackrel{~}{\sigma }`$ is the sum of all two legged diagrams without two legged particle insertions, with propagators $`G`$ and $`D_0`$. Thus this equation does not resum two legged interaction subgraphs (one wavy line goes in, one wavy line goes out). However resummation of these diagrams corresponds to a partial resummation of four legged particle subgraphs (for example the second equation in (5.4) below resums bubble diagrams), and is necessary in order to get the right behaviour, in particular for the many-electron system.
Another popular way of eliminating two legged subdiagrams (instead of using integral equations) is the use of counterterms. The underlying combinatorial identity is the following one. Let
$$𝒮(\psi ,\overline{\psi })=𝑑k\overline{\psi }_kG_0^1(k)\psi _k+𝒮_{\mathrm{int}}(\psi ,\overline{\psi })$$
$`(5.3)`$
be some action of a field theoretical model and let $`T(k)=T(G_0)(k)`$ be the sum of all amputated two legged particle diagrams without two legged particle subdiagrams, evaluated with the bare propagator $`G_0`$. Let $`\delta 𝒮(\psi ,\overline{\psi })=𝑑k\overline{\psi }_kT(k)\psi _k`$. Consider the model with action $`𝒮\delta 𝒮`$. Then a $`p`$-point function of that model is given by the sum of all $`p`$-legged diagrams which do not contain any two legged particle subdiagrams, evaluated with the bare propagator $`G_0`$. In particular, by construction, the two point function of that model is exactly given by $`G_0`$. Now, since the quadratic part of the model under consideration (given by the action $`𝒮\delta 𝒮`$) should be given by the bare Greens function $`G_0^1`$ and the interacting Greens function is $`G`$, one is led to the equation $`G^1T(G)=G_0^1`$ which coincides with (5.2).
Since the quantities $`\sigma `$ and $`\pi `$ in (5.1) are not explicitely given but merely are given by a sum of diagrams, we have to make an approximation in order to get a concrete set of integral equations which we can deal with. That is, we substitute $`\sigma `$ and $`\pi `$ by its lowest order contributions which leads to the system
$$G(k)=\frac{1}{G_0(k)^1+𝑑pD(p)G(kp)},D(q)=\frac{1}{D_0(q)^1+𝑑pG(p)G(p+q)}$$
$`(5.4)`$
This corresponds to the use of (1.3,4) retaining only the $`r=2`$ term. Thus we assume that the expansions for $`\sigma `$ and $`\pi `$ are asymptotic. A rigorous proof of that is of course a very difficult mathematical problem and this has not been adressed in this paper. Roughly one may expect this if each diagram contributing to $`\sigma `$ and $`\pi `$ allows a $`const^n`$ bound (no $`n!`$ and of course no divergent contributions). One may look in \[FKLT\] for an outline of proof for the many electron system with an anisotropic dispersion relation. In that case actually one obtains a series with a small positive radius of convergence instead of only an asymptotic one (because the model is fermionic), which simplifies the proof considerably.
Finally we remark that the equations (5.4) can be found in the literature. Usually they are derived from the Schwinger-Dyson equations which is the following non closed set of two equations for the three unknown functions $`G,D`$ and $`\mathrm{\Gamma }`$, $`\mathrm{\Gamma }`$ being the vertex function (see, for example, \[AGD\]):
$`G(k)`$ $`=`$ $`G_0(k)+G_0(k){\displaystyle 𝑑pG(p)D(kp)\mathrm{\Gamma }(p,kp)G(k)}`$
$`D(q)`$ $`=`$ $`D_0(q)+D_0(q){\displaystyle 𝑑pG(p)G(p+q)\mathrm{\Gamma }(p+q,q)D(q)}`$ (5.5)
The function $`\mathrm{\Gamma }(p,q)`$ corresponds to an off-diagonal inverse matrix element as it shows up for example in (4.22). Then application of (1.4) transforms (5.5) into (5.1). One may say that although the equations (5.4) are known, usually they are not really taken seriously. For our opinion this is due to two reasons. First of all these equations, being highly nonlinear, are not easy to solve. In particular, for models involving condensation phenomena like superconductivity or Bose-Einstein condensation, it seems to be apropriate to write them down in finite volume since some quantities may become macroscopic. And second, since they are usually derived from (5.5) by putting $`\mathrm{\Gamma }`$ equal to 1 (or actually -1, by the choice of signs in (5.5)), one may feel pretty suspicious about the validity of that approximation. The equations (5.1) tell us that this is a good approximation if the expansions for $`\sigma `$ and $`\pi `$ are asymptotic.
The applications of the method shown in this paper basically confirmed the common expectations for the particular models, thus one may say there are no really new results. However, we think it is fair to say that the computation of field theoretical correlation functions is an extremely difficult mathematical problem and therefore one should have welcome everything which sheds some new light on these problems. We hope that we could convince the reader that the method presented in this paper definitely does this.
References
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# Advection-Dominated Accretion with Infall and Outflows
## 1 Introduction
Advection-dominated accretion flows (ADAFs) have been invented to explain low-luminosity black hole candidates like Sgr A (Narayan et al., 1998) in our Galactic Center. The low-luminosity of this model is achieved in an optically thin plasma, where most of the energy is stored in hot ions, while electrons as potential radiators are inefficiently coupled to the heat source and remain relatively cold. The electrons become nonetheless mildly relativistic close to the central black hole and the inevitable synchrotron radiation is observed from most ADAF candidates. The presence of magnetic fields not far from energy equipartition with the gas is indicative of their origin in MHD-instabilities (Balbus & Hawley, 1991) leading to turbulence in the accretion flow and subsequent generation of an effective viscosity. On larger scales magnetic fields are likely to be responsible for the collimation of outflows from accretion disks into jets, seen in the cores of M87 (Reynolds et al., 1996) and NGC 4258 (Lasota et al., 1996; Herrnstein et al., 1998), which are prototypical ADAF candidates. Furthermore the model can explain the accretion in some low-luminosity AGNs of elliptical galaxies (di Matteo et al., 1999) and in NGC 4258 (Gammie, Narayan & Blandford, 1999). In a recent X-ray survey (Sambruna, Eracleous & Mushotzky, 1999) a few more examples have been found for low-luminosity core in radio-loud AGNs, which are candidates for advection-dominated accretion flows in their central engines, suggesting that ADAFs can be found even in radio-loud AGNs and that jet formation is a common feature.
Outflow models for ADAFs have been investigated by Blandford & Begelman (1999) and applied to several candidates (di Matteo et al., 1999; Quataert & Narayan,, 1999). On the theoretical side Igumenshchev & Abramowicz (1999) and Stone, Pringle & Begelman (1999) have performed time-dependent 2D calculations of accretion flows, which in some cases resemble ADAFs for certain viscosity parameters $`\alpha 0.1`$, but suggest the production of outflows for larger $`\alpha `$ (Igumenshchev & Abramowicz, 1999). It is found that the $`(rr)`$ stress tensor component, which was not included in the original description of vertically integrated models for accretions flows, is important in the cited calculations for ADAFs. The existence of self-similar solutions with a radial viscous force has been shown in previous work (Narayan & Yi, 1995a) and discussed for 2D solutions with a separation of variables.
In this paper we describe advection-dominated accretion flows in polar-integrated variables including the radial viscous braking force, which either produce outflows or are formed by wind infall. The wind infall is assumed to consist of free falling low angular momentum gas, which is cold with respect to the already existing accretion flow. One possible source for this gas are stellar winds from massive stars in young clusters as in the center of our galaxy. In the following we will refer to this wind infall into the accretion flow as infall. In our treatment the main difference between infall and outflow is the sign of the mass infall rate and we talk about winds, if it is not necessary to distinguish between infall and outflow. We restrict the discussion to an extension of the self-similar solutions given by Narayan & Yi (1994) for the Newtonian limit.
In §2 we present the equations, which describe the accretion flow including the reaction to winds. The role of $`(rr)`$ stresses and bulk viscosity is emphasised. We discuss angular momentum transport and viscosity in §3 and specify the possible equations of state and the resulting energy equation in §4. General features of self-similar solutions are presented in §5 and in §6 detailed solutions for the $`\alpha `$-viscosity law are shown. Consequences of the alternative $`\beta `$-viscosity are discussed in §7 and ADAFs with an intermediate shear-limited viscosity law follow in §8. We compare the solutions and draw our conclusions in §9.
## 2 Stationary accretion with infall or outflows
In a first step we have to establish the set of equations, which describe the accretion flow. Advection-dominated flows, in which we are interested, are known to be quasi-spherical (Narayan & Yi, 1995a) and therefore we will use spherical coordinates in our discussion. Consider a stationary, axisymmetric and rotating flow with angular velocity $`\mathrm{\Omega }`$ around a compact object of mass $`M`$. Instabilities in the flow, either hydrodynamic or magneto-hydrodynamic in origin, generate turbulence on small scales and lead to an effective viscosity much larger than the microscopic one. The effective turbulent viscosity $`\nu `$ in the flow will redistribute specific angular momentum $`\mathrm{}(r,\theta )=\mathrm{sin}^2\theta r^2\mathrm{\Omega }`$ by a local viscous torque between neighbouring rings or shells. Short term evolution on scales smaller than the mean free path of eddies $`\lambda `$ in the turbulent flow, which is related to the viscosity $`\nu =v_{\mathrm{eddy}}\lambda `$ is not resolved in this description. All quantities like density or accretion velocity must be understood in a local time averaged sense with probably large short term variations. Besides the time average, we will discuss accretion flows also as polar-averaged, one-dimensional flows with only a radial coordinate $`r`$. The average is taken over sections of shells occupied by the flow.
Polar motions can consequently not be traced in this treatment. We are left with the angular velocity $`\mathrm{\Omega }=\mathrm{\Omega }(r)`$ and the radial component of the velocity, which is the accretion velocity $`u=u(r)`$. Mass is accreted by the central object with a rate $`\dot{M}=2\pi ru\mathrm{\Sigma }`$. In this discussion $`\mathrm{\Sigma }`$ is a suitable polar integral of the density
$$\mathrm{\Sigma }=2_{\pi /2\vartheta }^{\pi /2+\vartheta }d\theta r\mathrm{sin}\theta \rho $$
(1)
so that $`H=r\vartheta `$ is the vertical thickness of the accretion flow and the integral in equation (1) is restricted to a region around the symmetry plane of the flow to make room for winds seen in Fig.1. Nonetheless any vertical motion has a radial and polar component and the later is considered as a source or sink of mass, momentum and energy. We regard these sources as external to our description of the flow and call it infall into the accretion flow or an outflow respectively.
For a non-relativistic flow the conservation of mass implies that the change in mass flux at every radius is balanced by mass exchange with the wind in the stationary case
$$\frac{1}{r}\frac{(ru\mathrm{\Sigma })}{r}=\dot{\mathrm{\Sigma }}_W,$$
(2)
where $`\dot{\mathrm{\Sigma }}_W`$ is the mass per surface area added to the flow from the infall per time at a given radius. According to equation (1) this happens at a height $`r\vartheta `$, but rapid mixing with the flow is assumed so that the equations averaged over polar section occupied by the accretion flow remain valid. Angular momentum is a conserved quantity in accretion flows without winds and is redistributed by torques generated by viscous stresses $`t_{r\varphi }`$
$$\frac{1}{r}\frac{}{r}(ru\mathrm{\Sigma }\mathrm{}r^2t_{r\varphi })=\dot{\mathrm{\Sigma }}_W\mathrm{}_W.$$
(3)
Here $`\mathrm{}`$ is a mass weighted polar average of $`\mathrm{}(r,\theta )`$. The infall contributes its own angular momentum to the flow $`\dot{\mathrm{\Sigma }}_W\mathrm{}_W`$ and exerts an additional torque, which depends on the relative angular momentum of infall and flow $`\mathrm{}_W\mathrm{}`$.
The accretion velocity $`u`$ follows from the imbalance of gravitational force, centrifugal barrier and radial pressure gradient. In addition to that we include the viscous braking force $`F_\nu `$ from the $`(rr)`$-component of the stress tensor and bulk viscosity
$`{\displaystyle \frac{1}{r}}{\displaystyle \frac{(ru\mathrm{\Sigma }u)}{r}}`$ $`=`$ $`\mathrm{\Sigma }\left(r\mathrm{\Omega }^2{\displaystyle \frac{GM}{r^2}}\right)r{\displaystyle \frac{𝒫}{r}}`$ (4)
$`+`$ $`F_\nu +\dot{\mathrm{\Sigma }}_Wu_W`$
with the integrated pressure $`𝒫=d\theta P`$ and the radial velocity of the wind $`u_W`$, which adds its momentum to the flow. The change in specific radial momentum induced by winds is proportional to the velocity difference $`u_Wu`$. The viscous force (see Appendix A)
$`F_\nu `$ $`=`$ $`F_{rr}+F_{\mathrm{bulk}}`$ (5)
$`F_{rr}`$ $`=`$ $`\eta _1[{\displaystyle \frac{4r}{3}}{\displaystyle \frac{}{r}}\left(\nu \mathrm{\Sigma }{\displaystyle \frac{}{r}}\left({\displaystyle \frac{u}{r}}\right)\right)`$ (6)
$`+`$ $`4\nu {\displaystyle \frac{}{r}}\left({\displaystyle \frac{u}{r}}\right)]`$
$`F_{\mathrm{bulk}}`$ $`=`$ $`\eta r{\displaystyle \frac{}{r}}\left({\displaystyle \frac{\nu \mathrm{\Sigma }}{r^3}}{\displaystyle \frac{(r^2u)}{r}}\right)`$ (7)
has a contribution proportional the $`(rr)`$ component of the shear tensor of the flow in equation (6) and the compression of the gas in equation (7). We make the crude assumption that the corresponding viscosities are the same for the $`(r\varphi )`$ and $`(rr)`$ components of the shear stress and also for the bulk viscosity. Our ignorance is cast into the parameters $`\eta _1`$ and $`\eta `$, which measure the strength of the $`(rr)`$-shear and bulk viscosity relative to the familiar $`(r\varphi )`$-shear viscosity respectively. For isotropic turbulence $`\eta _11`$ is expected, but no estimate on $`\eta `$ can be derived from this assumption. We will call $`\eta `$ and $`\eta _1`$ viscous force measures in the following. Finally an energy equation must be given
$`u\mathrm{\Sigma }{\displaystyle \frac{e}{r}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Sigma }}{r^2}}{\displaystyle \frac{(r^2u)}{r}}+Q^+\mathrm{\Lambda }`$ (8)
$`+`$ $`\dot{\mathrm{\Sigma }}_W\left(\omega _W\omega +{\displaystyle \frac{(𝐯_𝐖𝐯)^2}{2}}\right),`$
which describes the change of specific internal energy $`e`$ of the flow due to compression (first term right side), viscous heating $`Q^+`$, radiative cooling $`\mathrm{\Lambda }`$, the enthalpy difference between flow $`\omega `$ and wind $`\omega _W`$ and the kinetic energy associated with the velocity difference, which has to be dissipated into heat. Here $`𝐯`$ is the velocity vector of the flow and $`𝐯_𝐖`$ the corresponding vector for the wind.
## 3 Angular Momentum Transport and Viscosity
The $`(r\varphi )`$ component of the stress tensor is assumed to be proportional to the corresponding component of the shear tensor with the kinematic viscosity $`\mu =\nu \rho `$ as the factor of proportionality. For the viscosity $`\nu `$ we will adapt three different representations. The first and obvious one is the $`\alpha `$-prescription introduced by Shakura & Sunyaev (1973) and the resulting $`\alpha `$-disks will be discussed in §6. The $`\alpha `$-viscosity is not a unique choice and effects of the so-called $`\beta `$-viscosity introduced by Duschl, Strittmatter & Biermann are discussed as $`\beta `$-disks in §(7). For the $`\alpha `$-viscosity
$$\nu =\alpha \frac{c_s^2}{\mathrm{\Omega }_K}$$
(9)
we avoided to introduce the vertical scale-height $`H`$ of the disk, because we believe that it leads to the misunderstanding that angular momentum transport in the radial direction depends on the thickness of the disk. Only with the assumption of vertical hydrostatic equilibrium is it possible to introduce $`H=c_s/\mathrm{\Omega }_K`$ in (9).
The $`\alpha `$-viscosity can be recovered from the $`\beta `$-viscosity law in the case of shock-limited turbulence in a Keplerian disk. The $`\beta `$-viscosity assumes that a typical length scale in the direction of transport $`\mathrm{\Delta }r`$ and the typical velocity difference $`\mathrm{\Delta }v_\varphi `$ between interacting shells, which exchange eddies, equal the mean free path of eddies and their typical velocities relative to the mean flow respectively. A parametrisation with a new constant $`\beta `$ is suggested
$$\nu =\mathrm{\Delta }v_\varphi \mathrm{\Delta }r=\beta v_\varphi r.$$
(10)
If the typical velocity of eddies in a differentially rotating disk is limited by the sound speed, than the typical length scale $`\mathrm{\Delta }r`$ of communication between differentially rotating rings is estimated Duschl et al. to be
$$\mathrm{\Delta }r\frac{c_s}{\mathrm{\Omega }}.$$
(11)
Only in disks with Keplerian rotation does this length scale equal the vertical scale-height and the form $`\nu =\alpha c_sH`$ is recovered. As a third possibility we consider the shear-limited form of the $`\beta `$-viscosity with the length scale from equation (11), which has not been used for sub-keplerian accretion flows before. This is reasonable, if the mean free path of eddies is not limited by the vertical scale-height, but instead determined from the distance between shells, for which the velocity difference due to differential rotation equals the eddie velocity. This argument implies a larger effective viscosity for disks in sub-keplerian rotation. The $`(r\varphi )`$ component of the viscous stress tensor becomes in any case
$$t_{r\varphi }=\nu \mathrm{\Sigma }r\frac{\mathrm{\Omega }}{r}.$$
(12)
We assume that the same viscosity prescription can be applied to the $`(rr)`$-component of the stress and the bulk viscosity. Their relative strength is scaled to the $`(r\varphi )`$-viscosity by the force measures $`\eta _1`$ and $`\eta `$ introduced above. The heat generated by the described viscous forces has to be included in the energy equation and amounts to
$`Q^+`$ $`=`$ $`\nu \mathrm{\Sigma }\left(r{\displaystyle \frac{\mathrm{\Omega }}{r}}\right)^2+{\displaystyle \frac{4}{3}}\eta _1\nu \mathrm{\Sigma }\left(r{\displaystyle \frac{}{r}}\left({\displaystyle \frac{u}{r}}\right)\right)^2`$ (13)
$`+`$ $`\eta {\displaystyle \frac{\nu \mathrm{\Sigma }}{r^4}}\left({\displaystyle \frac{(r^2u)}{r}}\right)^2.`$
The internal energy is increased by viscous heat and compression and reduced by radiative cooling $`\mathrm{\Lambda }`$.
## 4 Equation of State
In general accretion flows are expected to exist for all reasonable equations of state for the accreted gas. The self-similar solutions we will use in what follows, are restricted to rather simple equations of state. For ADAFs the gas should be well described by an ideal gas law with $`\gamma =5/3`$ for ratio of specific heats, because the temperatures in ADAF solutions are so large ($`T>10^7`$K for most radii of interest) that the gas is completely ionised and only ions of a few heavy elements retain their highly bound inner electrons (see Narayan & Raymond (1999) for possible emission lines of these elements). The change of the equation of state due to partial ionization can be neglected, but the contribution of magnetic to the total pressure must be considered for two reasons. The first is the observational evidence that at least some candidates (e.g. Sgr A, see Narayan et al. (1998)) for ADAFs show a significant contribution of synchrotron emission to the total luminosity, which requires magnetic fields close to pressure equipartition with the gas. Given their existence, we might try a MHD-description of the accretion flow to separate the evolution of magnetic fields and gas, which might be different (Bisnovatyi-Kogan & Lovelace, 1997) unless magnetic diffusivity (Heyvaerts, Priest & Bardou, 1996) and/or reconnection require an even more complex model. The easy way out is to assume that turbulence generates small scale magnetic fields, which dominate the energy density in magnetic fields and produce an isotropic contribution to total pressure and energy. In doing so, one arrives at a hydrodynamic description of the accretion flow with a equation of state, which has to incorporate the magnetic pressure. Consider an equation of state of the quite general form
$$P\rho ^{\chi _\rho }T^{\chi _T}$$
(14)
with the exponents $`\chi _\rho ,\chi _T`$ defined according to Cox & Giuli (1968). For the self-similar solutions we have to restrict the internal energy per unit volume to be proportional to the total pressure
$$P=\frac{\chi _\rho }{\chi _T}(\gamma 1)e=(\mathrm{\Gamma }_31)e$$
(15)
with $`\gamma =c_P/c_V`$ being the ratio of specific heats and $`\mathrm{\Gamma }_3`$ the adiabatic coefficient defined by Chandrasekhar (1939). The restriction from equation (15) is that $`\chi _\rho ,\chi _T`$ and $`\mathrm{\Gamma }_3`$ have to be constants, which requires that the relative strength of the magnetic pressure must be constant so that the effective equation of state is an ideal gas law with a ratio of specific heats less than $`5/3`$. We define the isothermal sound speed $`c_s=\sqrt{r𝒫/\mathrm{\Sigma }}`$ from the integrated pressure and density. The sound speed used in the viscosity prescription is the total pressure divided by the mass density. We use thermodynamic relations (Cox & Giuli, 1968) to get an equation for the temperature from equation (8) for stationary accretion flows
$`uc_V{\displaystyle \frac{T}{r}}`$ $`=`$ $`{\displaystyle \frac{\chi _\rho rc_s^2}{\mathrm{\Sigma }}}u{\displaystyle \frac{}{r}}\left({\displaystyle \frac{\mathrm{\Sigma }}{r}}\right)+{\displaystyle \frac{Q^+\mathrm{\Lambda }}{\mathrm{\Sigma }}}`$ (16)
$`+`$ $`{\displaystyle \frac{\dot{\mathrm{\Sigma }}_W}{\mathrm{\Sigma }}}\left(\omega _W\omega +{\displaystyle \frac{(𝐯_𝐖𝐯)^2}{2}}\right).`$
Here $`c_V`$ is the specific heat at constant volume
$$c_V=\frac{\chi _T^2c_s^2}{\chi _\rho (\gamma 1)T}$$
(17)
and if we assume that gas pressure contributes a fraction $`\delta `$ to the total pressure and magnetic fields are responsible for the rest, we get
$$\mathrm{\Gamma }_31=\frac{\delta }{\frac{3}{2}\delta +(1\delta )}$$
(18)
for the adiabatic exponent $`\mathrm{\Gamma }_3`$ which is related to $`\gamma `$ by equation (15). For constant $`\chi _\rho `$ and $`\chi _T`$ we can rewrite the temperature gradient in equation (16) as a gradient of the sound speed and search for solutions in terms of the sound speed.
## 5 Self-similar ADAF solutions with a wind
The set of equations described in §2 and §3 allow self-similar power-law solutions for accretion velocity, angular velocity and sound speed in the way Narayan & Yi (1994) have shown
$$u=u_0s^{1/2}c\mathrm{\Omega }=\mathrm{\Omega }_0s^{3/2}\frac{c}{r_G}$$
(19)
$$c_s^2=a^2s^1c^2$$
(20)
with the radial coordinate scaled to the gravitational radius of the central mass
$$s=\frac{r}{r_G};r_G=\frac{GM}{c^2}$$
(21)
and the speed of light $`c`$. It is required that the cooling is either completely negligible or is a radius independent fraction of the heating rate
$$f=1\frac{\mathrm{\Lambda }}{Q^+}$$
(22)
so that $`1f`$ is the cooling efficiency, which will be small, if electrons are inefficiently coupled to the heat source. This happens, if ions are preferentially heated by viscous friction and the energy transfer rate between electrons and ions is smaller than the heating rate.
For the self-similar solution the radial component of the wind velocity must be proportional to the accretion velocity and we get a measure $`\xi _1=u_W/u`$, which tells us the relative radial velocity of the wind, but we will neglect it in the following discussion and assume $`\xi _1=0`$. This is motivated from the infall calculations of Coker, Melia & Falcke (1999), which predicts a very small radial velocity of the infalling material in the disk mid-plane for the infall from stellar winds in the Galactic Center. If a fraction $`f`$ larger than a few per cent of the viscously generated heat is not radiated away but kept in the flow as internal energy, the disk scale height $`H=c_s/\mathrm{\Omega }_K`$ is of the same order as the radius and the infall may have a significant radial velocity component. Nonetheless we will neglect it in the following discussion. The same is true for the enthalpy of the infalling material, but here we have better reasons to believe that the gas joining the ADAF is cold compared with the gas in the ADAF, which is heated by viscous friction, while the infalling gas maybe provided by stellar winds or the gas of the Hii region Sgr A West in case of the Galactic Center with temperatures of $`\mathrm{6\hspace{0.17em}10}^310^4`$ K (see Mezger, Duschl & Zylka (1996) for a review of the Galactic Center). This should be compared to typical temperatures of ADAFs of $`10^7`$ K at $`10^4`$ black hole radii. The situation with outflows might be different, where it is hard to imagine that the gas leaving the disk has a different temperature than the gas left in the flow. The gas might need some outward pointing radial momentum to get away from the ADAF, but again we will ignore it here. For infalling material we expect it to be in free fall so that their total velocity is $`v_W=v_{\mathrm{ff}}=\sqrt{2}cs^{1/2}`$. In the case of outflows we require that they are able to escape the gravitational attraction of the central object and their velocity must be at least the free fall velocity $`v_{\mathrm{ff}}`$.
In the same way the rotation of the infall has to be a constant fraction $`\xi =\mathrm{\Omega }_W/\mathrm{\Omega }`$ of the rotation of the accretion flow itself and again calculations by Coker et al. (1999) for the Galactic Center predict infall with only weak rotation. Outflows on the one hand might be expected to carry their initial angular momentum from the point of origin and therefore $`\xi =1`$. On the other hand magnetically driven outflows will exert a torque on the remaining accretion flow if the gas in the outflow has different angular velocity than the point in the flow, to which it is connected by magnetic field lines. In this case $`0<\xi <1`$ would be expected.
The rate of mass added by the infall per surface area of the flow in its central plane is constrained to steep radial profiles and introduces one free parameter $`p`$ (Blandford & Begelman, 1999) not to be confused with the pressure $`P`$
$$\dot{\mathrm{\Sigma }}_W=p\frac{u\mathrm{\Sigma }}{r},\mathrm{\Sigma }=\mathrm{\Sigma }_0s^{1/2+p}.$$
(23)
For infall into accretion flows $`p`$ has to be negative and the radial dependence of the infall rate must be steeper than $`r^2`$. Therefore most of the mass is added at small radii and the total mass infall $`d\dot{M}_Wdrr^{1+p}`$ diverges in the center. This solution is consequently not valid for small radii, where the infall must deviate from the solution (23). Positive values of $`p`$ correspond to outflows generated by the accretion flow and reversing the argument for the total mass now taken away from the ADAF, the solution is confined to small radii so that $`\dot{M}_W`$ is still smaller than the mass accretion rate supplied at the outer radius. The surface density of the flow implies a density $`\rho r^{3/2+p}`$, which restricts the specific equation of state (14) to those with $`\chi _\rho `$,$`\chi _T`$ and $`\mathrm{\Gamma }_3`$ being constants and the temperature as a function of radius follows
$$T^{\chi _T}r^{1+(1\chi _\rho )(3/2+p)}$$
(24)
so that the radial pressure gradient is independent of the thermodynamic exponents and $`Pr^{5/2+p}`$. The gas density in self-similar solutions for advection-dominated accretion is unconstrained and only the required inefficient cooling restricts the solutions to small mass accretion rates below $`10^2\dot{M}_{\mathrm{Edd}}`$ (Narayan & Yi, 1995b). $`\dot{M}_{\mathrm{Edd}}`$ is the Eddington accretion rate for an radiation efficiency of 10% in terms of the rest mass of the accreted material.
## 6 $`\alpha `$-ADAFs
Whether the parametrisation of the wind, which is necessary for the self-similar solution in §5, is justified or not, depends on boundary conditions, which the wind has to meet at large distances, or the internal physics of outflows. While outflows might naturally follow the self-similarity of ADAFs, the distribution of stellar wind sources will determine, if a power-law dependence of the infall rate is reasonable. The parametrisation of the viscosity (Shakura & Sunyaev, 1973) in equation (9) is based on dimensional arguments and besides numerical simulations, which do not exist for ADAF conditions, there is no way to determine what values for $`\alpha `$ are appropriate in the case of stationary ADAFs. Without information on the strength of the magnetic fields from MHD-calculations, the adiabatic exponent $`\mathrm{\Gamma }_3`$ is also undetermined. We will discuss the self-similar solutions as functions of $`\alpha `$ and $`\mathrm{\Gamma }_3`$ as well as for certain wind parameters and viscous force measures $`\eta `$ and $`\eta _1`$.
The self-similar ansatz in §5 reduces the dynamical equations (3), (4) and (16) to a set of non-linear algebraic equations for the coefficients of sound speed $`a`$, angular velocity $`\mathrm{\Omega }_0`$ and accretion velocity $`u_0`$. The accretion velocity is derived from equation (3) in terms of sound speed $`a`$ and $`\alpha `$
$$u_0=\frac{3}{2}\alpha a^2𝒜𝒜=\frac{1+2p}{1+2p(1\xi )}$$
(25)
with a constant $`𝒜`$, which reflects the torque exerted by the the wind on the flow. The combined wind parameter $`𝒜`$ is 1 for a vanishing wind or a non-rotating wind. For rotating infall the solution is restricted to moderate infall rates with $`p>1/2`$. Otherwise accretion would not be possible. For a given sound speed and accretion velocity, the centrifugal barrier and therefore the angular velocity follows from the radial momentum equation (4), if we know the viscous force measures. We find for the square of the angular velocity <sup>1</sup><sup>1</sup>1
$$=𝒜(12p(1\xi _1))+\eta (52p)+\frac{4}{3}\eta _1(1+2p)$$
$$\mathrm{\Omega }_0^2=c^2\left(\frac{5}{2}p\right)a^2\frac{9}{8}\alpha ^2𝒜\frac{a^4}{c^2}.$$
(26)
We are left with a quadratic equation for the square of the sound speed. One solution turns out to be irrelevant either by predicting $`a^2`$ to be negative or leaving $`\mathrm{\Omega }_0^2`$ in equation(26) negative. But even the second root not always gives reasonable solutions, because with increasing $`\alpha `$ and small changes in $`a^2`$ as seen in Fig.2 and Fig.3, $`\mathrm{\Omega }_0^2`$ decreases and becomes negative. Solutions for the accretion problem exist only for values of the viscosity parameter smaller than a critical $`\alpha _c`$. The wind parameters for the solutions in Fig. 2 and 3 have been chosen in a way that the infall or outflow is in free fall or leaves with the escape speed. No radial momentum of the wind is included and both infall and outflow are cold. While the infall has no proper angular momentum, the outflow rotates with the $`\mathrm{\Omega }`$ of the accretion flow. With the requirements of a cold outflow and total velocity of the wind being the escape speed, which both enter the energy equation (16) in the same way, we make a minimal energy assumption for the extraction of internal energy by the outflow. For the equation of state we take the natural choice $`\chi _\rho =1`$ so that the mix of gas and magnetic field behaves like an ideal gas. The ratio of specific heats equals $`\mathrm{\Gamma }_3`$, if we assume $`\chi _T=1`$. $`\chi _T`$ is not a parameter of our solutions, but determines the actual temperature of the gas through equation (24). For energy equipartition we have $`\mathrm{\Gamma }_3=1.4`$ from equation (18) and we use this value if not stated otherwise. For the cooling efficiency $`1f`$ we assume a low rate of 1%. For the dynamics of the flow not the energy radiated away matters, but the energy left in the flow. The actual cooling efficiency is unimportant as long as it is small and $`f1`$.
For an ADAF without winds it is possible to derive an analytic solution for the sound speed and determine the maximum allowed $`\alpha _c`$ given by<sup>2</sup><sup>2</sup>2 $``$ $`=`$ $`2(5\eta +{\displaystyle \frac{4}{3}}\eta _1)[{\displaystyle \frac{4}{9}}(\mathrm{\Gamma }_31)`$ (27) $`\times `$ $`(({\displaystyle \frac{7}{3}}\mathrm{\Gamma }_3)f({\displaystyle \frac{5}{3}}\mathrm{\Gamma }_3))]`$ $`+`$ $`{\displaystyle \frac{32}{3}}\eta _1f(\mathrm{\Gamma }1)\left({\displaystyle \frac{5}{3}}\mathrm{\Gamma }_3\right)`$
$$\alpha _c=\frac{\sqrt{8/92(\mathrm{\Gamma }_31)(7/3\mathrm{\Gamma }_3)+}}{f(\mathrm{\Gamma }_31)(3\eta +4\eta _1)}.$$
(28)
In that case $`𝒜`$ equals $`1`$ and the accretion velocity is $`u_0=(3/2)\alpha a^2`$. If the viscosity parameter is small, the limit $`\alpha 0`$ provides a good approximation for the sound speed
$$a^2\frac{6f(\mathrm{\Gamma }_31)}{3(\mathrm{\Gamma }_31)(5f2\chi _\rho )+106\chi _\rho }$$
(29)
and the angular velocity $`\mathrm{\Omega }=\sqrt{c^2(5/2)a^2}`$ as seen in Fig.2 and Fig.3. From equation (28) and (27) we immediately see that $`1<\mathrm{\Gamma }_3<5/3`$ is required for the no wind case and the familiar result arises that no solution exists for a non-relativistic ideal gas. It is also obvious that the upper limit on $`\alpha `$ increases and becomes irrelevant with vanishing radial viscous force. In other words, with increasing viscous force measures $`\eta ,\eta _1`$ ADAF solutions are restricted to reasonably small values of $`\alpha <\alpha _c`$.
One additional criterion for the realization of accretion solutions is the Bernoulli number
$$\mathrm{Be}=\frac{GM}{r}+r^2\mathrm{\Omega }^2+u^2+e+\frac{P}{\rho }$$
(30)
discussed by Narayan & Yi (1994, 1995a) for ADAFs and Blandford & Begelman (1999) for ADIOS. While the total energy decreases inwards the Bernoulli number being the total specific energy plus $`P/\rho `$ is positive and increases inwards for ADAFs. This changes with outflows due to their cooling effect by removing internal energy and reducing the Bernoulli number of the remaining flow. Correspondingly an infall increases the Bernoulli number of the flow. The combinations of allowed values for $`\alpha `$ and $`\mathrm{\Gamma }_3`$ as a function of the wind strength $`p`$ is shown in Fig.4 for corotating ($`\xi =1`$) and non-rotating winds. The statement that ADAFs exist for an ideal gas with negative Bernoulli numbers, if an outflow is present (Blandford & Begelman, 1999), is confirmed, even when the radial viscous force is included. The minimal possible wind strength depends on the viscosity parameter $`\alpha `$ as seen in Fig.4.
## 7 $`\beta `$-ADAFs with winds
Following the arguments in Duschl et al. (2000) the viscosity law should not depend on the sound speed, if the actual velocity of turbulent eddies is smaller than the sound speed in the gas. In the $`\beta `$-viscosity law (10) the eddy velocity is determined from the actual rotation velocity $`v_{\mathrm{eddy}}\sqrt{\beta }r\mathrm{\Omega }`$ with the lessening $`\sqrt{\beta }`$ equally distributed between velocity and length scale. The arguments given in §3 suggest that $`\alpha `$-viscosity can be seen as the shock-limited case of the $`\beta `$-viscosity and the $`\alpha `$-ADAFs discussed in §6 can be checked for consistency, if this hypothesis is correct. Specialising on the no wind case in the limit of small $`\alpha `$ we find the ratio
$$\frac{c_s^2}{(r\mathrm{\Omega })^2}=\frac{f(\mathrm{\Gamma }_31)}{(5/3\mathrm{\Gamma }_3\chi _\rho )}$$
(31)
and the $`\beta `$-viscosity can be applied for $`\beta <1.5f`$ with the adaption of $`\chi _\rho =1`$ and $`\mathrm{\Gamma }_3=1.4`$. Therefore the $`\beta `$-viscosity is valid for ADAFs with low cooling efficiencies, if the resulting flows are similar to $`\alpha `$-disks. The scale length of the $`\alpha `$-viscosity is the vertical scale height $`H=c_s/\mathrm{\Omega }_K`$, which turns out be $`H1.55r\sqrt{f/(1.6+6f)}0.6r`$ for $`\alpha `$-ADAFs in the no wind case and the limit $`\alpha 0`$. The length scale of the $`\beta `$-viscosity is $`\sqrt{\beta }r`$ and we expect the same sound speeds and angular velocities for $`\alpha `$\- and $`\beta `$-ADAFs below $`\beta 0.3`$. This is verified by comparing Fig.2 and 3 with Fig.5. The actual solution for $`\beta `$-ADAFs is derived in the same way as for $`\alpha `$-ADAFs in §6 and we find the accretion velocity depending explicitly on $`\mathrm{\Omega }`$ and not $`c_s^2`$.
$$u_0=\frac{3}{2}\beta \mathrm{\Omega }_0𝒜$$
(32)
The constant $`𝒜`$ is defined in equation (25). The increased accretion velocity seen in Fig.6 compared to $`\alpha `$-ADAFs is explained by the larger value of $`\mathrm{\Omega }_0`$ in equation (25) relative to the sound speed squared in case of the $`\alpha `$-viscosity. The angular velocity of the accretion flow can be given<sup>3</sup><sup>3</sup>3
$$=\left((52p)\eta +4/3(1+2p)\eta _1\right)𝒜$$
$$\mathrm{\Omega }_0^2=\frac{1(5/2p)a^2}{9/8\left[(12p(1\xi _1))𝒜^2+\right]\beta ^2+1}$$
(33)
including the reaction of the flow due to the presence of winds. Turning to the energy equation shows that two obvious solutions exist for the sound speed. First an non-rotating and non-accreting solution
$$a^2=\frac{2c^2}{52p},$$
(34)
where the gravitational attraction is completely balanced by the pressure gradient. The gas is in hydrostatic equilibrium and constitutes an optically thin atmosphere.
This is not a viable solution for a black hole, because general relativity does not allow a pressure supported atmosphere close to the horizon. The second solution is an accreting and rotating flow as shown in Fig.5. The gas is cooler than in the atmosphere solution and strong outflows in corotation with the flow or even non-rotating outflows reduce the sound speed further so that the accreted gas can be gravitationally bound to the central mass. For a corotating, cold outflow, which reaches infinity with the escape velocity from the radius of origin $`\xi _2=2`$, and a simple equation of state $`\chi _\rho =1`$ we can give the sound speed for a case of negligible radial viscosity $`\eta =\eta _1=0`$: <sup>4</sup><sup>4</sup>4 $`\mathrm{\Xi }`$ $`=`$ $`{\displaystyle \frac{𝒜}{24}}[4(6p)(12p)(1+\beta ^2𝒜^2)`$ $`+`$ $`\beta ^2𝒜^2(3(1+4p)^2p(17+2p))]`$ $`+`$ $`{\displaystyle \frac{𝒜(1p)}{\mathrm{\Gamma }_31}}\left({\displaystyle \frac{2}{3}}+{\displaystyle \frac{3}{4}}\beta ^2𝒜^2(12p)\right)`$
$$a^2=\frac{f\left[3(3/4p)\beta ^2𝒜^2+1\right]p𝒜}{(5/2p)f\mathrm{\Xi }}.$$
(35)
For non-rotating and cold infall, which is in free-fall for $`p<0`$, or outflows barely reaching infinity ($`\xi _2=1`$) with $`p<0.36`$ the ratio of specific heats has to be smaller than a limiting value
$$\mathrm{\Gamma }_31<2\frac{1p}{39p+2p^2}$$
(36)
at which the rotation of the flow stops and no solutions exist for $`\mathrm{\Gamma }_3`$ larger than that. This limit for $`\mathrm{\Gamma }_3`$ equals $`5/3`$ for ADAFs without winds or outflows and is the well known result that ADAFs do not exits for an ideal equation of state. The consistency of $`\beta `$-ADAFs follow from the above predictions based on $`\alpha `$-flows. With infall into the flow the sound speed increases and the rotation velocity decreases. The $`\beta `$-viscosity law is therefore more likely to be applicable to ADAFs with infall than with outflows.
## 8 ADAFs with shear-limited viscosity and winds
The motivation for a shear limited viscosity law is the observation that the standard $`\alpha `$-viscosity agrees with the general $`\beta `$-viscosity law only for Keplerian accretion disks with a shock-limited eddy velocity. In the shear-limited viscosity
$$\nu =\widehat{\alpha }\frac{c_s^2}{\mathrm{\Omega }}$$
(37)
we assume that the eddy velocity is indeed shock-limited and the mean free path of eddies is determined by the radial distance of differentially rotating shells, which allows the exchange of eddies along radial paths in a comoving frame of one shell. If the eddy velocity is the sound speed, the maximal distance of shells follows from a Taylor expansion of $`\mathrm{\Omega }`$ in $`\mathrm{\Delta }v_\varphi `$
$$\mathrm{\Delta }r\frac{c_s}{2\mathrm{\Omega }}.$$
(38)
The length scale over which shells can interact is larger for sub-keplerian rotation. This is the main difference of equation (37) to the original $`\alpha `$-viscosity. The consequence for the accretion velocity
$$u_0=\frac{3}{2}\widehat{\alpha }\frac{a^2}{\mathrm{\Omega }_0}𝒜$$
(39)
is an explicit dependence on the inverse of $`\mathrm{\Omega }_0`$ contrary to $`\beta `$-ADAFs in equation (32). The combined wind parameter $`𝒜`$ is defined in equation (25). The square of the angular velocity follows from a quadratic equation with two roots <sup>5</sup><sup>5</sup>5 $`\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \frac{9}{8}}\widehat{\alpha }^2𝒜[(12p(1\xi _1))𝒜`$ $`+`$ $`{\displaystyle \frac{4}{3}}\eta _1(1+2p)+\eta (52p)]`$
$$\mathrm{\Omega }_0^2=\frac{1}{2}\frac{52p}{4}a^2\pm \sqrt{\left(\frac{1}{2}\frac{52p}{4}a^2\right)^2a^4\mathrm{\Psi }}.$$
(40)
In the limit of vanishing viscosity parameter $`\widehat{\alpha }`$ the angular velocity has two obvious solutions
$$\widehat{\alpha }0\mathrm{\Omega }_0\{\begin{array}{c}\sqrt{1(5/2p)a^2}\hfill \\ \mathrm{or}\hfill \\ 0\hfill \end{array}.$$
(41)
For small $`\widehat{\alpha }`$ the rotation for the second solution is proportional to $`\widehat{\alpha }`$ and positive. The existence of these slowly or non-rotating ADAFs depends on non-vanishing force measures $`\eta ,\eta _10`$. The solution disappears without a radial viscous force. The viscosity in non-rotating ADAFs tends to a finite value at $`\widehat{\alpha }=0`$
$$\nu =a\sqrt{\frac{1(5/2p)a^2}{a^2\mathrm{\Psi }}}c\sqrt{sr_G}$$
(42)
and so does the accretion velocity. If we interpret $`\widehat{\alpha }c_s/\mathrm{\Omega }`$ as the viscous length scale, it does not diverge for small $`\widehat{\alpha }`$. The sound speed also tends to a positive value as seen in the upper solution branches for the sound speed in Fig.8 and accretion velocity in Fig.7 corresponding to the lower branches for the angular velocity in Fig.7. The energy equation is itself quadratic in $`a^2`$ for either choices of $`\mathrm{\Omega }_0`$ in equation (40). Only one solution satisfies the condition
$$a^2\frac{2}{52p}$$
(43)
for real values of $`\mathrm{\Omega }`$. We end up with two solution branches, which join at a critical $`\widehat{\alpha }_c`$.
Beyond $`\widehat{\alpha }_c`$ no physical solution exists for the accretion problem. Besides the slowly rotating solution discussed above, we find more familiar ADAFs, where the accretion velocity is linear in $`\widehat{\alpha }`$ as seen in Fig.7. For small viscosity parameters, the solution is independent of the force measures $`\eta ,\eta _1`$. Real solutions for the square of the sound speed in the no wind case are only possible for
$$\widehat{\alpha }\frac{(5/3\chi _\rho \mathrm{\Gamma }_3)}{3f(\mathrm{\Gamma }_31)\sqrt{\eta +(4/3)\eta _1}}.$$
(44)
The situation does not differ from the standard $`\alpha `$-viscosity in so far as no upper limit for $`\widehat{\alpha }`$ exists in the limit of vanishing radial viscosity $`\eta =\eta _1=0`$ in the radial momentum equation (4) as long as the ratio of specific heats is in the range $`1<\mathrm{\Gamma }_35/3`$. In the limit of vanishing viscosity $`\widehat{\alpha }0`$ the rotating solution coincides with the unique ADAF solution without winds and radial viscosity. For the simplest equation of state $`\chi _\rho =1`$ the sound speed in the limit $`\widehat{\alpha }0`$ is
$$a^2=\frac{2ϵf}{2ϵ^2+5ϵf},ϵ=\frac{5/3\mathrm{\Gamma }_3}{\mathrm{\Gamma }_31}.$$
(45)
The critical $`\widehat{\alpha }_c`$, which sets the upper limit for the existence of ADAFs with shear-limited viscosity, depends not only on the force measures $`\eta ,\eta _1`$, but also on the presence of winds. The same pattern as for the standard $`\alpha `$-viscosity appears, in which solutions with outflows are possible for a wider range of $`\widehat{\alpha }`$ than for infall solutions as seen in Fig.7 and Fig.8. It turns out that the slowly rotating branch of solutions produces smaller Bernoulli numbers than the fast rotators for infall and the no wind case $`p=0`$. Only for outflow solutions is the fast rotator preferred with smaller Bernoulli numbers. But in all cases shown in Fig.8, the Bernoulli number is positive and no gravitationally bound flow is found. This changes if more energy is extracted by the wind and the fast rotating outflow solution shown in the figures has negative Bernoulli numbers, if the outflow is cold and reaches a terminal velocity equal to the escape speed at its origin. The allowed combinations of $`\widehat{\alpha }`$ and $`\mathrm{\Gamma }_3`$ regardless of Bernoulli numbers are shown in Fig.9. The dividing lines in the $`(\widehat{\alpha },\mathrm{\Gamma }_3)`$-plane are vertical without a radial viscous force and bend in the way drawn in the figure, when the $`(rr)`$ component of the stress tensor is included.
## 9 Conclusions
We have investigated the importance of a radial viscous braking force for advection-dominated accretion flows (ADAFs) in the presence of infall or outflows. Under the assumption of almost isotropic turbulence, bulk viscosity and $`(rr)`$ component of the viscous stress tensor provide an efficient brake of the rapid radial accretion of hot gas in ADAFs and produce additional heat due to viscous friction. This opens a second channel for transfer of kinetic into internal energy and supports in part the radial motion so that the shear due to differential rotation and the centrifugal barrier is reduced in ADAFs with radial viscous braking.
We derived self-similar solutions of ADAFs for three different viscosity laws. The standard $`\alpha `$-viscosity produces solutions, which show more and more sub-keplerian rotation with increasing $`\alpha `$. The solutions terminate at a critical $`\alpha _c`$, which depends strongly on the ratio of specific heats so that the greatest possible $`\alpha `$ tends to zero for a non-relativistic ideal gas. At the critical $`\alpha _c`$ the rotation of the flow vanishes and a purely radial inflow appears, which is not only supported by a pressure gradient, but also by the viscous braking force, provided the flow is still turbulent in the absence of differential rotation. This limit is reminiscent of Bondi accretion (Bondi, 1952) in the presence of effective turbulent viscosity. The transition from ADAFs to Bondi accretion nonetheless affords a suspicious fine tuning in $`\alpha `$.
The $`\beta `$-viscosity law of Duschl et al. (2000) based on geometrical arguments in the absence of shock-limited turbulence allow ADAF solutions for all reasonable values of $`\beta `$. The estimates of $`\beta `$ as inverse of the critical Reynolds number of the flow suggests small $`\beta `$s, for which the solutions do not differ from $`\alpha `$-ADAFs. No transition to a Bondi like flow is possible in this case and no upper bound on $`\beta `$ exists.
We showed that it is possible to derive a shear-limited viscosity law from the $`\beta `$-viscosity mentioned above, for which the transition to non-rotating Bondi like accretion with turbulent viscosity occurs naturally. The now familiar ADAF solutions with sub-keplerian rotation at $`\widehat{\alpha }0`$ join with a second branch of solutions at a maximal sustainable $`\widehat{\alpha }`$. Similar to the $`\alpha `$-ADAFs no accretion flow with larger $`\widehat{\alpha }`$ are possible. The second solution branch is a hot, slowly rotating, and rapidly accreting solution even for small values of $`\widehat{\alpha }`$. The transition to viscous Bondi accretion occurs from the slowly rotating to the non-rotating solution in the limit $`\widehat{\alpha }0`$, where $`\mathrm{\Omega }`$ is linear in $`\widehat{\alpha }`$. A finite viscous force is present in these flows even in the limit $`\widehat{\alpha }=0`$. In all cases the connection of ADAFs to non-rotating flows depends on the presence of a radial viscous force.
The observation that ADAFs generally possess positive Bernoulli numbers lead to the idea (Narayan & Yi, 1995a; Blandford & Begelman, 1999) that ADAFs are good candidates for the production of outflows. In that way the accretion flow loses energy to the outflow and the remaining material is left with negative Bernoulli numbers and gravitationally bound to the central accreting mass. We confirm that statement in the presence of a radial viscous force for $`\alpha `$\- and $`\beta `$-ADAFs and show the back-reaction of outflows on ADAFs for different outflow characteristics. Most noticeably is the cooling effect and the increased accretion velocity of the remaining ADAF as angular momentum and internal energy is carried away. Outflows with a minimal energy assumption for the extracted energy have to be fairly massive $`p0.35`$ to lower the Bernoulli number and leave a bound flow in case of the shear-limited viscosity. It is much easier to get a bound flow, if the minimal energy assumption is violated and the terminal velocity of the outflow is the escape speed from the origin of the out-flowing material.
If a natural choice of the viscosity parameter—either $`\alpha `$ or $`\widehat{\alpha }`$—exists and provided the strength of the radial viscous force can be estimated from isotropic turbulence, then outflows are inevitable for certain equations of state with ratios for specific heats close to $`5/3`$. This conclusion is independent of arguments based on the positiveness of Bernoulli numbers for ADAFs.
The scenario of thin disk evaporation in binary systems (Liu et al., 1999) or the transition from cooling flows to ADAFs in low-luminosity cores of elliptical galaxies (Quataert & Narayan,, 1999b) can explain the existence of ADAFs in these systems. The formation of an ADAF, which is the most promising model for the spectral energy distribution of Sgr A in the Galactic Center, cannot proceed in either way. We suggest that an ADAF in the Galactic Center forms out of stellar wind infall (Coker et al., 1999) and the transfer of kinetic energy of the infall into internal energy of an advection-dominated flow. The even larger Bernoulli numbers produced in this way wound give rise to subsequent outflows and reduce the mass accretion rate inferred from infall calculations by Coker et al. (1999) to the smaller accretion rates predicted from spectral fitting (Quataert & Narayan,, 1999) of ADAF models to Sgr A.
I thank Ramesh Narayan and Wolfgang Duschl for helpful conversations and challenging discussions. This work has been supported through DAAD fellowship D/98/27005.
## Appendix A The radial viscous force
The radial viscous force in spherical coordinates is derived from the tensor divergence of the shear stress tensor $`t_{ij}`$ and the bulk viscous force from the divergence of bulk viscosity and compression. The stress tensor is assumed to be proportional to the shear tensor $`𝒟`$ which is defined as
$$𝒟_{ij}=\frac{1}{2}(v_{i;j}+v_{j;i})\frac{1}{3}g_{ij}v_{;l}^l$$
(A1)
with the kinematic viscosity $`\mu =\nu \rho `$ as a scalar function relating shear and stress $`t_{ij}=\mu 𝒟_{ij}`$. Here $`v_i`$ are velocity vector components, $`g_{ij}`$ the metric tensor and $`v_{i;j}`$ implies covariant differentiation of the velocity. The radial viscous force derived from the shear tensor is
$`f_\nu `$ $`=`$ $`(2\mu 𝒟^{rj})_{;j}`$
$`=`$ $`{\displaystyle \frac{4}{3}}\left({\displaystyle \frac{}{r}}\left[\mu r{\displaystyle \frac{}{r}}\left({\displaystyle \frac{u}{r}}\right)\right]+3\mu {\displaystyle \frac{}{r}}\left({\displaystyle \frac{u}{r}}\right)\right),`$
where axial-symmetry and $`v_\theta =0`$ is assumed and $`u`$ is the radial velocity as used throughout the paper. The corresponding heating rate due to this viscous friction is
$$q_{rr}^+=\frac{4}{3}\mu r^2\left(\frac{}{r}\left(\frac{u}{r}\right)\right)^2.$$
(A3)
The contribution from compression and bulk viscosity $`\zeta `$ to the radial viscous force is
$$f_{\mathrm{bulk}}=(\zeta v_{;l}^l)_{;r}=\frac{}{r}\left(\frac{\zeta }{r^2}\frac{}{r}(ur^2)\right)$$
(A4)
and the corresponding heating rate is
$$q_{\mathrm{bulk}}^+=\left(\frac{\zeta }{r^2}\frac{}{r}(ur^2)\right)^2.$$
(A5)
Multiplying by $`r`$ and performing polar integration of the force gives the radial viscous force used in (6) and (7)
$$F_{rr}=\frac{4r}{3}\left(\frac{}{r}\left[\nu \mathrm{\Sigma }\frac{}{r}\left(\frac{u}{r}\right)\right]+3\nu \mathrm{\Sigma }\frac{}{r}\left(\frac{u}{r}\right)\right)$$
(A6)
and
$$F_{\mathrm{bulk}}=r\frac{}{r}\left(\frac{\nu \mathrm{\Sigma }}{r^3}\frac{}{r}(ur^2)\right)$$
(A7)
with the replacement of the bulk viscosity by the effective turbulent viscosity $`\nu `$. The heating rates used in equation (13) are derived from equation (A3) and (A5) in the same way. The force measures $`\eta ,\eta _1`$ used in the text have been omitted here.
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# 1 Introduction
## 1 Introduction
One of the issues of string theory preventing phenomenological applications is the moduli that we get upon compactification. Since the critical dimension is $`D=10`$ for (worldsheet) $`N=1`$ superstrings, one has to compactify six of ten dimensions to get four-dimensional spacetime. Then there appear a number of massless scalars which characterize the size, the shape and other structures of the compactification space, although no so many (approximately) massless scalars are expected to be observed in reality. One must also break supersymmetries. It would therefore be interesting to explore any possibility to find a lower-dimensional string model with less supersymmetries which we could use as a starting point.
Recently, an interesting duality between strings on singular Calabi-Yau spaces and lower-dimensional non-gravitational (string) theories has been discussed -. In this paper, we will construct a modular invariant string partition function on four-dimensional Minkowski $`\times `$ two-dimensional black hole . This model may be thought of as a noncompact $`SL(2,\text{})/U(1)`$ version of Gepner models and corresponds to a conifold point on the moduli space of Calabi-Yau compactifications of type II string theories. We directly deal with the characters of $`N=2`$, $`c=9`$ unitary superconformal field theory realized as a Kazama-Suzuki model based on $`SL(2,\text{})/U(1)`$ at level $`k=3`$. Unlike the conventional Gepner models, which use a tensor product of minimal models (realized as a compact coset $`SU(2)/U(1)`$ ) as the internal CFT and is known to describe regular Calabi-Yau compactifications, the necessary central charge 9 is supplied by a single $`SL(2,\text{})/U(1)`$ conformal field theory. Thus we do not need to take a tensor product. If we regard the coset conformal field theory as a gauged $`SL(2,\text{})`$ WZW model , we end up with a six dimensional string theory on four-dimensional Minkowski space $`\times `$ two-dimensional black hole<sup>1</sup><sup>1</sup>1 This is a critical string theory since the total central charge is 0 and the Liouville mode decouples. However, in the linear dilaton region which is far from the tip of the cigar, the coordinate field along the cigar looks like the Liouville field which adjusts the total central to 0 if treated as a conformal field . The idea of replacing the Liouville $`\times `$ $`U(1)`$ system in noncritical string theories with two-dimensional Euclidean black hole was proposed in . .
The main obstacle in constructing modular invariants using $`N=2`$, $`c>3`$ superconformal characters - was its bad modular behavior (See for an earlier attempt.); for example, the nondegenerate NS character is given by
$`\text{Tr}q^{L_0}=q^{h+1/8}{\displaystyle \frac{\vartheta _3(0|\tau )}{\eta ^3(\tau )}},`$ (1)
whose monomial factor $`q^{h+1/8}`$ makes the modular behavior awful. To overcome this problem, we use the following two ideas : First, we consider (a countable set of) infinitely many primary fields so that the sum of their $`q^{h+1/8}`$ factors form a certain theta function. Which theta function to choose has to be examined carefully, and will be determined later. The modular property of the total partition function can be thus improved. Consistent CFTs with an infinite number of primary fields are known (e.g. $`c=1`$ CFTs ) and not surprising. On the contrary, they are required for modular invariance .
There is still a problem even after the monomial of $`q`$ is replaced by a theta function; the number of eta and theta functions do not balance between the denominator and the numerator. It is a problem because in the modular $`S`$ transformation the $`\sqrt{\tau }`$ factors do not cancel. The situation is similar for a free scalar partition function, in which, however, the $`|\tau |`$ factor from the left and right eta functions is compensated by the modular transformation of the zeromode integral. Thus our second proposal is to assume that the “internal” $`N=2`$ superconformal system constructed from the $`SL(2,\text{})/U(1)`$ coset has a degree of freedom of the center-of-mass motion along a certain direction, which has to be integrated over in the partition function. This assumption just agrees with the picture of CFTs on singular Calabi-Yau spaces advocated by Witten . With these two ideas we construct a modular invariant, and it turns out that the integration over the continuous set of ensembles above is nothing but the “Liouville-momentum” integration along the cigar (See for an earlier analysis of strings on two-dimensional black hole.).
Another important question in constructing a model is how to restore the spacetime supersymmetry. In type II string theories the key role was played by Jacobi’s abstruse identity. Some similar useful theta identities were found by Bilal and Gervais long time ago and were used to construct interesting noncritical superstring models. In fact, our model is closely related to their six-dimensional ($`d=5`$) model, in particular contains the latter as a subsector, although the interpretation is somewhat different.
The remainder of this paper is organized as follows. In sect. 2, we briefly review the construction of $`N=2`$ superconformal algebra based on the noncompact coset $`SL(2,\text{})/U(1)`$ and the geometrical meaning of the free fields used in the realization. In sect. 3, we use theta identities to construct a modular invariant partition function, clarify the $`N=2`$ character content and summarize the features of our model. In sect. 4 we study the lowest mass spectra. Finally, we present our conclusions and suggestions for further work in sect. 5.
## 2 $`SL(2,\text{})/U(1)`$ Kazama-Suzuki model
### 2.1 Free-field realizations
Let us begin with a review of the $`SL(2,\text{})/U(1)`$ Kazama-Suzuki model . We use the following free-field realization of the $`SL(2,\text{})`$ current algebra:
$`J^3(z)`$ $`=`$ $`i\sqrt{{\displaystyle \frac{k}{2}}}\varphi ,`$
$`J^\pm (z)`$ $`=`$ $`i\left(\sqrt{{\displaystyle \frac{k}{2}}}\theta \pm i\sqrt{{\displaystyle \frac{k2}{2}}}\rho \right)\mathrm{exp}\left(\pm i\sqrt{{\displaystyle \frac{2}{k}}}(\theta \varphi )\right).`$ (2)
The OPEs of the free scalars are $`\rho (z)\rho (0)\theta (z)\theta (0)\mathrm{log}z`$ but $`\varphi (z)\varphi (0)+\mathrm{log}z`$. This realization may be obtained by bosonizing the $`\beta `$-$`\gamma `$ system of the Wakimoto realization followed by a redefinition of the scalars. The same realization was used in . The energy-momentum tensor is
$`T_{SL(2,\text{})}(z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\rho )^2+{\displaystyle \frac{1}{\sqrt{2(k2)}}}^2\rho {\displaystyle \frac{1}{2}}(\theta )^2+{\displaystyle \frac{1}{2}}(\varphi )^2.`$ (3)
The central charge is $`c_{SL(2,\text{})}=3k/(k2)`$.
The $`SL(2,\text{})`$ parafermions are defined by removing the exponentials of $`\varphi `$ from $`J^\pm `$:
$`\psi ^\pm (z)`$ $`=`$ $`i\left(\sqrt{{\displaystyle \frac{1}{2}}}\theta \pm i\sqrt{{\displaystyle \frac{k2}{2k}}}\rho \right)\mathrm{exp}\left(\pm i\sqrt{{\displaystyle \frac{2}{k}}}\theta \right),`$ (4)
where the parafermion fields $`\psi ^\pm `$ are $`\psi _1`$ and $`\psi _1^{}`$ in the usual notation. (See also for a realization of the $`SL(2,\text{})`$ parafermion algebra.)
The energy-momentum tensor of the parafermion theory is
$`T_{SL(2,\text{})/U(1)}(z)`$ $`=`$ $`T_{SL(2,\text{})}(z)\left(+{\displaystyle \frac{1}{2}}(\varphi )^2\right).`$ (5)
The $`N=2`$ superconformal algebra can be obtained by adding back another free boson $`\phi (z)`$ with the OPE $`\phi (z)\phi (0)\mathrm{log}z`$. The currents are given by
$`T_{N=2}(z)`$ $`=`$ $`T_{SL(2,\text{})/U(1)}(z){\displaystyle \frac{1}{2}}(\phi )^2,`$
$`T_F^\pm (z)`$ $`=`$ $`\sqrt{{\displaystyle \frac{2k}{k2}}}\psi ^\pm \mathrm{exp}\left(\pm i\sqrt{{\displaystyle \frac{k2}{k}}}\phi \right),`$ (6)
$`J_{N=2}(z)`$ $`=`$ $`i\sqrt{{\displaystyle \frac{k}{k2}}}\phi .`$
The central charge is $`c_{N=2}=3k/(k2)`$ again. $`k`$ will be set to 3 later.
### 2.2 Representations of classical SL(2,)
A unitary $`SL(2,\text{})/U(1)`$ coset module is constructed by forbidding the $`J^3`$ excitation in a unitary highest (or rather lowest) weight module of the $`SL(2,\text{})`$ current algebra. Thus the lowest $`L_0`$ level states form a unitary representation of (classical) $`SL(2,\text{})`$. (See e.g. for the representations of the classical $`SL(2,\text{})`$ Lie algebra. See also .) The latter can be realized as the differential operators acting on functions on $`S^1`$: <sup>2</sup><sup>2</sup>2These equations correct the inconsistency in the sign convention in ref. .
$`H^3`$ $`=`$ $`ϵ+i{\displaystyle \frac{d}{dx}},`$
$`H^\pm `$ $`=`$ $`ie^{ix}{\displaystyle \frac{d}{dx}}+(ϵ\pm l)e^{ix},`$ (7)
$`C`$ $`=`$ $`{\displaystyle \frac{1}{2}}(H^+H^{}+H^{}H^+)(H^3)^2.`$
$`\{e^{imx}|m\text{}\}`$ form a complete set and
$`H^3e^{imx}`$ $`=`$ $`(m+ϵ)e^{imx},`$
$`H^\pm e^{imx}`$ $`=`$ $`(m+ϵ\pm l)e^{i(m\pm 1)x},`$ (8)
$`Ce^{imx}`$ $`=`$ $`l(l1)e^{imx}.`$
Thus the representations of $`SL(2,\text{})`$ are labeled by $`(l,ϵ)`$. The corresponding equations in the $`SL(2,\text{})`$ current algebra module are
$`J_0^3|m+ϵ>`$ $`=`$ $`(m+ϵ)|m+ϵ>,`$
$`J_0^\pm |m+ϵ>`$ $`=`$ $`(m+ϵ\pm l)|m+ϵ\pm 1>,`$ (9)
$`𝑱|m+ϵ>`$ $`=`$ $`l(l1)|m+ϵ>`$
with $`𝑱=1/2(J_0^+J_0^{}+J_0^{}J_0^+)(J_0^3)^2`$.
Unitary representations $`(l,ϵ)`$ of (classical) $`SL(2,\text{})`$ are known to be classified into the following four cases:<sup>3</sup><sup>3</sup>3 We do not need to consider the universal covering of $`SL(2,\text{})`$ in our $`c_{N=2}=9`$ case. Then $`ϵ`$ takes discrete values.
* Principal unitary series: $`(+\frac{1}{2}+ip,ϵ)`$; $`p\text{}`$, $`ϵ\{0,\frac{1}{2}\}`$.
* Complementary series: $`(l,0)`$; $`0<l<1`$.
* Discrete series $`𝒟_n^+`$: $`(n+ϵ,ϵ)`$; $`ϵ\{0,\frac{1}{2}\}`$, $`n\text{}_{>0}`$ if $`ϵ=0`$, $`n\text{}_0`$ if $`ϵ=\frac{1}{2}`$.
* Discrete series $`𝒟_n^{}`$: $`(nϵ,ϵ)`$; $`n\text{}_0`$, $`ϵ\{0,\frac{1}{2}\}`$.
* Trivial representation.
In the cases i, ii, the whole module consists of the states $`|m+ϵ>`$, $`m\text{}`$. There are neither highest weight states nor lowest-weight states. In the cases iii$`\pm `$, the module spanned by $`|m+ϵ>`$, $`m\text{}`$ turns reducible due to the appearance of null states. In this case the irreducible submodules $`𝒟_{lϵ}^\pm `$ are unitary. $`𝒟_{lϵ}^+`$ ($`𝒟_{l+ϵ}^{}`$ ) has a lowest- (highest-)weight state $`|l>`$ ($`|l>`$). The quotient module divided by $`𝒟_{lϵ}^+𝒟_{l+ϵ}^{}`$ is finite, and non-unitary if $`l\frac{1}{2},1`$. If $`l=\frac{1}{2}`$ the quotient module is empty; if $`l=1`$ the quotient module consists of a single state, and hence corresponds to the trivial representation. One can also construct unitary representations by starting from negative $`n`$, but they only give equivalent representations. ($`𝒟_0^+`$ is equivalent to $`𝒟_1^+`$ if $`ϵ=0`$.) This is manifest in the symmetry of the Casimir $`ll+1`$. Finally, the trivial representation maps any element of the Lie algebra to 0.
The vertex operator of the $`SL(2,\text{})`$ current algebra at level $`k`$ corresponding to the state $`|m+ϵ>`$ in the $`(l,ϵ)`$ representation is
$`|m+ϵ>e^{+\sqrt{\frac{2}{k2}}l\rho +i\sqrt{\frac{2}{k}}(m+ϵ)(\theta \varphi )}.`$ (10)
It has
$`L_0^{SL(2,\text{})}={\displaystyle \frac{l^2l}{k2}},J_0^3=m+ϵ.`$ (11)
The corresponding $`N=2`$ vertex operator is
$`e^{+\sqrt{\frac{2}{k2}}l\rho +i\sqrt{\frac{2}{k}}(m+ϵ)\theta +i\frac{2(m+ϵ)}{\sqrt{k(k2)}}\phi }.`$ (12)
It has
$`h=L_0^{N=2}={\displaystyle \frac{(l^2l)+(m+ϵ)^2}{k2}},Q=J_0^{N=2}={\displaystyle \frac{2(m+ϵ)}{k2}}.`$ (13)
### 2.3 Interpretation — the $`SL(2,\text{})`$ WZW model
One of the nice features of the realization (2) is its clear geometrical meanings. To see this, let us write out the $`SL(2,\text{})`$ WZW action
$`S_{\text{WZW}}`$ $`=`$ $`{\displaystyle \frac{k}{8\pi }}{\displaystyle d^2x\text{Tr}_\mu g^\mu g^1}+k\mathrm{\Gamma }(g),`$ (14)
$`\mathrm{\Gamma }(g)`$ $`=`$ $`{\displaystyle \frac{1}{12\pi }}{\displaystyle d^3xϵ^{\overline{\mu }\overline{\nu }\overline{\rho }}\text{Tr}\overline{g}^1_{\overline{\mu }}\overline{g}\overline{g}^1_{\overline{\nu }}\overline{g}\overline{g}^1_{\overline{\rho }}\overline{g}}`$ (15)
using the parameterization
$`g`$ $`=`$ $`\left[\begin{array}{cc}\hfill u+w& v+y\hfill \\ \hfill v+y& uw\hfill \end{array}\right],`$ (18)
$$\begin{array}{cc}u=\mathrm{cosh}𝝆\mathrm{cos}𝒕,\hfill & v=\mathrm{cosh}𝝆\mathrm{sin}𝒕,\hfill \\ w=\mathrm{sinh}𝝆\mathrm{cos}\stackrel{\mathbf{~}}{𝜽},\hfill & y=\mathrm{sinh}𝝆\mathrm{sin}\stackrel{\mathbf{~}}{𝜽}.\hfill \end{array}$$
(The signature of the worldsheet is $`\eta ^{\mu \nu }=\text{diag}[1,+1]`$.) The result is
$`S_{\text{WZW}}`$ $`=`$ $`{\displaystyle \frac{k}{8\pi }}{\displaystyle }d^2x[\sqrt{h}h^{\mu \nu }(_\mu 𝝆_\nu 𝝆\mathrm{cosh}^2𝝆_\mu 𝒕_\nu 𝒕+\mathrm{sinh}^2𝝆_\mu \stackrel{\mathbf{~}}{𝜽}_\nu \stackrel{\mathbf{~}}{𝜽})`$ (19)
$`2ϵ^{\mu \nu }\mathrm{sinh}^2𝝆_\mu 𝒕_\nu \stackrel{\mathbf{~}}{𝜽}]`$
$`=`$ $`{\displaystyle \frac{k}{4\pi }}{\displaystyle d^2x\left[_+𝝆_{}𝝆+_+𝒕_{}𝒕+\mathrm{sinh}^2𝝆_+(𝒕\stackrel{\mathbf{~}}{𝜽})_{}(𝒕+\stackrel{\mathbf{~}}{𝜽})\right]}.`$
The last line is the expression on the flat worldsheet with $`x^\pm =(x^1\pm x^0)/\sqrt{2}`$. $`SL(2,\text{})`$ invariant metric is given by
$`ds^2`$ $`=`$ $`d𝝆^2\mathrm{cosh}^2𝝆d𝒕^2+\mathrm{sinh}^2𝝆d\stackrel{\mathbf{~}}{𝜽}^2`$ (20)
in this parameterization. We now dualize $`\stackrel{\mathbf{~}}{𝜽}`$ by adding
$`{\displaystyle \frac{k}{4\pi }}{\displaystyle d^2xϵ^{\mu \nu }_\mu 𝜽_\nu \stackrel{\mathbf{~}}{𝜽}}`$ (21)
to $`S_{\text{WZW}}`$. The dual action reads
$`S_{\text{WZW}}^{\text{dual}}`$ $`=`$ $`{\displaystyle \frac{k}{8\pi }}{\displaystyle d^2x\sqrt{h}\left[(_\mu 𝝆)^2(_\mu (𝒕+𝜽))^2+\mathrm{tanh}^2𝝆(_\mu 𝜽)^2\right]}`$ (22)
with the dilaton field
$`\mathrm{\Phi }=\mathrm{log}\mathrm{cosh}𝝆+\text{const.}`$ (23)
Thus the target space of the dual sigma model is $`\text{}\times `$ two-dimensional black hole. The cigar geometry may be obtained by simply dropping the “time” coordinate $`𝒕`$+$`𝜽`$. In the region $`𝝆\mathrm{}`$, $`𝝆`$, $`𝜽`$ and $`𝒕+𝜽`$ correspond to the free fields $`\rho `$, $`\theta `$ and $`\varphi `$, respectively. Clearly, $`\rho `$ plays the role of the Liouville field.
## 3 Modular Invariant Partition Function
Having reviewed the $`SL(2,\text{})/U(1)`$ coset construction, we will now construct a modular invariant partition function. In type II string theories, the key equation was Jacobi’s abstruse identity:
$`\vartheta _3^4(0|\tau )\vartheta _4^4(0|\tau )\vartheta _2^4(0|\tau )=0.`$ (24)
Are there any such nice theta identities for us, too? In fact, we may find one in the works on noncritical string theory done by Bilal and Gervais long time ago:
$`\mathrm{\Lambda }_1(\tau )`$ $``$ $`\mathrm{\Theta }_{1,1}(\tau ,0)\left(\vartheta _3^2(0|\tau )+\vartheta _4^2(0|\tau )\right)\mathrm{\Theta }_{0,1}(\tau ,0)\vartheta _2^2(0|\tau )=0,`$ (25)
where
$`\mathrm{\Theta }_{m,1}(\tau ,\nu )`$ $`=`$ $`{\displaystyle \underset{n\text{}}{}}q^{(n+m/2)^2}z^{2(n+m/2)}(m=0,1),`$ (26)
$$q=\mathrm{exp}(2\pi i\tau ),z=\mathrm{exp}(2\pi i\nu )$$
are the level-1 $`SU(2)`$ theta functions. We also use another identity :
$`\mathrm{\Lambda }_2(\tau )`$ $``$ $`\mathrm{\Theta }_{0,1}(\tau ,0)\left(\vartheta _3^2(0|\tau )\vartheta _4^2(0|\tau )\right)\mathrm{\Theta }_{1,1}(\tau ,0)\vartheta _2^2(0|\tau )=0,`$ (27)
which is nothing but the modular $`S`$ transform of (25). The determinant of the coefficient matrix of $`\mathrm{\Theta }_{m,1}(\tau ,0)`$ in (25)(27) consistently vanishes (for they are nonzero functions) due to (24). Their modular properties are
$`\mathrm{\Lambda }_1(\tau +1)`$ $`=`$ $`i\mathrm{\Lambda }_1(\tau ),`$
$`\mathrm{\Lambda }_2(\tau +1)`$ $`=`$ $`\mathrm{\Lambda }_2(\tau ),`$ (28)
and
$`\mathrm{\Lambda }_1(\tau )`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}(3\pi i/4)}{\sqrt{2}\tau ^{3/2}}}\left(\mathrm{\Lambda }_1(1/\tau )+\mathrm{\Lambda }_2(1/\tau )\right),`$
$`\mathrm{\Lambda }_2(\tau )`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}(3\pi i/4)}{\sqrt{2}\tau ^{3/2}}}\left(+\mathrm{\Lambda }_1(1/\tau )+\mathrm{\Lambda }_2(1/\tau )\right).`$ (29)
Thus
$`\left|\mathrm{\Lambda }_1(\tau )/\eta ^3(\tau )\right|^2+\left|\mathrm{\Lambda }_2(\tau )/\eta ^3(\tau )\right|^2`$ (30)
is modular invariant.
Having found a building block, we now construct a partition function. We set $`k=3`$ from now on. For the NS sector, we collect the unitary representations
$`h={\displaystyle \frac{Q^2+1}{4}}+p^2,Q\text{}`$ (31)
and integrate over $`p\text{}`$ with the same weight (Figure.1). Except when $`Q`$ is odd and $`p=0`$, they all have non-degenerate characters:
$`\text{Tr}_{\mathrm{NS}}q^{L_0^{N=2}}z^{J_0^{N=2}}=\mathrm{ch}_{\mathrm{NS}}(h,Q)=q^{p^2+\frac{Q^2+1}{4}}z^Qf^{\mathrm{NS}}(\tau ,z),`$ (32)
where
$`f^{\mathrm{NS}}(\tau ,z)`$ $``$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1+zq^{n1/2})(1+z^1q^{n1/2})}{(1q^n)^2}}\left(=q^{1/8}{\displaystyle \frac{\vartheta _3(\nu |\tau )}{(\eta (\tau ))^3}}\right).`$ (33)
Note that this class of representations precisely correspond to the ones made out of the principal unitary series only.
When
$`Q=Q_m=(2m+1),h=h_m={\displaystyle \frac{Q_m^2+1}{4}}=m^2+m+{\displaystyle \frac{1}{2}}(m\text{})`$ (34)
(that is, $`p=0`$), the representation reaches the boundary of the unitarity region, where its irreducible subspace becomes smaller. The non-degenerate character then decomposes into a sum of two degenerate characters of A<sub>3</sub> type:
$`\mathrm{ch}_{\mathrm{NS}}(h_m,Q_m)=\mathrm{ch}_{\mathrm{A}_3\mathrm{deg}}(h_m,Q_m;m)+\mathrm{ch}_{\mathrm{A}_3\mathrm{deg}}(h_m+|m+1/2|,Q_m+\text{sign}(Q_m);m),`$ (35)
where
$`\mathrm{ch}_{\mathrm{A}_3\mathrm{deg}}(h,Q;r)=q^hz^Q{\displaystyle \frac{f^{\mathrm{NS}}(\tau ,z)}{1+q^{|r+1/2|}z^{\mathrm{sign}(Q)}}}.`$ (36)
The first term of (35) is the irreducible character of the representation $`(h_m,Q_m)`$, while the second is that of the adjacent integer $`Q`$ degenerate representation on the same boundary line (Figure.1). These representations are also summed over. As a result, the points $`(h_m,Q_m)`$ may also be thought of as if they were generic (non-degenerate) ones. Such degenerate representations on the boundary lines of the unitarity region are made out of the discrete series of $`SL(2,\text{})`$.
For the R sector, we similarly consider the following set of representations:
$`h={\displaystyle \frac{Q^2}{4}}+p^2,Q\text{}.`$ (37)
In this case, unless $`Q`$ is even and $`p=0`$, the characters are given by the generic ones:
$`\text{Tr}_\mathrm{R}q^{L_0^{N=2}}z^{J_0^{N=2}}=\mathrm{ch}_\mathrm{R}(h,Q)=q^{p^2+\frac{Q^2}{4}}z^Qf^\mathrm{R}(\tau ,z),`$ (38)
where
$`f^\mathrm{R}(\tau ,z)`$ $``$ $`(z^{1/2}+z^{1/2}){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1+zq^n)(1+z^1q^n)}{(1q^n)^2}}\left(={\displaystyle \frac{\vartheta _2(\nu |\tau )}{(\eta (\tau ))^3}}\right).`$ (39)
We do not shift the $`U(1)`$ charge by $`\pm 1/2`$ in the definition of the Ramond character; if the representation has two degenerate lowest $`L_0`$ states (which is the generic ($`h3/8`$) case), $`Q`$ represents the mean value of $`U(1)`$ charges of the two states. If $`h=3/8`$, the lowest $`L_0`$ state is unique because the other becomes null. Then the $`U(1)`$ charge of the lowest $`L_0`$ state is $`Q1/2`$ (P<sup>+</sup> module).
If $`Q`$ is even and $`p=0`$, the non-degenerate character again can be written as a sum of two irreducible degenerate characters of P$`{}_{}{}^{+}{}_{3}{}^{}`$ type:
$`\mathrm{ch}_\mathrm{R}(h_m,Q_m)=\mathrm{ch}_{\mathrm{P}_3^+\mathrm{deg}}(h_m,Q_m;m)+\mathrm{ch}_{\mathrm{P}_3^+\mathrm{deg}}(h_m+|m|,Q_m+\text{sign}(Q_m+1/2);m),`$ (40)
where
$`Q=Q_m=2m,h=h_m={\displaystyle \frac{Q_m^2}{4}}=m^2(m\text{})`$ (41)
and
$`\mathrm{ch}_{\mathrm{P}_3^+\mathrm{deg}}(h,Q;r)=q^hz^Q{\displaystyle \frac{f^\mathrm{R}(\tau ,z)}{1+q^{|r|}z^{\mathrm{sign}(Q+1/2)}}}`$ (42)
(Figure.2). Again, taking into account the extra degenerate representations (the second term of (40)), the points $`(h_m,Q_m)`$ can be thought of as generic.
The two sets of representations (31) and (37), as well as the extra degenerate representations added to them, transform into each other by a spectral flow.
Which GSO projection should we take? We wish to construct a partition function in which the fermion theta in the four-dimensional Minkowski + ghost sector and the characters of the internal $`N=2`$ sector are combined into the form like (30). Thus we need to have the following GSO projection: Let $`F`$ ($`\overline{F}`$) be the right (left) fermion number of the four-dimensional Minkowski + ghost sector, $`Q`$ ($`\overline{Q}`$) the right (left) $`N=2`$ $`U(1)`$ charge and $`ϵ`$ ($`\overline{ϵ}`$) the parameter which distinguishes the parity of the $`U(1)`$ charge of the right (left) $`N=2`$ ground state (See eq. (12).). Then for the NS sector, we only keep the states with both $`F+Q`$ and $`\overline{F}+\overline{Q}`$ odd, and $`2ϵ+2\overline{ϵ}`$ even. In other words, we keep odd fermion excited states if $`ϵ=\overline{ϵ}=0`$, while we do even states if $`ϵ=\overline{ϵ}=1/2`$. The two parameters $`m`$ and $`ϵ`$ labeling the $`U(1)`$ charge are separately GSO projected. On the other hand, for the R sector, the states with the same $`N=2`$ vacuum $`U(1)`$ parity (i.e. $`ϵ=\overline{ϵ}`$) are similarly paired, but the left-right chirality may or may not be the same. If the chirality is the same, we get a IIB-like model, while if it is opposite, we get a IIA-like model. In addition, the left-right diagonal $`p=\overline{p}`$ are required for modular invariance for both sectors.
With this GSO projection, the total partition function reads
$`Z(\tau )={\displaystyle \frac{d\tau d\overline{\tau }}{\text{Im}\tau }(\text{Im}\tau )^2\left|\eta (\tau )\right|^4(\text{Im}\tau )^{\frac{1}{2}}\left|\eta (\tau )\right|^2\left[\left|\mathrm{\Lambda }_1(\tau )/\eta ^3(\tau )\right|^2+\left|\mathrm{\Lambda }_2(\tau )/\eta ^3(\tau )\right|^2\right]},`$ (43)
where the factor $`(\text{Im}\tau )^{\frac{1}{2}}`$ comes from the diagonal “Liouville-momentum” $`p`$ integration, and $`\left|\eta (\tau )\right|^2`$ from the transverse fermions. The transverse fermion theta has already been taken into account in the last factor and GSO projected with the $`N=2`$ theta together. This is the main result of this paper.
We will now list some of the notable features of our partition function $`Z(\tau )`$ (43):
* It is modular invariant. Modular invariance has been achieved by integrating the “Liouville momenta” $`p`$. This may be understood as a summation over the radial momenta on the cigar. Indeed, the principal unitary series is the only class of representations that corresponds to an $`N=2`$ vertex operator with real $`\rho `$ momentum (apart from the imaginary background charge $`i/\sqrt{2}`$). Consequently, spacetime, which was supposed to be four-dimensional, turns five-dimensional effectively; any “particle” in the four-dimensional world has a continuous spectrum. This agrees with the picture of singular CFTs advocated in ref. .
* It is unitary and tachyon free.
* It is spacetime supersymmetric. $`Z(\tau )`$ is zero in reality. Since the vanishing $`\mathrm{\Lambda }_{1,2}(\tau )`$ are a consequence of the ordinary spectral flow, the spacetime supercharge must be given by the usual one using the bosonized ghost, the fermion-number current and the $`N=2`$, $`U(1)`$ current. This is in contrast to the one in ref. containing a contribution from the “longitudinal” boson.
* It has a graviton. The tensor product of the NS transverse fermion excitations yields a graviton, a dilaton and an anti-symmetric two-form field. They survive the GSO projection. This is the most significant difference between Bilal-Gervais’s model and ours; the graviton comes from the $`F+Q2ϵ`$ odd sector ($`\mathrm{\Lambda }_2(\tau )`$), which is missing in the former. Those fields are massive in the sense that they have $`L_0=\overline{L_0}=1/4`$ even when $`p=\overline{p}=0`$.
* It contains bound states in the spectrum. As we have shown in (35) and (40), the partition function has a contribution from the representations made out of the discrete series of $`SL(2,\text{})`$. They do not have a momentum along the cigar and are regarded as the bound states .
## 4 Mass Spectra
Let us now discuss more in detail the lightest mass spectra of our model. In the last section, we have seen that any four-dimensional particle exhibits a continuous mass spectrum, which is attributed to its momentum along the cigar. Therefore we consider the value of total $`L_0`$ of particles “at rest” along the cigar (i.e. $`p=0`$). This is equivalent to studying masses in five dimensions.
Since the transverse fermion theta and the $`N=2`$ fermion theta are GSO projected together and enter in the partition function symmetrically, $`SO(4)`$ (acting on four real fermions) plays an analogous role to the transverse rotational group (or the little group for massive states) in six dimensions, although our model does not have six-dimensional Poincare invariance. The four-dimensional field content can be conveniently obtained by a dimensional reduction (since the $`N=2`$ fermions carry no spacetime indices).
We first consider the states coming from $`|\mathrm{\Lambda }_1(\tau )|^2`$. The lightest NS-NS fields are four scalars. They are of course common in both types of GSO projection in the R sector. On the other hand, the doubly-degenerate lowest $`L_0`$ states in the R sector cannot be a spinor of $`SO(4)`$ because a pseudo-real spinor needs four components. Thus it can only be a nonchiral Majorana $`SO(2)`$ spinor. Then the lightest R-R fields are a four-dimensional vector and two scalars in either projection. (The IIB- and IIA-like projections yield the same R-R fields here because either of the eigenspaces of the $`SO(4)`$ chirality operator decompose into a direct sum of $`+`$ and $``$ $`SO(2)`$ chirality eigenspaces.) They are massless ($`L_0=\overline{L_0}=0`$). Including fermions, they form a four-dimensional $`N=2`$ $`U(1)`$ vector multiplet \+ a hypermultiplet. They are the same field content as a single $`N=4`$, $`U(1)`$ vector multiplet has, although we have only eight spacetime supersymmetries from the standard supercharge construction. It is interesting that those fields are formally obtained by a dimensional reduction of a six-dimensional (2,0) tensor multiplet or a (1,1) vector multiplet on a flat background. They are the lightest fields of Bilal-Gervais’s closed string model .
The lightest NS-NS bosons from $`|\mathrm{\Lambda }_2(\tau )|^2`$ are a graviton, a dilaton and a 2-form (self-dual + anti-self-dual) as we have seen in the previous section. In the R sector, we have now twice as many states as those in $`|\mathrm{\Lambda }_1(\tau )|^2`$ at the lowest level, and hence may regard them as the components of a single $`SO(4)`$ Weyl spinor. Then in the R-R sector, the IIB-like projection yields four anti-self-dual 2-forms and four scalars, while the IIA-like projection gives four vectors. They are the bosonic fields of a six-dimensional $`N=2`$ (“$`𝒩=1`$”) graviton + a self-dual tensor \+ four anti-self-dual tensor multiplets (IIB-like), and a graviton + a self-dual tensor + four vector multiplets (IIA-like), respectively. Again, they combine into a (2,0) graviton + a tensor multiplets in the former case, and a single (1,1) graviton multiplet in the latter. The four-dimensional field contents are obtained by the dimensional reduction of those fields. They have $`L_0=\overline{L_0}=1/4`$.
## 5 Conclusion
In this paper we have constructed a modular invariant partition function of superstrings on four-dimensional Minkowski space $`\times `$ two-dimensional black hole using the $`N=2`$, $`c=9`$ superconformal characters. Our model may be thought of as a modular invariant extension of Bilal-Gervais’s $`d=5`$ noncritical string model and describes type II strings on a conifold. It is unitary, tachyon free and has a continuous spectrum in the four-dimensional sense.
In the $`\alpha ^{}0`$ limit the cigar becomes very thin and can be replaced by a thin cylinder $`\text{}\times S^1`$ since one would need very high energy to see the effect of the sharp tip. In this case one is left with a four-dimensional $`N=2`$ non-gravitational theory with a vector multiplet and a hypermultiplet. This will give an example of holography proposed in , and the $`SL(2,\text{})`$ coset in our model will provide a regularization of the strong coupling singularity of the linear-dilaton vacuum.
It is also interesting that the four-dimensional massless spectrum (with $`p=0`$) of our model coincides with that of the tensionless strings which arise on the four-dimensional intersection of two M5-branes . To understand its implication we recall that the (0,4) tensionless strings arise in type IIB theory on K3 when K3 gets an ADE singularity , while IIB on such a singular K3 is known to be T-dual to a system of type IIA 5-branes . Thus, in view of this, the coincidence of the spectra may suggest that type II strings on a Calabi-Yau threefold with a conifold singularity have a dual description in terms of two intersecting NS5-branes.
On the other hand, $`\alpha ^{}\mathrm{}`$ means that the area near the tip of the cigar is zoomed in on, and the whole target space looks like a six-dimensional Minkowski space. In this case all the towers of particles affect the low-energy physics, of which we cannot expect any local quantum field theory descriptions.
In the middle region of $`\alpha ^{}`$ in between, the lowest level fields do not completely decouple but interact with some other light fields. In the previous section, we have seen that the first two lightest fields in the IIA-like GSO projection are the dimensional reduction of a six-dimensional (1,1) vector multiplet and a graviton multiplet. Remarkably, they are the field content of $`D=6`$, $`N=4`$ ($`=(1,1)`$) gauged supergravity ! Although the graviton multiplet has $`L_0=\overline{L_0}=1/4`$, perhaps this is linked to the well-known subtlety in defining masslessness of a particle in a curved background (See e.g. .). Indeed, $`L_0+\overline{L_0}`$ corresponds to the Klein-Gordon operator in a flat space, but the correspondence becomes less clear in the Minkowski $`\times `$ cigar geometry. If the massless point is shifted by $`1/4`$, the graviton becomes massless but the lightest fields from $`|\mathrm{\Lambda }_1(\tau )|^2`$ then have negative mass square. They can, however, still remain stable if they are above the Breitenlohner-Freedman bound . It would be interesting to explore $`D=6`$, $`N=4`$ gauged supergravity on this background. The supergravity interpretation of the IIB-like projection model remains an open question.
Finally, some generalization of our construction to other singular Calabi-Yau spaces can be done when the corresponding analogue of Jacobi’s abstruse identity is known. For example, an interesting theta identity was found in , which seems to be related to a Calabi-Yau four-fold with a conifold singularity. Very recently, new theta identities corresponding to Calabi-Yau $`n`$-folds with an ADE singularity have been systematically obtained by Eguchi and Sugawara .
I wish to thank T. Eguchi, A. Fujii, T. Kawai, Y. Satoh, Y. Yamada and S.-K. Yang for valuable discussions. I also wish to thank K. Fujikawa, Y. Matsuo, Y. Sugawara and all participants of the seminar given at Dept. of Physics, University of Tokyo in Nov. 1999 for questions and comments, which were useful for completing the final version. I am also grateful for the hospitality of Dept. of Math., Kobe University where a part of this work was done.
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# Series expansions for lattice Green functions
## 1 Introduction
Green functions, also known as two-point functions or propagators, are ubiquitous in physics. They appear in classical field theories, quantum mechanics and quantum field theories to name but a few areas. When a discrete space-time approach is used as a regulator, the continuum Green functions become lattice functions. In the problems of condensed matter physics both the continuum and lattice approach are natural. Random walks on a lattice provide another field of application for Green functions, where they are generating functions for the return probability of a random walker on a lattice . Given an isotropic random walk occurring in steps of one unit on a hypercubic lattice, one considers the probability of visiting a given site, or of returning to the starting site, after $`n`$ steps. This problem was first considered by Pólya who showed that this probability tends to one only in one and two dimensions, for $`n`$ tending to infinity. This is the recurrence/transience transition which is an important result in the study of random walks.
The free propagator appears naturally in the perturbative expansion of field theories. A lattice provides a regulator for the infinities in such expansions, and the four-dimensional discrete Green functions appear naturally . The three-dimensional function appears in the study of effective three-dimensional theories . A lattice regularization has also been used in to study non-relativistic quantum scattering in two and three space dimensions. For dimensions greater than one, the continuum free Green function is singular at the origin, while the lattice version is finite. This provides a regularization of the UV divergences. In yet another context, the propagator was used to find the resistance of a network of resistors . (See also in relation to conformal invariance.)
Although the integral representation of the lattice Green function for the hypercubic lattice is well-known, the analytic calculation of this integral remains a challenge. A closed form appears to exist only in one dimension and for a special case in two dimensions. In higher dimensions there are several partial results . A rather recent numerical approach was used in and developed in for the massless free propagator. It was applied to the two-dimensional case in . Similar position space methods were considered earlier in , and later in .
In this paper I consider the free propagator as a function of the mass squared which is allowed to vary in the whole complex plane. The space coordinates are treated as parameters. Series expansions are first obtained for large values of the mass. In one and two dimensions these series are hypergeometric while in higher dimensions they are not. This large mass expansion follows from the multidimensional integral representation of the lattice Green function. However this representation does not allow one to find the series expansion around other values of the mass. To this end recurrence relations and corresponding differential equations in the mass squared are derived. This approach gives the singularities of the Green function and allows one to expand around any point. Non-trivial monodromies are found around the singular points.
The paper is organized as follow. Section 2 introduces the notation, and the well-known integral representation for the minimal anisotropic lattice Green function for the simple cubic lattice in $`d`$ dimensions. Some general properties are also discussed. The one-dimensional integral representation in terms of Bessel functions is given. A simple derivation for the location of the singularity in the complex mass plane is found. A general expansion of the lattice Green function in the inverse of the mass is obtained, and a recurrence relation between dimensions for the coefficients of this expansion is noted. In section 3 the one-dimensional case is studied in detail as a simple illustration of the methods used. This also serves as a guide to the features common to all dimensions. The same methods are applied to the two dimensional Green function and two novel expansions are found (section 4). The recurrence relation and the differential equation are then obtained for the three-dimensional case (section 5). Section 6 concludes with general remarks for four and higher dimensions. Appendix A contains a set of formulæ used in the text. In appendix B, the method of is extended to arbitrary mass and anisotropies.
## 2 General results
Consider the hypercubic lattice in $`d`$ dimensions, with unit vectors $`\widehat{e}_j`$, $`j=1,\mathrm{},d`$. The discrete anisotropic Green equation is
$$Hf(\stackrel{}{x})\underset{j=1}{\overset{d}{}}\alpha _j\left[f(\stackrel{}{x}+\widehat{e}_j)+f(\stackrel{}{x}\widehat{e}_j)\right]=2\beta f(\stackrel{}{x})+\delta _{\stackrel{}{x},\stackrel{}{0}}$$
(1)
The integers $`x_j`$ label the lattice sites, and the anisotropies $`\alpha _j`$ are arbitrary but non-vanishing complex numbers. When the $`\alpha _j`$’s are omitted, they should be assumed to be all equal to 1. The action of the discrete anisotropic $`d`$-dimensional Laplace operator $`\mathrm{\Delta }`$ on $`f`$ is given by $`(H2_{j=1}^d\alpha _j)f(\stackrel{}{x})`$. Here the contribution of the latter term has been absorbed in $`\beta `$ which is, depending on the context, the mass squared, the energy eigenvalue or a formal expansion parameter for a generating function. This definition of $`\beta `$ gives a natural parity symmetry (6), and renders the sets of singularities symmetric with respect to the origin. The free massless scalar propagator of lattice gauge theories corresponds to the isotropic case, $`\alpha _j=1`$ for all $`j`$, and $`\beta =d`$. A completely isotropic solution of the isotropic equation satisfies
$$2df(\widehat{e}_j)=2\beta f(\stackrel{}{0})+1,j=1,\mathrm{},d$$
(2)
Adding to a solution of (1) any solution of the homogeneous equation,
$$\underset{j=1}{\overset{d}{}}\alpha _j\left[f(\stackrel{}{x}+\widehat{e}_j)+f(\stackrel{}{x}\widehat{e}_j)\right]=2\beta f(\stackrel{}{x}),$$
(3)
yields another Green function. A large class of such solutions can be written as
$$h(\stackrel{}{x}|\stackrel{}{\alpha },\beta )=_\pi ^{+\pi }\mathrm{}_\pi ^{+\pi }\frac{d^d\stackrel{}{q}}{(2\pi )^d}\mathrm{exp}(i\stackrel{}{q}\stackrel{}{x})\stackrel{~}{h}(\stackrel{}{q})\delta (\underset{j=1}{\overset{d}{}}\alpha _j\mathrm{cos}q_j\beta )$$
(4)
where the function $`\stackrel{~}{h}`$ is arbitrary but well-behaved.
One can define a decoupling point at $`\beta =0`$, for which equation (1) becomes an equation for two sub-lattices. A lattice where $`_{j=1}^dx_j`$ is even and one where $`_{j=1}^dx_j`$ is odd.
The minimal solution of (1) is given by
$$G_\pm ^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha },\beta )=_\pi ^{+\pi }\mathrm{}_\pi ^{+\pi }\frac{d^d\stackrel{}{q}}{(2\pi )^d}\frac{\mathrm{exp}(i\stackrel{}{q}\stackrel{}{x})}{2_{j=1}^d\alpha _j\mathrm{cos}q_j2\beta \pm iϵ}$$
(5)
where $`ϵ0^+`$. The $`ϵ`$ prescription removes integration ambiguities at the poles.
When only $`d^{}`$ anisotropy parameters are equal to each other, the solution (5) is invariant under the $`2^dd^{}!`$ parity transformations and permutations of the $`x_j`$’s. For $`d^{}=d`$, the symmetry is that of the $`d`$-dimensional hypercubic group, and the functions $`G_\pm ^{(d)}`$ are completely symmetric in the absolute values of their arguments $`x_j`$. There is also a parity symmetry relating $`\beta `$ to $`\beta `$,
$$G_\pm ^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha },\beta )=e^{i\pi \left(1+X\right)}G_{}^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha },\beta )$$
(6)
where $`X_{j=1}^d|x_j|`$, and a complex conjugation symmetry
$$\left(G_\pm ^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha },\beta )\right)^{}=G_{}^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha }^{},\beta ^{})$$
(7)
The latter symmetry shows that the two Green functions are complex conjugates of each other.
The exponentiation formula
$$\frac{n!}{A^{n+1}}=_0^{\mathrm{}}𝑑tt^n\mathrm{exp}(tA),\mathrm{Re}(A)>0,n=0,1,2,\mathrm{}$$
(8)
and the integral representations (74) for the Bessel functions yield a one-dimensional integral representation for (5):
$$G_\pm ^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha },\beta )=\frac{(\pm i)^{1+X}}{2}_0^+\mathrm{}𝑑t\mathrm{exp}\left(\frac{tϵ}{2}it\beta \right)J_{|x_1|}(\alpha _1t)\mathrm{}J_{|x_d|}(\alpha _dt)$$
(9)
In this expression $`ϵ`$ can be set to zero. For $`\alpha _jR`$, the integrals converge in the domains $`\mathrm{Im}\beta 0`$, without a finite number of real points. Note that (9) satisfies (1) by virtue of properties (7881).
To find the singularities of $`G_\pm ^{(d)}`$ in $`\beta `$, consider first the initial integral representation (5). For $`|\beta |>_j|\alpha _j|`$, the integrand is a continuous function, without singularities, integrated over a compact domain. Thus the Green function has no singularities for these values of $`\beta `$. This includes the point at infinity. Now consider, for simplicity, the case without anisotropies ($`\alpha _j=1`$). Equation (5) implies that the possible singularities are real. When $`d=1`$, the integral $`(\text{9})`$ converges provided the oscillating cosine of the asymptotic expansion (76) is not “canceled” by $`\mathrm{exp}(it\beta )`$. For $`\beta =\pm 1`$, and only for these values, the integral has a diverging contribution of the form $`^{\mathrm{}}𝑑t/\sqrt{t}`$, which results in the branch points $`\beta =\pm 1`$. This is confirmed by the explicit expressions given in section 3. The same reasoning holds for $`d=2`$. The product of the two cosines yields one divergent contribution, $`𝑑t/t`$ for three values of $`\beta `$: $`0`$ and $`\pm 2`$. This is confirmed in section 4. For $`\beta =0`$, note that this approach also predicts a lack of divergence for points on the odd sub-lattice. One has $`\beta \mathrm{ln}\beta `$, which vanishes as $`\beta `$ tends to 0. For $`d3`$, there are enough powers of $`\sqrt{t}`$ to give a converging integral at all $`\beta `$. This corresponds to the recurrence/transience transition in the context of random walks. Let $`[n]`$ denote the integer part of $`n`$. Taking $`[(d1)/2]`$ $`\beta `$-derivatives of (9), and using the asymptotic expansion for the Bessel function, yields $`d+1`$ singularities: $`\beta _0=d,d+2,\mathrm{},d2,d`$. For $`d1`$ and odd, they are of the branch point type: $`(\beta \beta _0)^{p\frac{1}{2}}`$, where $`p`$ is a non-negative integer. For $`d2`$ and even, the singularities are logarithmic of the type: $`(\beta \beta _0)^p^{}\mathrm{ln}^q^{}(\beta \beta _0)`$, where $`p^{}`$ is a non-negative integer and $`q^{}`$ a positive integer. A slightly more complicated but essentially similar analysis applies for arbitrary anisotropies. The location of the singularities will depend explicitly on the $`\alpha _j`$’s. These results can be made rigorous through the use of a Tauberian theorem.
The $`ϵ`$ prescription becomes a convergence factor when one uses the series expansions (75) of the Bessel functions. This gives
$$G_\pm ^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha },\beta )=\frac{1}{2\beta ^{1+X}}\underset{n=0}{\overset{\mathrm{}}{}}c_{X+2n}^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha })\beta ^{2n}$$
(10)
where
$$c_n^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha })_\pi ^{+\pi }\mathrm{}_\pi ^{+\pi }\frac{d^d\stackrel{}{q}}{(2\pi )^d}\mathrm{exp}(i\stackrel{}{q}\stackrel{}{x})\left(\underset{j=1}{\overset{d}{}}\alpha _j\mathrm{cos}q_j\right)^n$$
(11)
and
$`c_n^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha })`$ $`=`$ $`0,n=0,\mathrm{},X1`$ (12)
$`c_{X+2n}^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha })`$ $`=`$ $`\left(X+2n\right)!{\displaystyle \underset{j=1}{\overset{d}{}}}\left({\displaystyle \frac{\alpha _j}{2}}\right)^{|x_j|}{\displaystyle \underset{k_10}{}}\mathrm{}{\displaystyle \underset{k_d0}{}}{\displaystyle \underset{j=1}{\overset{d}{}}}\left({\displaystyle \frac{\alpha _j}{2}}\right)^{2k_j}`$
$`\times {\displaystyle \frac{\delta _{k_1+\mathrm{}+k_d,n}}{k_1!\mathrm{}k_d!(|x_1|+k_1)!\mathrm{}(|x_d|+k_d)!}},n0`$
In particular one finds
$$G_\pm ^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha },\beta )\frac{X!_{j=1}^d\alpha _j^{|x_j|}}{|x_1|!\mathrm{}|x_d|!(2\beta )^{1+X}}\beta \mathrm{}$$
(14)
which shows the Green function to vanish faster than the simple estimate $`1/(2\beta )`$ obtained from the integral representation (5). A more compact form of the coefficients (2) can be obtained by the pairwise replacement of Bessel functions through identity (82). This is done for $`d=2,3,4`$ in the following sections.
The result (10) does not depend on the sign of the $`ϵ`$ prescription. This is easily understood by noticing that for $`|\beta |>_{j=1}^d|\alpha _j|`$ the denominator in (5) does not have poles and therefore $`ϵ`$ can be set to zero. The large $`\beta `$ expansion then yields (10) with the $`c_n`$ coefficients given by (11). One can also conclude that the expansion (102) converges at least for $`|\beta |>_{j=1}^d|\alpha _j|`$. Convergence at the generalized massless point ($`\beta =_{j=1}^d|\alpha _j|`$) depends on the dimensionality of the lattice. This is related to the recurrence/transience of the random walk. In one and two dimensions the series diverge at this point, despite the $`ϵ`$ prescription. In higher dimensions the series converge.
The random walk interpretation of the $`c_n`$’s is the following. For a random walker starting from the origin, let $`P_n(\stackrel{}{x})`$ be the probability of visiting the site $`\stackrel{}{x}`$ after $`n`$ unit steps on the $`d`$-dimensional hypercubic lattice. Take the anisotropies to be positive and such that $`_{j=1}^d\alpha _j=d`$. The probability of jumping from $`\stackrel{}{x}`$ to $`\stackrel{}{x}+\widehat{e}_j`$ or to $`\stackrel{}{x}\widehat{e}_j`$ is $`\frac{\alpha _j}{2d}`$. One then has: $`P_n(\stackrel{}{x})=\frac{1}{d^n}c_n^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha })`$. The vanishing of $`c_n`$ for $`n<X`$ is therefore natural since $`X`$ is the minimal number of steps required to reach point $`\stackrel{}{x}`$.
General relations between the coefficients $`c_n`$ for different dimensions can be simply found. Let $`d^{}`$ be any positive integer smaller than $`d`$. Expanding $`_{j=1}^d\alpha _j\mathrm{cos}q_j`$ using the binomial formula readily yields
$$c_n^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha })=\underset{k=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{k}\right)c_k^{(d^{})}(x_1,\mathrm{},x_d^{}|\alpha _1,\mathrm{},\alpha _d^{})c_{nk}^{(dd^{})}(x_{d^{}+1},\mathrm{},x_d|\alpha _{d^{}+1},\mathrm{},\alpha _d)$$
(15)
Note that these relations are valid for arbitrary anisotropies. One can also get other equations by expanding with the multinomial formulæ.
## 3 The one-dimensional Green function
In one dimension it is possible to obtain closed form expressions, and corresponding series expansions. The anisotropy parameter $`\alpha _1`$ is set to one as it is an irrelevant overall factor.
The closed forms are obtained by integration in the complex $`q`$-plane. The integration contour is the rectangular path $`]\pi +i\mathrm{},\pi ][\pi ,\pi ][\pi ,\pi +i\mathrm{}[]\pi +i\mathrm{},\pi +i\mathrm{}[`$, for the upper half-plane, and its reflection about the real axis for the lower half-plane. The part at infinity gives a vanishing contribution, while the two side contributions cancel each other because $`xx_1`$ is integer. One finds
$$G_\pm ^{(1)}(x|\beta )=\pm \frac{e^{\pm ik|x|}}{2i\mathrm{sin}k},\beta =\mathrm{cos}k,k]0,\pi [$$
(16)
One can compare (16) to its continuum counter-part:
$$g_\pm ^{(1)}(x|k)=_{\mathrm{}}^{\mathrm{}}\frac{dq}{2\pi }\frac{\mathrm{exp}(iqx)}{k^2q^2\pm iϵ}=\pm \frac{e^{\pm ik|x|}}{2ik}$$
(17)
One also finds
$$G_+^{(1)}(x|\beta )=G_{}^{(1)}(x|\beta )=\{\begin{array}{ccc}\hfill +\frac{e^{ik|x|}}{2i\mathrm{sin}k},k_1[0,\pi ],k_2>0& & \\ & & \\ \hfill \frac{e^{ik|x|}}{2i\mathrm{sin}k},k_1]0,\pi [,k_2<0& & \end{array}$$
(18)
where $`\beta =\mathrm{cos}(k_1+ik_2)`$ and $`k=k_1+ik_2`$. Note that the positive exponential $`\frac{e^{+k_2|x|}}{2\mathrm{sinh}k_2}`$ ($`k_2>0`$), despite satisfying equation (1), is not obtained. This was to be expected from the integral representation (5), since for $`|\beta |>1`$ the integrand has no singular point and the integral must vanish as $`|x|\mathrm{}`$.
The general expression (102) reduces to
$`G_\pm ^{(1)}(x|\beta )`$ $`=`$ $`{\displaystyle \frac{1}{(2\beta )^{|x|+1}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(|x|+2k)!}{k!(|x|+k)!}}(2\beta )^{2k}`$
$`=`$ $`{\displaystyle \frac{1}{(2\beta )^{|x|+1}}}_2F_1({\displaystyle \frac{|x|}{2}}+{\displaystyle \frac{1}{2}},{\displaystyle \frac{|x|}{2}}+1;|x|+1;\beta ^2)`$
This series converges uniformly for $`|\beta |>1`$ and simply for $`|\beta |=1,\beta \pm 1`$; it diverges at $`\beta =\pm 1`$ and for $`|\beta |<1`$. The equality of this hypergeometric series to the closed form (18) was otherwise known . The recurrence relation for the coefficients of the Green function, with $`c_nc_n^{(1)}(x|1)`$, reads
$$(n^2x^2)c_nn(n1)c_{n2}=0$$
(20)
The function $`G_\pm ^{(1)}`$ satisfies a second-order differential equation in $`\beta `$, of the hypergeometric type:
$$(\beta ^21)y^{\prime \prime }+3\beta y^{}+(1x^2)y=0$$
(21)
The indices at the three regular singular points $`\beta =\pm 1`$ and $`\mathrm{}`$ are $`(\frac{1}{2},0)`$ and $`(1|x|,1+|x|)`$, respectively. The two branch points $`\beta =\pm 1`$ imply the existence of monodromies in the complex-mass plane. Solution (3) corresponds to $`(\beta =\mathrm{};s=1+|x|)`$. It is then natural to investigate the properties of the solution corresponding to $`(\beta =\mathrm{};s=1|x|)`$:
$`y_1(\beta )`$ $`=`$ $`2^{|x|2}\beta _2^{|x|1}F_1({\displaystyle \frac{1|x|}{2}},1{\displaystyle \frac{|x|}{2}};1|x|;\beta ^2),x0`$ (22)
$`y_{l1}(\beta )`$ $`=`$ $`{\displaystyle \frac{1}{4\beta }}{\displaystyle \frac{1}{\sqrt{1\beta ^2}}}\mathrm{ln}\left({\displaystyle \frac{1}{\beta }}\right)+{\displaystyle \frac{1}{4}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\beta ^{(2n+1)}{\displaystyle \frac{d}{ds}}\left({\displaystyle \frac{\mathrm{\Gamma }\left(n+\frac{s}{2}\right)\mathrm{\Gamma }\left(\frac{s+1}{2}\right)}{\mathrm{\Gamma }\left(n+\frac{s+1}{2}\right)\mathrm{\Gamma }\left(\frac{s}{2}\right)}}\right)_{|s=1}`$ (23)
with $`{}_{2}{}^{}F_{1}^{}(0,\frac{1}{2};0,\beta ^2)1`$. The derivative of the coefficients in (23) can be written using the function $`\psi (z)=\frac{d}{dz}\mathrm{ln}\mathrm{\Gamma }(z)`$. The appearance of the logarithm for this second solution is due the degeneracy of the indices at $`x=0`$. The hypergeometric series (22) truncates to polynomials in $`\beta ^{+1}`$. Note that this solution almost provides a Green function. The a priori arbitrary normalization of a solution was chosen to be equal to $`2^{|x|2}`$ in (22) so that the one-dimensional Green equation, $`f(x+1|\beta )+f(x1|\beta )2\beta f(x|\beta )=\delta _{x,0}`$, is satisfied by $`f(x|\beta )`$ given in (22) and $`f(0|\beta )`$ set to 0.
The expansion around $`\beta =1`$ can be carried out similarly. The two solutions are
$`y_0(\beta )`$ $`=`$ $`{}_{2}{}^{}F_{1}^{}(1|x|,1+|x|;{\displaystyle \frac{3}{2}};{\displaystyle \frac{1\beta }{2}})`$ (24)
$`=`$ $`{\displaystyle \frac{\mathrm{sin}\left(2|x|\mathrm{Arcsin}\left(\sqrt{\frac{1\beta }{2}}\right)\right)}{|x|\sqrt{1\beta ^2}}}\mathrm{for}\beta ]1,1]`$ (25)
$`y_{1/2}(\beta )`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\beta 1}}}_2F_1({\displaystyle \frac{1}{2}}|x|,{\displaystyle \frac{1}{2}}+|x|;{\displaystyle \frac{1}{2}};{\displaystyle \frac{1\beta }{2}})`$ (26)
$`=`$ $`i{\displaystyle \frac{\sqrt{2}\mathrm{cos}\left(2|x|\mathrm{Arcsin}\left(\sqrt{\frac{1\beta }{2}}\right)\right)}{\sqrt{1\beta ^2}}}\mathrm{for}\beta ]1,1[`$ (27)
These solutions are valid for all values of $`x`$. In (27) the choice $`\sqrt{1}=i`$ was made. The solution $`y_0`$ is regular at $`\beta =1`$ and is a solution to the Green equation when appropriately normalized, while $`y_{1/2}`$ is singular at $`\beta =1`$, and is a solution to the homogeneous equation (3). Linear combinations of these solutions provide the two analytic continuations to the function defined by (3):
$$\frac{1}{2\sqrt{2}}(\sqrt{2}|x|y_0\pm y_{1/2})=\pm \frac{1}{2i\mathrm{sin}k}e^{\pm ik|x|},\beta =\mathrm{cos}k,k]0,\pi [$$
(28)
Moreover one finds
$`\beta _2^{|x|1}F_1({\displaystyle \frac{1|x|}{2}},1{\displaystyle \frac{|x|}{2}};1|x|;\beta ^2)`$
$`=2^{1|x|}|x|_2F_1(1|x|,1+|x|;{\displaystyle \frac{3}{2}};{\displaystyle \frac{1\beta }{2}}),|x|1`$ (29)
This corresponds to the polynomial solutions (22) and (24), which have to match since, as polynomials, they are defined on the whole complex $`\beta `$-plane. (An amusing by-product of the foregoing analysis is the identity: $`1=_{n=1}^{\mathrm{}}\frac{(2n+1)!}{2^{3n}(2n1)(n!)^2}`$) The expansion around $`\beta =1`$ yields similar results as can be expected from parity. This symmetry is not explicit on the series representations because the expansion point is not $`\beta =0`$.
The limits $`\beta \pm 1`$ for the Green function do not exist. However the following well-defined limit
$$G^{(1)}(x)\frac{1}{2}\underset{\beta 1^{}}{lim}\left(G_+^{(1)}(x|\beta )+G_{}^{(1)}(x|\beta )\right)=\frac{1}{2}|x|$$
(30)
is also a solution to the Green equation. In fact the latter equation is, for any $`\beta `$, a one-dimensional recurrence relation which can be solved directly by the standard method. The points $`\beta =\pm 1`$ are degeneracy points and the direct solution in this case is: $`f(x)=(\pm 1)^x(f(0)\pm \frac{1}{2}|x|)`$, for $`\beta =\pm 1`$, and $`f(0)`$ arbitrary. For $`\beta =1`$, and with the choice $`f(0)=0`$, one recovers (30).
## 4 The two-dimensional Green function
For arbitrary anisotropies, the identity (82) for the product of two Bessel functions can be used to obtain the following series expansion for the Green function:
$`G_\pm ^{(2)}(\stackrel{}{x}|\stackrel{}{\alpha },\beta )`$ $`=`$ $`{\displaystyle \frac{1}{|x_2|!\mathrm{\hspace{0.17em}2}\beta ^{X+1}}}\left({\displaystyle \frac{\alpha _1}{2}}\right)^{|x_1|}\left({\displaystyle \frac{\alpha _2}{2}}\right)^{|x_2|}`$
$`\times {\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(X+2k)!}{k!(|x_1|+k)!}}_2F_1(k,|x_1|k;|x_2|+1;\left({\displaystyle \frac{\alpha _2}{\alpha _1}}\right)^2)\left({\displaystyle \frac{\alpha _1}{2\beta }}\right)^{2k}`$
When $`\alpha _1=\alpha _2`$, these anisotropies can be set to 1 and the preceding expression becomes
$`G_\pm ^{(2)}(\stackrel{}{x}|\beta )`$ $`=`$ $`{\displaystyle \frac{1}{(2\beta )^{X+1}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(X+k+1)_k(X+2k)!}{k!(|x_1|+k)!(|x_2|+k)!}}(2\beta )^{2k}`$ (32)
$`=`$ $`{\displaystyle \frac{X!}{(2\beta )^{X+1}|x_1|!|x_2|!}}`$
$`\times _4F_3({\displaystyle \frac{X+1}{2}},{\displaystyle \frac{X+1}{2}},{\displaystyle \frac{X}{2}}+1;{\displaystyle \frac{X}{2}}+1;X+1,|x_1|+1,|x_2|+1;4\beta ^2)`$
The series (32) converges uniformly for $`|\beta |>2`$, simply for $`|\beta |=2`$ and $`\beta \pm 2`$ and diverges at $`\beta =\pm 2`$ and $`|\beta |<2`$.
The recurrence relation for the coefficients of the Green function read
$$(n^2X^2)(n^2x^2)c_n4n^2(n1)^2c_{n2}=0$$
(34)
where $`x|x_1||x_2|`$. The Green function $`G_\pm ^{(2)}`$ satisfies a $`4^{\mathrm{th}}`$ order differential equation in $`\beta `$:
$`\beta ^2(\beta ^24)y^{\prime \prime \prime \prime }+\beta (10\beta ^216)y^{\prime \prime \prime }+\left(\beta ^2(2(x_1^2+x_2^2)+25)8\right)y^{\prime \prime }`$
$`+3\beta \left(2(x_1^2+x_2^2)+5\right)y^{}+(1X^2)(1x^2)y=0`$ (35)
The four regular singular points and their indices are:
$`\beta `$ $`=`$ $`0:s=1,1,0,0`$
$`\beta `$ $`=`$ $`\pm 2:s=2,1,0,0`$
$`\beta `$ $`=`$ $`\mathrm{}:s=1+X,\mathrm{\hspace{0.33em}1}X,\mathrm{\hspace{0.33em}1}+x,\mathrm{\hspace{0.33em}1}x`$
The three finite singular points were predicted in section 2. The series (32) is the regular solution corresponding to $`(\beta =\mathrm{};s=1+X)`$. The degeneracy of the indices at the other points signals the existence of logarithmic solutions.
The radius of convergence of the series expansion around $`\beta =0`$ is equal to 2, as $`\beta =\pm 2`$ are the closest singularities. Thus this expansion and the expansion around infinity cover the whole complex $`\beta `$-plane. The recurrence relation for the coefficients of the series expansions around $`\beta =0`$ contains two terms, just like (34) for the expansion around infinity. This permits the explicit determination of the series. One also has to find the connection coefficients involved in the linear combinations of the four solutions which reproduce the Green functions at hand, or the analytic continuations of (32). In the previous section, in (28), one had $`\frac{|x|}{2}`$ and $`\pm \frac{1}{2\sqrt{2}}`$. But the determination of the appropriate functions of $`\stackrel{}{x}`$ is not an easy task and there is no general systematic method which gives a closed form result. Here it is possible to carry out this analysis completely, and I have found the following expansions and connection coefficients:
$$G_\pm ^{(2)}(\stackrel{}{x}|\beta )=[f_0(\stackrel{}{x})\pm l_0(\stackrel{}{x})]y_0(\beta )+[f_1(\stackrel{}{x})\pm l_1(\stackrel{}{x})]y_1(\beta )\pm h_0(\stackrel{}{x})y_{l0}(\beta )\pm h_1(\stackrel{}{x})y_{l1}(\beta )$$
(36)
where
$`y_0(\beta )=_4F_3({\displaystyle \frac{1+X}{2}},{\displaystyle \frac{1X}{2}},{\displaystyle \frac{1+x}{2}};{\displaystyle \frac{1x}{2}};1,{\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{2}};{\displaystyle \frac{\beta ^2}{4}})`$ (37)
$`y_1(\beta )=\beta _4F_3({\displaystyle \frac{2+X}{2}},{\displaystyle \frac{2X}{2}},{\displaystyle \frac{2+x}{2}};{\displaystyle \frac{2x}{2}};{\displaystyle \frac{3}{2}},{\displaystyle \frac{3}{2}},1;{\displaystyle \frac{\beta ^2}{4}})`$ (38)
$`y_{l0}(\beta )=y_0(\beta )\mathrm{ln}\beta +{\displaystyle \frac{d}{ds}}_5F_4({\displaystyle \frac{s+1+X}{2}},{\displaystyle \frac{s+1X}{2}},{\displaystyle \frac{s+1+x}{2}};{\displaystyle \frac{s+1x}{2}},1;`$ (39)
$`{\displaystyle \frac{s+2}{2}},{\displaystyle \frac{s+2}{2}},{\displaystyle \frac{s+1}{2}},{\displaystyle \frac{s+1}{2}};{\displaystyle \frac{\beta ^2}{4}})_{|s=0}`$
$`y_{l1}(\beta )=y_1(\beta )\mathrm{ln}\beta +\beta {\displaystyle \frac{d}{ds}}_5F_4({\displaystyle \frac{s+1+X}{2}},{\displaystyle \frac{s+1X}{2}},{\displaystyle \frac{s+1+x}{2}};{\displaystyle \frac{s+1x}{2}},1;`$ (40)
$`{\displaystyle \frac{s+2}{2}},{\displaystyle \frac{s+2}{2}},{\displaystyle \frac{s+1}{2}},{\displaystyle \frac{s+1}{2}};{\displaystyle \frac{\beta ^2}{4}})_{|s=1}`$
are the four independent solutions of the differential equation around $`\beta =0`$. The $`s`$-derivatives can be easily expressed in terms of $`\psi (z)`$, the logarithmic derivative of the Gamma function. The logarithms are taken real for positive $`\beta `$. The six connection functions are given by:
$`f_0(\stackrel{}{x})`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{cos}\left({\displaystyle \frac{\pi }{2}}X\right)\mathrm{cos}\left({\displaystyle \frac{\pi }{2}}x\right)+{\displaystyle \frac{1}{4}}\mathrm{sin}\left({\displaystyle \frac{\pi }{2}}X\right)\mathrm{sin}\left({\displaystyle \frac{\pi }{2}}|x|\right)`$ (41)
$`f_1(\stackrel{}{x})`$ $`=`$ $`{\displaystyle \frac{1}{2}}|x_1^2x_2^2|f_0(\stackrel{}{x})={\displaystyle \frac{1}{2}}|x|Xf_0(\stackrel{}{x})`$ (42)
$`h_0(\stackrel{}{x})`$ $`=`$ $`h_0\mathrm{cos}\left({\displaystyle \frac{\pi }{2}}X\right)\mathrm{cos}\left({\displaystyle \frac{\pi }{2}}x\right)`$ (43)
$`h_1(\stackrel{}{x})`$ $`=`$ $`{\displaystyle \frac{h_0}{2}}|x_1^2x_2^2|\mathrm{sin}\left({\displaystyle \frac{\pi }{2}}X\right)\mathrm{sin}\left({\displaystyle \frac{\pi }{2}}|x|\right)={\displaystyle \frac{h_0}{2}}xX\mathrm{sin}\left({\displaystyle \frac{\pi }{2}}X\right)\mathrm{sin}\left({\displaystyle \frac{\pi }{2}}x\right)`$ (44)
$`l_0(\stackrel{}{x})`$ $`=`$ $`\left(\psi \left({\displaystyle \frac{X+1}{2}}\right)+\psi \left({\displaystyle \frac{|x|+1}{2}}\right)\right)h_0(\stackrel{}{x})`$ (45)
$`l_1(\stackrel{}{x})`$ $`=`$ $`\left(\psi \left({\displaystyle \frac{X}{2}}+1\right)+\psi \left({\displaystyle \frac{|x|}{2}}+1\right)2\right)h_1(\stackrel{}{x})`$ (46)
$`{\displaystyle \frac{h_0}{2}}(|x|+X)\mathrm{sin}\left({\displaystyle \frac{\pi }{2}}X\right)\mathrm{sin}\left({\displaystyle \frac{\pi }{2}}|x|\right)`$
where $`h_0=\frac{i}{2\pi }`$ These functions satisfy a number of equations of which the simplest are:
$`Hf_0(\stackrel{}{x})=\delta _{\stackrel{}{x},\stackrel{}{0}},Hh_0(\stackrel{}{x})=0,Hl_0(\stackrel{}{x})=0`$ (47)
$`Hf_1(\stackrel{}{x})=2f_0(\stackrel{}{x}),Hh_1(\stackrel{}{x})=2h_0(\stackrel{}{x}),Hl_1(\stackrel{}{x})=2l_0(\stackrel{}{x})`$ (48)
$`x_1[f_0(x_1+1,x_2)f_0(x_11,x_2)]=x_2[f_0(x_1,x_2+1)f_0(x_1,x_21)]`$ (49)
Thus $`f_0`$ is a Green function at the decoupling point while $`h_0`$ and $`l_0`$ are solutions of the homogeneous equation at the same point. The functions $`f_1`$, $`h_1`$ and $`l_1`$ appear as potentials for the sources $`f_0`$, $`h_0`$ and $`l_0`$, respectively. The last equation for $`f_0`$ is a strong form of a directional-independence relation. The $`f_0f_1`$ part in (36) is a solution, regular at $`\beta =0`$, of the Green equation. The remaining part is a singular solution of the homogeneous equation, independently from the value of $`h_0`$. Finally, $`h_0(\stackrel{}{x})y_0(\beta )+h_1(\stackrel{}{x})y_1(\beta )`$ is a regular solution of the homogeneous equation.
The parity property is satisfied with $`\mathrm{ln}(1)=\pm i\pi `$, and the complex conjugation symmetry for both $`\beta `$ and $`\beta ^{}`$ not on the branch cut.
Consider now the expansion about $`\beta =2`$, and let $`v=\beta 2`$. The index $`s`$ can take one of the three values already found: $`0`$, $`1`$ or 2. A generic solution is given by
$$y(v)=\underset{n0}{}a_n(s)v^{n+s}$$
(50)
where
$`16(n+s)^2(n+s1)(n+s2)a_n`$ (51)
$`+4(n+s1)(n+s2)\left[5n^2+(10s9)n+5+5s^29s2(x_1^2+x_2^2)\right]a_{n1}`$
$`+2(2n+2s3)(n+s2)\left[2n^2+(4s6)n+5+2s^26s2(x_1^2+x_2^2)\right]a_{n2}`$
$`+(n+s2+X)(n+s2X)(n+s2+x)(n+s2x)a_{n3}=0`$
with $`n0`$ and $`a_{<0}0`$. For $`a_0(s)1`$, the solutions
$`y_2(v)={\displaystyle \underset{n0}{}}a_n(s=2)v^{n+2}`$ (52)
$`y_1(v)={\displaystyle \underset{n0}{}}a_n(s=1)v^{n+1},a_10`$ (53)
$`y_0(v)={\displaystyle \underset{n0}{}}a_n(s=0)v^n,a_1=a_20`$ (54)
are linearly independent and regular at $`\beta =2`$. The fourth solution has the expected logarithmic singularity:
$$y_{l0}(v)=y_0(v)\mathrm{ln}v+\underset{n0}{}v^n\frac{d}{ds}a_n(s)_{|s=0}$$
(55)
The complete Green function is a priori a linear combination of the four solutions. I have not looked for the corresponding connection coefficients. However one can conclude that the logarithmic solution corresponds to a solution of the homogeneous equation (3), and can therefore be dropped without altering the Green property. The remaining piece is the natural regularization at $`\beta =2`$ of the Green function defined by (5). The situation at $`\beta =2`$ is similar.
For $`x_1=x_2`$ and arbitrary anisotropies, Montroll found an expression in terms of Legendre functions . In terms of the definitions adopted here one has:
$$G_\mathrm{M}^{(2)}((x_1,x_1)|(\alpha _1,\alpha _2),\beta )=\frac{1}{2\pi (\alpha _1\alpha _2)^{\frac{1}{2}}}Q_{|x_1|\frac{1}{2}}\left(\frac{\beta ^2(\alpha _1^2+\alpha _2^2)}{2\alpha _1\alpha _2}\right)$$
(56)
where $`Q_\nu (z)`$ is the $`\nu ^{\mathrm{th}}`$ Legendre function of the second kind. These functions have logarithmic singularities at $`z=\pm 1`$. But for $`\nu `$ a half odd-integer, only $`z=1`$ is a singularity:
$$Q_{|x_1|\frac{1}{2}}(1ϵ)\frac{1}{2}\mathrm{ln}\left(\frac{ϵ}{2}\right)\gamma \psi \left(|x_1|+\frac{1}{2}\right)+O(ϵ),ϵ0^+$$
(57)
where $`\gamma =0.577216\mathrm{}`$. Thus one finds a logarithmic divergence at $`\beta =\pm 2`$:
$`G_\mathrm{M}^{(2)}((x_1,x_1)|(1,1),\beta )={\displaystyle \frac{1}{4\pi }}\mathrm{ln}\left({\displaystyle \frac{4\beta ^2}{4}}\right)`$ (58)
$`+{\displaystyle \frac{\gamma }{2\pi }}+{\displaystyle \frac{1}{2\pi }}\psi \left(|x_1|+{\displaystyle \frac{1}{2}}\right)+O\left({\displaystyle \frac{\beta ^24}{2}}\right),\beta 2^{}\mathrm{or}\beta 2^+`$
One also has $`G_\mathrm{M}^{(2)}((x_1,x_1)|(1,1),0)=\frac{1}{4}`$
Montroll’s result should be qualified. It is general in terms of anisotropies, but partial as it applies only to the line $`x_1=x_2`$ and $`|\beta |<2`$. Also, it is only a part of the full Green function, even under these restrictions. Consider $`x_1=x_2`$ in the Green function given by (36). At $`\beta =0`$ this expression diverges, as expected from the analysis of section 2, and thus contradicts (56). However taking the half-sum of the two Green functions one finds:
$$G^{(2)}(\stackrel{}{x}|\beta )=\frac{1}{2}\left(G_+^{(2)}(\stackrel{}{x}|\beta )+G_{}^{(2)}(\stackrel{}{x}|\beta )\right)=f_0(\stackrel{}{x})y_0(\beta )+f_1(\stackrel{}{x})y_1(\beta )$$
(59)
This solution of the Green equation is regular at $`\beta =0`$, and
$$G^{(2)}((x_1,x_1)|\beta )=\frac{1}{4}\mathrm{cos}(\pi x_1)_2F_1(\frac{1}{2}+|x_1|,\frac{1}{2}|x_1|;1;\frac{b^2}{4})$$
(60)
For $`x_1=0`$, one has the following identity:
$${}_{2}{}^{}F_{1}^{}(\frac{1}{2},\frac{1}{2};1;\frac{b^2}{4})=\frac{2}{\pi }K\left(\frac{\beta ^2}{4}\right)$$
(61)
where $`K`$ is the complete elliptic function of the first kind . One also has
$$Q_{|x_1|\frac{1}{2}}\left(\frac{\beta ^22}{2}\right)=\frac{\pi }{2}\mathrm{cos}(\pi x_1)_2F_1(\frac{1}{2}+|x_1|,\frac{1}{2}|x_1|;1;\frac{b^2}{4}),|\beta |<2$$
(62)
which shows that $`G^{(2)}((x_1,x_1)|\beta )`$ and $`G_\mathrm{M}^{(2)}((x_1,x_1)|(1,1),\beta )`$ are equal, up to a factor of $`1`$. The origin of this sign in (56) is unclear. Compare now (56) to the half-sum of the two functions (32) at $`x_1=x_2`$. The former function is even in $`\beta `$ while the latter is odd. Therefore they cannot be equal, and (56) holds only for $`|\beta |<2`$.
## 5 The three-dimensional Green function
Using the identity (82) for the product of two Bessel functions one obtains the following series expansion for the three-dimensional function:
$`G_\pm ^{(3)}(\stackrel{}{x}|\stackrel{}{\alpha },\beta )`$ $`=`$ $`{\displaystyle \frac{1}{|x_2|!\mathrm{\hspace{0.33em}2}\beta ^{X+1}}}\left({\displaystyle \frac{\alpha _1}{2}}\right)^{|x_1|}\left({\displaystyle \frac{\alpha _2}{2}}\right)^{|x_2|}\left({\displaystyle \frac{\alpha _3}{2}}\right)^{|x_3|}`$
$`\times {\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(X+2k)!\left({\displaystyle \frac{\alpha _3}{2\beta }}\right)^{2k}{\displaystyle \underset{p=0}{\overset{k}{}}}{\displaystyle \frac{1}{p!(kp)!(|x_1|+p)!(|x_3|+kp)!}}`$
$`\times \left({\displaystyle \frac{\alpha _1}{\alpha _3}}\right)_2^{2p}F_1(p,|x_1|p;|x_2|+1;\left({\displaystyle \frac{\alpha _2}{\alpha _1}}\right)^2)`$
When the anisotropies $`\alpha _j`$ are equal to 1 the preceding expression simplifies to
$`G_\pm ^{(3)}(\stackrel{}{x}|\beta )`$ $`=`$ $`{\displaystyle \frac{1}{(2\beta )^{X+1}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(X+2k)!}{(2\beta )^{2k}}}`$
$`\times {\displaystyle \underset{p=0}{\overset{k}{}}}{\displaystyle \frac{(|x_1|+|x_2|+p+1)_p}{p!(kp)!(|x_1|+p)!(|x_2|+p)!(|x_3|+kp)!}}`$
The series (5) converges for $`|\beta |3`$. For $`\stackrel{}{x}=\stackrel{}{0}`$ one can write
$`G_\pm ^{(3)}(\stackrel{}{0}|\beta )`$ $`=`$ $`{\displaystyle \frac{1}{2\beta }}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2k)!}{4^k(k!)^2}}u_k\beta ^{2k}`$ (65)
$`u_k`$ $`=`$ $`{\displaystyle \underset{p=0}{\overset{k}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{k}{p}}\right)^2\left({\displaystyle \genfrac{}{}{0pt}{}{2p}{p}}\right)`$ (66)
The value of this series at $`\beta =3`$ was calculated by Watson :
$$G_\pm ^{(3)}(\stackrel{}{0}|3)=\frac{2}{\pi ^2}\left[18+12\sqrt{2}10\sqrt{3}7\sqrt{6}\right]K^2\left((2\sqrt{3})^2(\sqrt{3}\sqrt{2})^2\right)$$
(67)
The numerical value of (67) is: $`0.2527\mathrm{}`$ (The definition of is adopted for the function $`K`$.) At the singular point $`\beta =3`$, this series converges, rather slowly, to the known value (67). The sum of the first 1001 terms gives: $`0.2502\mathrm{}`$. It is amusing to note that the $`u_k`$ ($`k1`$) appear to be divisible by 3, the number of dimensions.
It is possible to derive a recurrence relation for the coefficients of the $`3`$-dimensional Green function. Define the following even homogeneous polynomials
$`\mathrm{\Sigma }_{222}=x_1^2x_2^2x_3^2`$
$`\mathrm{\Sigma }_{422}=x_1^2x_2^2x_3^2(x_1^2+x_2^2+x_3^2)`$
$`\mathrm{\Sigma }_{2i}=x_1^{2i}+x_2^{2i}+x_3^{2i},i=1,2,3,4`$
$`\mathrm{\Sigma }_{(2i)(2i)}=x_1^{2i}x_2^{2i}+x_2^{2i}x_3^{2i}+x_1^{2i}x_3^{2i},i=1,2`$
$`\mathrm{\Sigma }_{(2i)2}=x_1^{2i}x_2^2+x_1^2x_2^{2i}+x_2^{2i}x_3^2+x_2^2x_3^{2i}+x_1^{2i}x_3^2+x_1^2x_3^{2i},i=2,3`$
For $`c_nc_n^{(3)}(\stackrel{}{x}|(1,\mathrm{},1))`$, I have found
$`(n2)(n4)[\mathrm{\Sigma }_84\mathrm{\Sigma }_{62}+6\mathrm{\Sigma }_{44}+4\mathrm{\Sigma }_{422}`$ (68)
$`4n^2(\mathrm{\Sigma }_6\mathrm{\Sigma }_{42}+10\mathrm{\Sigma }_{222})+2n^4(3\mathrm{\Sigma }_4+2\mathrm{\Sigma }_{22})4n^6\mathrm{\Sigma }_2+n^8]c_n`$
$`+4n(n1)^2(n4)[\mathrm{\Sigma }_6+\mathrm{\Sigma }_{42}+6\mathrm{\Sigma }_{222}(n^24n2)\mathrm{\Sigma }_42(3n^24n+2)\mathrm{\Sigma }_{22}`$
$`+n(n2)(5n^26n+8)\mathrm{\Sigma }_2n(n2)(n^2+2)(3n^26n+4)]c_{n2}`$
$`+2n(n1)(n2)(n3)[(n^24n+12)\mathrm{\Sigma }_4+2(5n^220n+12)\mathrm{\Sigma }_{22}`$
$`2n(n4)(7n^224n+28)\mathrm{\Sigma }_2+n(n4)(15n^496n^3+268n^2384n+248)]c_{n4}`$
$`+4n^2(n1)(n2)(n3)(n4)(n5)\left[3(n3)\mathrm{\Sigma }_2(7n^357n^2+158n162)\right]c_{n6}`$
$`+9n^2(n1)(n2)^2(n3)(n4)(n5)(n6)(n7)c_{n8}=0`$
This translates into a $`10^{\mathrm{th}}`$ order differential equation for $`y=G_\pm ^{(3)}`$:
$`\beta ^2(\beta ^21)^3(\beta ^29)y^{(10)}+\beta (\beta ^21)^2(61\beta ^4418\beta ^2+45)y^{(9)}`$ (69)
$`(\beta ^21)\left[\beta ^6(4\mathrm{\Sigma }_21433)+\beta ^4(16\mathrm{\Sigma }_2+7511)+\beta ^2(12\mathrm{\Sigma }_22673)+27\right]y^{(8)}`$
$`4\beta \left[\beta ^6(42\mathrm{\Sigma }_24167)+\beta ^4(146\mathrm{\Sigma }_2+17284)+\beta ^2(122\mathrm{\Sigma }_211683)18\mathrm{\Sigma }_2+1140\right]y^{(7)}`$
$`+[\beta ^6(6\mathrm{\Sigma }_4+4\mathrm{\Sigma }_{22}2552\mathrm{\Sigma }_2+102963)+\beta ^4(4\mathrm{\Sigma }_424\mathrm{\Sigma }_{22}+5740\mathrm{\Sigma }_2261972)`$
$`+\beta ^2(2\mathrm{\Sigma }_4+20\mathrm{\Sigma }_{22}2444\mathrm{\Sigma }_2+90750)+48\mathrm{\Sigma }_21536]y^{(6)}`$
$`\beta [\beta ^4(162\mathrm{\Sigma }_4108\mathrm{\Sigma }_{22}+17640\mathrm{\Sigma }_2337617)`$
$`+\beta ^2(76\mathrm{\Sigma }_4+392\mathrm{\Sigma }_{22}23180\mathrm{\Sigma }_2+470364)+14\mathrm{\Sigma }_4140\mathrm{\Sigma }_{22}+3812\mathrm{\Sigma }_264194]y^{(5)}`$
$`2[\beta ^4(2\mathrm{\Sigma }_62\mathrm{\Sigma }_{42}+20\mathrm{\Sigma }_{222}714\mathrm{\Sigma }_4476\mathrm{\Sigma }_{22}+28742\mathrm{\Sigma }_2278744)`$
$`+\beta ^2(2\mathrm{\Sigma }_62\mathrm{\Sigma }_{42}12\mathrm{\Sigma }_{222}+202\mathrm{\Sigma }_4+892\mathrm{\Sigma }_{22}`$
$`18870\mathrm{\Sigma }_2^{}+179166)+17\mathrm{\Sigma }_474\mathrm{\Sigma }_{22}+590\mathrm{\Sigma }_25055]y^{(4)}`$
$`4\beta [\beta ^2(16\mathrm{\Sigma }_616\mathrm{\Sigma }_{42}+160\mathrm{\Sigma }_{222}1230\mathrm{\Sigma }_4820\mathrm{\Sigma }_{22}+20776\mathrm{\Sigma }_2103135)`$
$`+8\mathrm{\Sigma }_6^{}8\mathrm{\Sigma }_{42}48\mathrm{\Sigma }_{222}+150\mathrm{\Sigma }_4+612\mathrm{\Sigma }_{22}5190\mathrm{\Sigma }_2+23022]y^{(3)}`$
$`+[\beta ^2(\mathrm{\Sigma }_84\mathrm{\Sigma }_{62}+6\mathrm{\Sigma }_{44}+4\mathrm{\Sigma }_{422}276\mathrm{\Sigma }_6+276\mathrm{\Sigma }_{42}2760\mathrm{\Sigma }_{222}`$
$`+6246\mathrm{\Sigma }_4+4164\mathrm{\Sigma }_{22}44916\mathrm{\Sigma }_2+108681)`$
$`40\mathrm{\Sigma }_6+40\mathrm{\Sigma }_{42}+240\mathrm{\Sigma }_{222}120\mathrm{\Sigma }_4^{}720\mathrm{\Sigma }_{22}+2280\mathrm{\Sigma }_24680]y^{(2)}`$
$`+9\beta [\mathrm{\Sigma }_8^{}4\mathrm{\Sigma }_{62}+6\mathrm{\Sigma }_{44}+4\mathrm{\Sigma }_{422}36\mathrm{\Sigma }_6`$
$`+36\mathrm{\Sigma }_{42}360\mathrm{\Sigma }_{222}+246\mathrm{\Sigma }_4^{}+164\mathrm{\Sigma }_{22}676\mathrm{\Sigma }_2+681)]y^{(1)}`$
$`+15\left[\mathrm{\Sigma }_8^{}4\mathrm{\Sigma }_{62}+6\mathrm{\Sigma }_{44}+4\mathrm{\Sigma }_{422}4\mathrm{\Sigma }_6+4\mathrm{\Sigma }_{42}40\mathrm{\Sigma }_{222}+6\mathrm{\Sigma }_4+4\mathrm{\Sigma }_{22}4\mathrm{\Sigma }_2+1\right]y=0`$
This equation has six regular singular points, 0, $`\pm 1`$, $`\pm 3`$ and $`\mathrm{}`$. The corresponding indices are:
$`\beta `$ $`=`$ $`0:s=7,7,6,5,5,4,3,2,1,0`$
$`\beta `$ $`=`$ $`\pm 1:s=6,5,4,3,2,1,0,{\displaystyle \frac{5}{2}}\text{,}{\displaystyle \frac{3}{2}}\text{,}{\displaystyle \frac{1}{2}}`$
$`\beta `$ $`=`$ $`\pm 3:s=8,7,6,5,4,3,2,1,0,{\displaystyle \frac{1}{2}}`$
$`\beta `$ $`=`$ $`\mathrm{}:s=5,3,\mathrm{\hspace{0.33em}1}+|x_1|+|x_2|+|x_3|,\mathrm{\hspace{0.33em}1}|x_1||x_2||x_3|,`$
$`1+|x_1|+|x_2||x_3|,\mathrm{\hspace{0.33em}1}+|x_1||x_2|+|x_3|,\mathrm{\hspace{0.33em}1}|x_1|+|x_2|+|x_3|,`$
$`1+|x_1||x_2||x_3|,\mathrm{\hspace{0.33em}1}|x_1||x_2|+|x_3|,\mathrm{\hspace{0.33em}1}|x_1|+|x_2||x_3|`$
The Green function is regular at $`\beta =\mathrm{}`$ and corresponds to the index $`1+|x_1|+|x_2|+|x_3|`$. The appearance of indices differing by integer values can result in logarithms. However the 10-term recurrence for the series expansion around $`\beta =\pm 3`$ shows that all 10 solutions do not contain logarithms. The same conclusion holds at $`\beta =\pm 1`$, with an 8-term recurrence. Therefore $`\beta =\pm 1`$ and $`\beta =\pm 3`$ are branch point singularities which are free of logarithms. At $`\beta =0`$ the indices indicate that some solutions contain logarithms. For the foregoing Green function $`\beta =0`$ is a regular point, and one should consider a linear combination of the logarithm-free solutions.
The reason for the appearance of $`\beta =0`$ as a singularity of the differential equation is unclear. Perhaps the fact that this point is a fixed point of the parity symmetry, or its status as a decoupling point may be relevant here.
At its four singular points, the Green function does not diverge as all the indices are non-negative; only the solutions which correspond to the vanishing indices give non-vanishing contributions. However, around a given singular point, one can still consider a natural regularization by dropping all the solutions in the corresponding linear combination which are associated with non-integer indices. Such solutions combine into a solution of the homogeneous equation (3). Finally note that some results were obtained by Joyce , at $`\stackrel{}{x}=\stackrel{}{0}`$ and around $`\beta =3`$.
## 6 Concluding remarks
Using equation (82) twice gives the large-$`\beta `$ series expansion of the four-dimensional Green function:
$`G_\pm ^{(4)}(\stackrel{}{x}|\stackrel{}{\alpha },\beta )`$ $`=`$ $`{\displaystyle \frac{_{j=1}^4\left(\frac{\alpha _j}{2}\right)^{|x_j|}}{|x_2|!|x_4|!}}{\displaystyle \frac{1}{2\beta ^{X+1}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(X+2k)!\left({\displaystyle \frac{\alpha _3}{2\beta }}\right)^{2k}`$
$`\times {\displaystyle \underset{p=0}{\overset{k}{}}}{\displaystyle \frac{1}{p!(kp)!(|x_1|+p)!(|x_3|+kp)!}}\left({\displaystyle \frac{\alpha _1}{\alpha _3}}\right)^{2p}`$
$`\times _2F_1(p,|x_1|p;|x_2|+1;\left({\displaystyle \frac{\alpha _2}{\alpha _1}}\right)^2)`$
$`\times _2F_1((kp),|x_3|(kp);|x_4|+1;\left({\displaystyle \frac{\alpha _4}{\alpha _3}}\right)^2)`$
When the anisotropies are all equal to 1 the preceding expression becomes
$`G_\pm ^{(4)}(\stackrel{}{x}|\beta )`$ $`=`$ $`{\displaystyle \frac{1}{(2\beta )^{(X+1)}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(X+2k)!}{(2\beta )^{2k}}}`$
$`\times {\displaystyle \underset{p=0}{\overset{k}{}}}{\displaystyle \frac{(|x_1|+|x_2|+p+1)_p(|x_3|+|x_4|+kp+1)_{kp}}{p!(kp)!(|x_1|+p)!(|x_2|+p)!(|x_3|+kp)!(|x_4|+kp)!}}`$
The series (6) converges for $`|\beta |4`$. For $`\stackrel{}{x}=\stackrel{}{0}`$ one can write
$`G_\pm ^{(4)}(\stackrel{}{0}|\beta )`$ $`=`$ $`{\displaystyle \frac{1}{2\beta }}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2k)!}{4^k(k!)^2}}v_k\beta ^{2k}`$ (72)
$`v_k`$ $`=`$ $`{\displaystyle \underset{p=0}{\overset{k}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{k}{p}}\right)^2\left({\displaystyle \genfrac{}{}{0pt}{}{2k2p}{kp}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2p}{p}}\right)`$ (73)
For $`\beta =4`$, this series converges to the known value of $`0.1549\mathrm{}`$ . The sum of the first 31 terms gives: $`0.1541\mathrm{}`$. Similarly to the $`u_k`$ in the preceding section, the $`v_k`$ ($`k1`$) also appear to be divisible by the number of dimensions, here equal to 4.
A derivation of the recurrence relation for the coefficients $`c_n^{(d)}(\stackrel{}{x})`$ can be done as for the lower dimensions. The order of the recurrence appears to be larger than 7. From this recurrence a differential equation can be derived. The 5 singular points are of the logarithmic type. The logarithmic solutions, at one given singular point, can be dropped leaving a regular Green function.
These features were seen to be common to the lowest dimensions. They also hold for all the higher dimensions. The coefficients $`c_n^{(d)}(\stackrel{}{x})`$ satisfy recurrence relations for all dimensions. The general form of these relations is easily inferred from the results for the lower dimensions. The coefficients appearing in the recurrence relations are polynomials in $`n`$ and the $`x_j^2`$’s. From such relations one can then derive the differential equation as was done for the lowest dimensions. Another common feature is the possibility of dropping the singular part, around one given singularity. This part is a solution of the homogeneous equation. Finally, at $`\beta =d`$, subtracting $`G_\pm ^{(d)}(\stackrel{}{0}|d)`$ from (5) provides another regularization (see for $`d=2`$).
Determining the explicit recurrence relation and the differential equation is however not a trivial task. It is also difficult to find the $`\stackrel{}{x}`$-dependent coefficients appearing in the linear combination of the solutions around a given singularity. These techniques were applied for the lower dimensions and new results were obtained. While the low dimensions studied in this paper seem to be at the limit of tractability of these methods, the knowledge obtained about the analytic structure of the lattice Green functions in all dimensions is an important step in their study.
Acknowledgements: I would like to thank Ivan Horváth for bringing to my attention reference and for a critical reading of the manuscript. Special thanks to Tim Newman for bringing to my attention reference , for a critical reading of the manuscript, for his lightning derivation of $`\frac{i}{2\pi }`$ and for enjoyable coffee breaks. I thank P. Arnold, P. Fendley and H. Thacker for discussions, and S. Adhikari and G. Moore for communications.
Appendix A: Bessel functions and other formulæ
The cylindrical Bessel functions $`J_n(z)`$ have both integral representations
$$J_n(z)=\frac{i^n}{\pi }_0^\pi 𝑑q\mathrm{exp}(iz\mathrm{cos}q)\mathrm{cos}(nq),nZ$$
(74)
and series expansions
$$J_n(z)=\left(\frac{z}{2}\right)^n\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1)^k}{k!\mathrm{\Gamma }(n+k+1)}\left(\frac{z}{2}\right)^{2k},$$
(75)
with an infinite radius of convergence. The asymptotic behavior at infinity is given by
$$J_n(z)\sqrt{\frac{2}{\pi z}}\left(\mathrm{cos}\chi \frac{4n^21}{8z}\mathrm{sin}\chi \right),|z|\mathrm{},|\mathrm{arg}(z)|<\pi $$
(76)
where $`\chi =z\frac{\pi }{2}n\frac{\pi }{4}`$ One also has
$$_0^{\mathrm{}}J_n(x)𝑑x=1,n0$$
(77)
The Bessel functions have the following properties
$`J_n(z)=(1)^nJ_n(z)nZ`$ (78)
$`J_n(z)=(1)^nJ_n(z)nZ`$ (79)
$`J_0(0)=+1,J_n(0)=0,n0`$ (80)
$`J_{n1}(z)J_{n+1}(z)=2J_n^{}(z)nZ`$ (81)
A particular formula for the product of two Bessel functions is
$$J_m(az)J_n(bz)=\frac{\left(\frac{az}{2}\right)^m\left(\frac{bz}{2}\right)^n}{\mathrm{\Gamma }(n+1)}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1)_2^kF_1(k,mk;n+1;\frac{b^2}{a^2})}{k!\mathrm{\Gamma }(m+k+1)}\left(\frac{az}{2}\right)^{2k}$$
(82)
(A typographical error in has been corrected.) Note that $`{}_{2}{}^{}F_{1}^{}`$ is in fact a polynomial in $`b^2/a^2`$. When $`a=b=1`$ this formula simplifies to
$$J_m(z)J_n(z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1)^k(m+n+k+1)_k}{k!\mathrm{\Gamma }(m+k+1)\mathrm{\Gamma }(n+k+1)}\left(\frac{z}{2}\right)^{m+n+2k}$$
(83)
The generalized hypergeometric series $`{}_{p}{}^{}F_{q}^{}`$ are defined by :
$${}_{p}{}^{}F_{q}^{}(a_1,\mathrm{},a_p;b_1,\mathrm{},b_q;z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(a_1)_k\mathrm{}(a_p)_k}{(b_1)_k\mathrm{}(b_q)_k}\frac{z^k}{k!}$$
(84)
where the Pochhammer symbol $`(a)_k`$ is defined by
$$(a)_0=1,(a)_k=\frac{\mathrm{\Gamma }(a+k)}{\mathrm{\Gamma }(a)}=a(a+1)\mathrm{}(a+k1)$$
(85)
and $`pq+1`$. Define a differential operator $`\delta =z\frac{d}{dz}`$. The function (84) satisfies the differential equation
$$\left[\delta (\delta +b_11)\mathrm{}(\delta +b_q1)z(\delta +a_1)\mathrm{}(\delta +a_p)\right]y(z)=0$$
(86)
For $`p=q+1`$ this equation is Fuchsian with three regular singular points at $`0`$, $`1`$ and $`\mathrm{}`$.
Appendix B: The WLW approach for arbitrary mass
Starting from an observation of C. Vohwinkel, Lüscher and Weisz have developed a powerful algorithmic method for the numerical calculation of the massless lattice Green function in four dimensions . Their method applies immediately to any dimension. Here I show that this method generalizes to arbitrary mass and anisotropies. A conserved quantity for the two-dimensional Green function at the massless and decoupling points is also derived.
An integration by parts of the left-hand side of (87) yields the right-hand side:
$`\alpha _j\left(G_\pm ^{(d)}(\stackrel{}{x}+\widehat{e}_j|\stackrel{}{\alpha },\beta )G_\pm ^{(d)}(\stackrel{}{x}\widehat{e}_j|\stackrel{}{\alpha },\beta )\right)=x_j(\stackrel{}{x}|\stackrel{}{\alpha }),j=1,\mathrm{},d`$ (87)
$`(\stackrel{}{x}|\stackrel{}{\alpha })={\displaystyle _\pi ^\pi }\mathrm{}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^d\stackrel{}{q}}{(2\pi )^d}}\mathrm{exp}(i\stackrel{}{q}.\stackrel{}{x})\mathrm{ln}(2\beta 2{\displaystyle \underset{j=1}{\overset{d}{}}}\alpha _j\mathrm{cos}q_jiϵ)`$ (88)
Equation (1) allows one to find another expression for $``$:
$$(\stackrel{}{x}|\stackrel{}{\alpha })=\frac{2}{_{j=1}^dx_j}\left(\underset{j=1}{\overset{d}{}}\alpha _jG_\pm ^{(d)}(\stackrel{}{x}\widehat{e}_j|\stackrel{}{\alpha },\beta )\beta G_\pm ^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha },\beta )\right)$$
(89)
provided $`_{j=1}^dx_j0`$. This gives the value of $`G_\pm ^{(d)}(\stackrel{}{x}+\widehat{e}_j|\stackrel{}{\alpha },\beta )`$ in terms of $`G_\pm ^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha },\beta )`$ and $`G_\pm ^{(d)}(\stackrel{}{x}\widehat{e}_k|\stackrel{}{\alpha },\beta )`$. The repeated use of these recurrence relations, coupled with the $`\pm x_j`$ invariance, shows that $`G_\pm ^{(d)}(\stackrel{}{x}|\stackrel{}{\alpha },\beta )`$ is a linear combination of the $`2^d`$ values corresponding to $`x_j=0,1`$. Note that all the vertices of the unit hypercube are needed when the anisotropies are arbitrary. These $`2^d`$ values can be calculated numerically, and the particular “$`\pm `$” branch obtained depending on the given value of $`\beta `$ in the complex plane. This generalizes the approach developed in .
One can look for additional conserved quantities as was done in . However this method depends rather strongly on the dimension. For the isotropic two-dimensional case, define
$$g_0(n)=G_\pm ^{(2)}((n,0)|\beta ),g_1(n)=G_\pm ^{(2)}((n,1)|\beta ),n0$$
(90)
The Green equation gives
$$g_0(n+1)+g_0(n1)+2g_1(n)2\beta g_0(n)=0,n1$$
(91)
and an equation inferred from the above approach is
$$g_1(n+1)=\frac{2n}{n+1}\left(\beta g_1(n)g_0(n)\right)\frac{n1}{n+1}g_1(n1),n1$$
(92)
One can look for a conserved quantity in the following form
$`C(n)`$ $`=`$ $`ng_0(n)+a_1ng_1(n)+b_0(n1)g_0(n1)+b_1(n1)g_1(n1)`$ (93)
$`+c_0g_0(n)+c_1g_1(n)+d_0g_0(n1)+d_1g_1(n1),n1`$
Using (91) and (92) one finds that $`C(n)`$ is independent of $`n`$ provided
$$b_0=c_0=d_0=1,a_1=b_1=\frac{1}{\beta 1},c_1=d_1=0$$
(94)
and $`\beta =0`$ or $`\beta =2`$. This form does not allow for other values of $`\beta `$, but a conserved quantity at $`\beta =2`$ can be obtained from the one for $`\beta =2`$ through the parity symmetry. The new quantity is a priori conserved for $`n`$ odd and $`n`$ even separately. It would interesting to find out whether arbitrary values of $`\beta `$ accommodate conserved quantities.
Note that $`\beta =0,\pm 2`$ are the three singular values. Therefore $`C(n)`$ may not be well-defined. However, using the explicit expression of section 4, and taking the limit $`\beta 0`$, one finds that the infinities cancel exactly, leaving $`C_\pm C_\pm (n)=\pm \frac{i}{\pi }`$ for all $`n1`$. For the half-sum the conserved quantity is therefore $`C=0`$. The situation at $`\beta =2`$ is similar. The conserved quantities can be finite through cancellations, and the divergence corresponds to a solution of the homogeneous equation and can therefore be dropped. (See also the remark in the conclusion of concerning this conserved quantity, and .)
Note added: Lattice Green functions also arise in the study of the statistical mechanics of the spherical model . Complex temperature singularities of this system were studied in . Lattice Green functions were also examined for the cases where factorizations in two complete elliptic integrals occur . I would like to thank P. Butera for bringing to my attention the four works cited in this note.
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# On the Structure of Advective Accretion Disks At High Luminosity
## 1 Introduction
Accretion discs are formed when the matter with a large angular momentum is falling into a black hole or another gravitating body. The well known objects where the accretion disks are found are protostar nebulae, binary X-ray sources, cataclysmic variables, active galactic nuclei and others. In this paper we discuss accretion disks around black holes. The standard model of geometrically thin accretion disk has been developed by Shakura (1972), Novikov & Thorne (1973) and Shakura & Sunyaev (1973), and has played a significant role in the development of accretion theory (see Pringle 1981; Frank, King, & Rain 1992; Kato, Mineshige, & Fukue 1998 for reviews). The standard model bases on the vertically averaged approach to equilibrium, and a suggestion of the local thermal balance in which the viscous heating of the gas is balanced by the local radiative cooling. Non-local effects, like the radial advection of thermal energy and the transonic nature of accretion flow, are neglected in the standard model. This simplified approach allows to reduce the general problem to a set of algebraical equations. Such a simple description becomes possible due to an approximate parameterization of the viscosity stress tensor with one non-zero component,
$$t_{r\varphi }=\alpha P,$$
$`(1)`$
suggested by Shakura (1972). The standard model gives a satisfactory appropriate solution of the problem at low accretion rates $`\dot{M}16L_{Edd}/c^2`$, where $`L_{Edd}`$ is the Eddington luminosity.
Simplified solution with inclusion of the advective terms into equations described the vertically integrated models of accretion disks was obtained by Paczyński & Bisnovatyi-Kogan (1981). This approach with some modifications have been used by many researchers to study transonic accretion flows around black holes (Muchotrzeb & Paczyński 1982; Muchotrzeb 1983; Matsumoto et al. 1984; Fukue 1987; Abramowicz et al. 1988; Chen & Taam 1993; Beloborodov 1998). The importance of the transonic nature of the accretion flows on the disk structure has been emphasized by Höshi & Shibazaki (1977), Liang & Thompson (1980) and Abramowicz & Zurek (1981), and later studied in more details by Abramowicz & Kato (1989).
Despite a significant progress in the study of optically thick accretion disks obtained during almost three decades there are a number of unsolved problems still posed in the theory. The problems are connected with a possible non-uniqueness of a solution at $`\alpha 0.01`$ and a non-standard behavior of a singular point type. It was reported by Matsumoto et al. (1984), Muchotrzeb-Czerny (1986) and Abramowicz et al. (1988) that in the case of viscosity prescription (1) the singular point changes its type from a saddle to node when one increases $`\alpha `$. The presence of the nodal-type singular point leads to creating of a possibility of multiple solutions as the authors have claimed. A similar change of the singular point type was reported by Chen & Taam (1993), who used the angular velocity gradient-dependent viscous stress,
$$t_{r\varphi }=\rho \nu r\frac{d\mathrm{\Omega }}{dr},$$
$`(2)`$
where $`\nu `$ is the kinematic viscosity coefficient defined by (14). Narayan, Kato, & Honma (1997) have compared two forms of viscosity (1) and (2) in the case of radiatively inefficient advection-dominated accretion flows. They concluded that the structures of flows corresponded to both viscosities are similar at $`\alpha <0.15`$.
In this paper we show that the mentioned problems have been created by several inconsistencies in the preceding studies. Some problems are connected with an inaccurate averaging of the equations over a disk thickness (Chen & Taam 1993), another ones appear due to an incomplete investigation of the singular points (Abramowicz et al. 1988). We have found that in the case of viscosity prescription (1) a set of equations describing the vertically averaged advective accretion disks has two singular points, independed of $`\alpha `$ and accretion rate. Note, that multiplicity of singular points in solutions for accretion flows in Paczyński-Wiita potential (3) was revealed by Fukue (1987), Chakrabarti & Molteni (1993) and Chakrabarti (1996) in a somewhat different context. We have shown, that at $`\alpha 0.01`$ the inner and outer (with respect to the black hole location) singular points are of the saddle type, and only one integral curve (“separatrix”) which crosses the inner point simultaneously crosses the outer one. This separatrix corresponds to the unique global solution which is determined by two parameters, $`\alpha `$ and $`\dot{m}=\dot{M}c^2/L_{Edd}`$, for a given black hole mass. In Figure 1a the structure of integral curves is schematically represented in the vicinity of the global solution which is shown by the thick line. At larger $`\alpha 0.1`$ the inner singular point is changed to a nodal-type one, while the outer point remains of a saddle-type. There is still one integral curve which goes continuously through both singular points providing a unique global solution, as it is shown in Figure 1b.
In the case of viscosity prescription (2) we have found that there is only one singular point which is always a saddle, and only one physical solution which passes through this point exists. Solutions which correspond to both forms of viscosity (1) and (2) are very close at low $`\alpha `$ limit, $`\alpha 0.1`$.
We have developed a numerical method to solve the set of equations describing the vertically averaged advective accretion disks. The method is based on the standard relaxation technique and explicitly uses conditions at the inner singular point and its vicinity. We have obtained these conditions by expanding a solution into power series around the singular point. Such a modification of the method allows to construct solutions which smoothly pass the singular points and satisfy the regularity conditions at these points with high computer precision in wide range of parameters $`\alpha `$ and $`\dot{m}`$.
The paper is organized as follows. In §2 we formulate a mathematical approach to the problem, write a set of equations, and formulate boundary conditions. In §3 we investigate critical points and discuss behavior of physical values in their vicinity. In §4 we describe our numerical results and discuss them in §5. Details of the numerical method and explicit expansion of physical quantities in the vicinity of the critical points are represented in Appendixes A and B, respectively.
## 2 Problem formulation
We consider a steady geometrically thin accretion disk around a non-rotating black hole. For simplicity, we use the pseudo-Newtonian approach to describe the disk structure in the vicinity of a black hole. In the approach the general relativistic effects are simulated by using Paczyński-Wiita potential (Paczyński & Wiita 1980)
$$\mathrm{\Phi }(r)=\frac{GM}{r2r_g},$$
$`(3)`$
where $`M`$ is the black hole mass and $`2r_g=2GM/c^2`$ is the gravitational radius. The disk self-gravity is neglected.
A general problem of investigation of two-dimensional structure of the accretion disks (in the radial and vertical directions) can be reduced to a one-dimension problem by averaging the disk structure in the vertical direction. In this formulation equations which are described the radial disk structure are written for the midplane density $`\rho `$, pressure $`P`$, radial velocity $`v`$ and angular velocity $`\mathrm{\Omega }`$. The mass conservation equation takes the form,
$$\dot{M}=4\pi rh\rho v,$$
$`(4)`$
where $`\dot{M}`$ is the accretion rate, $`\dot{M}>0`$, and $`h`$ is the disk half-thickness, which is expressed in terms of the isothermal sound speed $`c_s=\sqrt{P/\rho }`$ of gas,
$$h=\frac{c_s}{\mathrm{\Omega }_K}.$$
$`(5)`$
The equations of motion in the radial and azimuthal directions are
$$v\frac{dv}{dr}=\frac{1}{\rho }\frac{dP}{dr}+(\mathrm{\Omega }^2\mathrm{\Omega }_K^2)r,$$
$`(6)`$
$$\frac{\dot{M}}{4\pi }\frac{d\mathrm{}}{dr}+\frac{d}{dr}(r^2ht_{r\varphi })=0,$$
$`(7)`$
where $`\mathrm{\Omega }_K`$ is the Keplerian angular velocity, $`\mathrm{\Omega }_K^2=GM/r(r2r_g)^2`$, $`\mathrm{}=\mathrm{\Omega }r^2`$ is the specific angular momentum and $`t_{r\varphi }`$ is the ($`r`$, $`\varphi `$)-component of the viscous stress tensor. Other components of the stress tensor are assumed to be negligiblly small.
The vertically averaged energy conservation equation can be written in the form<sup>1</sup><sup>1</sup>1The vertical averaging in equation (8) have been done in different way by different authors (compare e.g. Shakura & Sunyaev 1973, and Abramowicz et al. 1988). Our choice of coefficients in (9)-(11), following Chen & Taam (1993) may be not the optimal one. Aposteriory analysis had shown that using 4 in the denominator of (9) instead of 2, would be a more consistent choice, but this change has a little influence on our numerical results.,
$$F_{adv}=F^+F^{},$$
$`(8)`$
where
$$F_{adv}=\frac{\dot{M}}{2\pi r}\left[\frac{dE}{dr}+P\frac{d}{dr}\left(\frac{1}{\rho }\right)\right],$$
$`(9)`$
$$F^+=ht_{r\varphi }r\frac{d\mathrm{\Omega }}{dr},$$
$`(10)`$
$$F^{}=\frac{2aT^4c}{3\kappa \rho h},$$
$`(11)`$
are the advective energy flux, the viscous dissipation rate and the cooling rate per unit surface, respectively, $`T`$ is the midplane temperature, $`\kappa `$ is the opacity and $`a`$ is the radiation constant.
The equation of state for accretion matter consisted of a gas-radiation mixture is
$$c_s^2=T+\frac{1}{3}\frac{aT^4}{\rho },$$
$`(12)`$
where $``$ is the gas constant. The specific energy of the mixture is
$$E=\frac{3}{2}T+\frac{aT^4}{\rho }.$$
$`(13)`$
We consider two prescriptions of viscosity in our models. In one case we adopt a simple relation (1) between the viscous stress and pressure. In another case we assume the angular velocity gradient-dependent viscous stress (2), where the viscosity $`\nu `$ is taken in the form
$$\nu =\frac{2}{3}\alpha c_sh.$$
$`(14)`$
In the limit $`\mathrm{\Omega }\mathrm{\Omega }_K`$ both prescriptions (1) and (2) coincide.
Integrating equation (7) we obtain
$$r^2ht_{r\varphi }=\frac{\dot{M}}{4\pi }(\mathrm{}\mathrm{}_{in}),$$
$`(15)`$
where $`\mathrm{}_{in}`$ is an integration constant and has a meaning of the specific angular momentum of accreting matter at the black hole horizon. Depending on the used viscosity prescription (1) or (2) expression (15) results in an algebraical equation or a first order differential equation, respectively. So, in the case of viscosity prescription (1) we have only two first order differential equations (6) and (8), which require to set two parameters as boundary conditions. In the case of viscosity prescription (2) we obtain additional differential equation from (15) and have to set three parameters as boundary conditions. The integration constant $`\mathrm{}_{in}`$ is chosen to obtain a global transonic solution with a subsonic part at large radii and a supersonic part close to the black hole horizon.
## 3 Investigation of singular points
### 3.1 $`\alpha P`$ viscosity prescription
We consider first the case of viscosity prescription (1). From (15) we obtain the algebraical expression for $`\mathrm{\Omega }`$,
$$\mathrm{\Omega }=\frac{\mathrm{}_{in}}{r^2}+\alpha \frac{c_s^2}{vr}.$$
$`(16)`$
Using (16) the system of equations (6) and (8) can be reduced to the following form,
$$r\frac{v^{}}{v}=\frac{N_1}{D_1},$$
$`(17)`$
$$r\frac{c_s^{}}{c_s}=(1^2)\frac{N_1}{D_1}+1r\frac{\mathrm{\Omega }_K^{}}{\mathrm{\Omega }_K}+\frac{\mathrm{\Omega }^2\mathrm{\Omega }_K^2}{c_s^2}r^2,$$
$`(18)`$
where
$$N_1=\left(1r\frac{\mathrm{\Omega }_K^{}}{\mathrm{\Omega }_K}+\frac{\mathrm{\Omega }^2\mathrm{\Omega }_K^2}{c_s^2}r^2\right)\left(7\frac{3}{2}\beta \frac{1+\beta }{43\beta }\frac{\alpha ^2}{^2}\right)+\left(1r\frac{\mathrm{\Omega }_K^{}}{\mathrm{\Omega }_K}\right)\left(1+\frac{3}{2}\beta \frac{1\beta }{43\beta }\right)+$$
$$\alpha \frac{\mathrm{}_{in}}{vr}+\frac{1}{2}\frac{\alpha ^2}{^2}\frac{1\beta }{\dot{m}}\frac{\mathrm{\Omega }_Kr^2}{c_sr_g},$$
$`(19)`$
$$D_1=(^21)\left(7\frac{3}{2}\beta \frac{1+\beta }{43\beta }\frac{\alpha ^2}{^2}\right)\left(1+\frac{3}{2}\beta \frac{1\beta }{43\beta }+\frac{1}{2}\frac{\alpha ^2}{^2}\right).$$
$`(20)`$
In equations (17)-(20) we use the following notations: $`v^{}dv/dr`$, $`c_s^{}dc_s/dr`$, $`\mathrm{\Omega }_K^{}d\mathrm{\Omega }_K/dr`$, $`\beta =T/c_s^2`$ and $`=v/c_s`$. From (4) and (16) it follows the algebraical equation for $`\beta `$,
$$\beta ^4(1\beta )\frac{3}{4\pi }\frac{\dot{M}\mathrm{\Omega }_K^4}{arvc_s^7}=0.$$
$`(21)`$
The equation $`D_1=0`$ defines singular points of the differential equations (17) and (18), and can be reduced to the following form,
$$\left(7\frac{3}{2}\beta \frac{1+\beta }{43\beta }\right)^4\left(\alpha ^2+8\frac{3\beta ^2}{43\beta }\right)^2+\frac{\alpha ^2}{2}=0.$$
$`(22)`$
Equation (22) is a quadratic equation with respect to $`^2`$ and has two positive roots, which correspond to two singular points:
$$_{1,2}^2=\frac{1}{2}\left[\alpha ^2+8\frac{3\beta _s^2}{43\beta _s}\pm \sqrt{\left(\alpha ^2+8\frac{3\beta _s^2}{43\beta _s}\right)^22\alpha ^2\left(7\frac{3}{2}\beta _s\frac{1+\beta _s}{43\beta _s}\right)}\right]\left(7\frac{3}{2}\beta _s\frac{1+\beta _s}{43\beta _s}\right)^1,$$
$`(23)`$
where $`\beta _s`$ is the value of $`\beta `$ taken at the singular point.
In $`\alpha 1`$ limit we have:
$$_1^2=\left(8\frac{3\beta _s^2}{43\beta _s}\right)\left(7\frac{3}{2}\beta _s\frac{1+\beta _s}{43\beta _s}\right)^1$$
$`(24)`$
and
$$_2^2=\frac{\alpha ^2}{2}\left(8\frac{3\beta _s^2}{43\beta _s}\right)^2.$$
$`(25)`$
The first singular point, in which $`_s=_1`$, locates close to the black hole last stable orbit at $`r=6r_g`$. The corresponding values of $`_1`$ are $`1.118`$ and $`1.069`$ for the gas pressure supported ($`\beta =1`$) and the radiation pressure supported ($`\beta =0`$) accretion flows, respectively. This point is an analogy of the singular point in a spherical flow, where the point divides the subsonic and supersonic regions of accretion flow. The second singular point, in which $`_s=_2`$, located at larger radius, is the result of simplified viscosity prescription (1). We will use the notations $`(r_s)_{in}`$ and $`(r_s)_{out}`$ for positions of the inner and outer singular points, respectively.
At the singular points the numerator $`N_1`$ and denominator $`D_1`$ must simultaneously vanish to provide a regular behavior for a global solution. The type of the singular points must be consistent with a transonic nature of solution. For example, a spiral-type singular point does not satisfy the latter requirement, but a saddle or nodal-type point does it. Detailed analysis of topology in vicinity of singular points was done for thin accretion disks under isothermal approximation by Abramowicz & Kato (1989). They showed that saddle, nodal or spiral types are formally possible, but only saddle and nodal points are physically relevant. This study had confirmed the previously obtained numerical results by Matsumoto et al. (1984). The latter authors demonstrated in a framework of the isothermal accretion disks that the type of the inner singular point is defined by value of $`\mathrm{}_{in}`$. At larger $`\mathrm{}_{in}`$ the point is a spiral, at smaller $`\mathrm{}_{in}`$ there is no inner singular point at all, and only unique choice of $`\mathrm{}_{in}`$ of moderate values corresponds to a saddle or nodal-type singular point. The choice between saddle or nodal-type singular points can be done only by constructing a global model of the disk.
### 3.2 $`\mathrm{\Omega }`$-gradient-dependent viscous stress
In the case of viscosity prescription (2) the differential equations (6), (8) and (15) can be reduced to the following form,
$$r\frac{\mathrm{\Omega }^{}}{\mathrm{\Omega }}=\frac{3}{2}\frac{\mathrm{\Omega }_Krv}{\alpha c_s^2}\left(1\frac{\mathrm{}_{in}}{\mathrm{\Omega }r^2}\right),$$
$`(26)`$
$$r\frac{v^{}}{v}=\frac{N_2}{D_2},$$
$`(27)`$
$$r\frac{c_s^{}}{c_s}=(1^2)\frac{N_2}{D_2}+1r\frac{\mathrm{\Omega }_K^{}}{\mathrm{\Omega }_K}+\frac{\mathrm{\Omega }^2\mathrm{\Omega }_K^2}{c_s^2}r^2,$$
$`(28)`$
where $`\mathrm{\Omega }^{}d\mathrm{\Omega }/dx`$, $`\beta `$ is defined by (21) and
$$N_2=\left(1r\frac{\mathrm{\Omega }_K^{}}{\mathrm{\Omega }_K}+\frac{\mathrm{\Omega }^2\mathrm{\Omega }_K^2}{c_s^2}r^2\right)\left(7\frac{3}{2}\beta \frac{1+\beta }{43\beta }\right)+\frac{3}{4}\frac{\mathrm{\Omega }^2\mathrm{\Omega }_Kr^3v}{\alpha c_s^4}\left(1\frac{\mathrm{}_{in}}{\mathrm{\Omega }r^2}\right)^2+$$
$$\left(1r\frac{\mathrm{\Omega }_K^{}}{\mathrm{\Omega }_K}\right)\left(1+\frac{3}{2}\beta \frac{1\beta }{43\beta }\right)\frac{1\beta }{\dot{m}}\frac{\mathrm{\Omega }_Kr^2}{c_sr_g},$$
$`(29)`$
$$D_2=(^21)\left(7\frac{3}{2}\beta \frac{1+\beta }{43\beta }\right)\left(1+\frac{3}{2}\beta \frac{1\beta }{43\beta }\right).$$
$`(30)`$
There is only one singular point of equations (26)-(28) defined by the equation $`D_2=0`$. The point is an analogy to the inner singular point discussed in §3.1. To be consistent with §3.1 we use notation $`(r_s)_{in}`$ for the position of the point. At $`(r_s)_{in}`$ we have
$$_s^2=\left(8\frac{3\beta _s^2}{43\beta _s}\right)\left(7\frac{3}{2}\beta _s\frac{1+\beta _s}{43\beta _s}\right)^1.$$
$`(31)`$
Note, that the expression for $`_1`$ given by (24) coincides with (31). The latter could mean that the properties of global solutions in the case of viscosity prescriptions (1) and (2) are very similar in the inner part of flow at the limit of small viscosity, $`\alpha 1`$. Our numerical results confirm this conclusion.
Abramowicz & Kato (1989) studied analytically the type of singular point in the isothermal disks in the case of viscosity prescription (2). They showed that the point is always a saddle, and there is no case of a node. This conclusion differs from one obtained in the case of viscosity prescription (1). Our numerical models confirm this dependence of the singular point type on a form of viscosity.
## 4 Numerical results
To be used in the numerical method the sets of differential equations (17)-(18) and (26)-(28) have been re-written in the dimensionless form using the following dimensionless quantities: $`\stackrel{~}{r}=r/r_g`$, $`\stackrel{~}{v}=v/c`$, $`\stackrel{~}{c}_s=c_s/c`$, $`\stackrel{~}{\mathrm{\Omega }}=\mathrm{\Omega }r_g/c`$, $`\stackrel{~}{\mathrm{}}=\mathrm{}/r_gc`$, $`m=M/M_{}`$. In the subsequent discussions we will use these dimensionless quantities skipping in the notations the ‘tilde’ mark. The used method is described in Appendix A. We have calculated a number of models varied by the viscosity prescriptions and parameters $`\alpha `$ and $`\dot{m}`$. The black hole mass $`m`$ contributes into the dimensionless equations in the combinations
$$D=\frac{ac^4\kappa GM_{}}{^4}\frac{m}{\dot{m}}.$$
The parameter $`D`$ was taken to be inversely proportional to $`\dot{m}`$ with $`m=10`$, $`=1.6510^8ergg^1K^1`$ and $`\kappa =0.4cm^2g^1`$. The numerical grid covers the radial range form $`r_{in}`$ located at the inner singular point, $`r_{in}=(r_s)_{in}`$, till $`r_{out}10^4r_g`$.
We discuss first the influence of the numerical outer boundary conditions on our models. We have found that the models are insensitive to the specific values of the outer boundary conditions. By fixing $`\alpha `$ and $`\dot{m}`$ the unique global transonic solution is fully determined. This solution also uniquely determines the outer boundary values: two or three values depending on the used viscosity prescription (1) or (2), respectively. In general, we do not know a priori these ‘correct’ boundary values which the global solution passes through, and consequently, our assumed numerical boundary values are quite arbitrary. But, this values should be close enough to the ‘correct’ ones due to reason of numerical stability. Also, the ‘correct’ outer values are very close, but not exactly equal, to the values obtained from the standard model for the Keplerian accretion disks. Calculations show that all numerical solutions which have the same $`\alpha `$ and $`\dot{m}`$, but different boundary values at different $`r_{out}`$, converge to the ‘common’ solution which is not affected by the outer boundary. This ‘common’ solution represents the global solution which we seek. Significant differences between some numerical solution and the ‘common’ one (with relative errors $`10^4`$) are observed only in $`23`$ grid points before the last outer point at $`r_{out}`$. Such a behaviour of the numerical solutions can be explained by special properties of difference equations (A1) at $`\epsilon 1`$. At small $`\epsilon `$ the method is unstable.
Each model is characterized by value of $`\mathrm{}_{in}`$ \[see eq.(15)\] which has a sense of a specific angular moment of matter infalling into black hole. Figure 2 \[panels (a) and (b)\] shows the dependence of $`\mathrm{}_{in}`$ on accretion rate $`\dot{m}`$ for three values of $`\alpha =0.01`$, $`0.1`$ and $`0.5`$, and two forms of viscosity prescription (1) and (2). At low $`\dot{m}1`$ the value of $`\mathrm{}_{in}`$ is independent of $`\dot{m}`$, but weakly varies with $`\alpha `$. In the low $`\alpha `$ case, $`\alpha =0.01`$ and $`0.1`$, the values of $`\mathrm{}_{in}`$ are close to the minimum value of the Keplerian angular momentum, $`(\mathrm{}_K)_{min}=3.6742`$. At high $`\dot{m}0.1`$ the values of $`\mathrm{}_{in}`$ deviate from $`(\mathrm{}_K)_{min}`$ to larger or smaller values depending on $`\alpha `$. In the case $`\alpha =0.01`$ and $`0.1`$ one can see only minor differences between models with different forms of viscosity. But, for large $`\alpha =0.5`$ the difference in values of $`\mathrm{}_{in}`$ increases. Unfortunately, we had been able to calculated only a limited number of models in the case of viscosity prescription (2) due to technical reason (see Appendix A for details), and our comparison of both prescriptions is not complete in this respect.
Figure 2 \[panels (c) and (d)\] shows locations of the inner singular points $`(r_s)_{in}`$ as a function of $`\dot{m}`$ for different values of $`\alpha `$ and two different viscosity prescriptions. Similar to the case of $`\mathrm{}_{in}`$ the models at low $`\dot{m}`$ show a weak dependence of $`(r_s)_{in}`$ on $`\dot{m}`$. In the low $`\alpha `$ models (squares and circles in Figure 2) the values of $`(r_s)_{in}`$ are close to the location of the black hole last stable orbit at $`r=6`$. At high $`\dot{m}0.1`$ the values of $`(r_s)_{in}`$ are decreasing functions of $`\dot{m}`$ in the case of low $`\alpha =0.01`$ and $`0.1`$, and non-monotonically behave in the case of $`\alpha =0.5`$ (triangles in Figure 2).
Figure 3 shows locations of the outer singular points $`(r_s)_{out}`$ in the case of viscosity prescription (1) as a function of $`\dot{m}`$ for different values of $`\alpha `$. Values of $`(r_s)_{out}`$ are an increasing function of $`\dot{m}`$ and show a power-law behaviour at $`\dot{m}3`$. It is interesting to note that values of $`(r_s)_{out}`$ are almost independent of $`\alpha `$.
Examples of the specific angular momentum distribution $`\mathrm{}(r)`$ are shown in Figure 4 for $`\dot{m}=160`$ and three values of $`\alpha =0.01`$ (short-dashed line), $`0.1`$ (solid line) and $`0.5`$ (dotted line). The distributions correspond to viscosity prescription (1). The location of the inner singular points are indicated by the correspondent points on the curves. The Keplerian angular momentum is displayed by the long-dashed line for comparison. Only the low viscosity model with $`\alpha =0.01`$ has a super-Keplerian part in $`\mathrm{}(r)`$. Models with larger viscosity are everywhere sub-Keplerian. Note, that the singular point in the low viscosity model (short-dashed line) locates in the inner sub-Keplerian part of the disk.
Figure 5 shows the dependence of $`\beta _s`$ on $`\dot{m}`$ at the inner singular points. The change of value of $`\beta `$ form 1 to 0 corresponds to the change of a state from the gas pressure to radiative pressure dominated one. The thin disks with $`\beta 1`$ are locally stable, whereas the parts of the disk in which $`\beta 0`$ are thermally and viscously unstable at $`\dot{m}100`$ (Pringle, Rees, & Pacholczyk 1973). At larger $`\dot{m}100`$ the instability can be suppressed by the advection effect (Abramowicz et al. 1988). We have found a weak dependence of $`\beta _s(\dot{m})`$ on the assumed viscosity prescriptions.
Using analysis discussed in Appendix B we have determined a type of singular points in our numerical solutions. In Figures 2-5 the saddle-type points are indicated by the solid squares, circles and triangles. The nodal-type points are represented by the corresponding empty dots in the same figures. In the case of viscosity prescription (1) the solutions have two singular points (see §3.1). The inner singular points, $`(r_s)_{in}`$, can be saddles or nodes depending on values of $`\alpha `$ and $`\dot{m}`$. Note that the change of type from a saddle to nodal one does not introduce any features in solutions. The outer singular points, $`(r_s)_{out}`$, are always of a saddle-type. In the case of viscosity prescription (2) the solutions have only inner singular points (see §3.2) which are always of a saddle-type.
In models with low $`\dot{m}16`$ and low $`\alpha 0.1`$ the values of $`r_s`$ and $`\mathrm{}_{in}`$ are quantitatively very close to the last stable orbit location ($`r_{in}=6`$) and value of $`\mathrm{}_{in}`$ \[$`\mathrm{}_{in}=(\mathrm{}_K)_{min}`$\] assumed in the standard model (Shakura & Sunyaev 1973). The radial structure in our models is also very close to the one for the standard model in the same range of $`\dot{m}`$ and $`\alpha `$. Such a good coincidence means that the advective terms in equations (6) and (8) are negligiblly small in considered models. However, the high $`\alpha `$ models show quite significant deviation from the standard model independently of $`\dot{m}`$ (see Figure 2).
At high accretion rates, $`\dot{m}16`$, the effect of advection becomes significant. We illustrate it by calculating the luminosity $`L`$ of disk in the case of viscosity prescription (1),
$$L=4\pi _{(r_s)_{in}}^{\mathrm{}}F^{}r𝑑r=2L_{Edd}_{(\stackrel{~}{r}_s)_{in}}^{\mathrm{}}\frac{(1\beta )\stackrel{~}{c}_s\stackrel{~}{r}^{1/2}}{\stackrel{~}{r}2}𝑑\stackrel{~}{r},$$
$`(32)`$
where $`F^{}`$ is given by (11). Figure 6 shows calculated dependences of $`L/L_{Edd}`$ on $`\dot{m}`$ for different values of $`\alpha =0.01`$, $`0.1`$, $`0.5`$ (short-dashed, solid and dotted lines, respectively). There is a simple linear relation $`L/L_{Edd}=\eta \dot{m}`$ in the standard model, in which the advection is neglected. The radiative efficiency $`\eta `$ is a constant and equals $`\eta =1/16`$ in the case of gravitational potential (3). We plot this relation by the straight long-dashed line in Figure 6. One can clearly see from the figure that effect of advection results in reduction of luminosities with respect to the one for the standard model at $`\dot{m}16`$. There is a weak dependence of the luminosity on value of viscosity in disks.
Finally note, that our numerical solutions corresponding to viscosity prescription (1) have some resemblance to the results of Abramowicz et al. (1988), but they show important quantitative differences, especially for large $`\alpha `$ and $`\dot{m}`$.
## 5 Discussion
We have obtained unique solutions for structure of advective accretion disk in a wide range of accretion rates and $`\alpha `$-parameters. Both viscosity prescriptions (1) and (2) have been investigated. The solutions corresponding to both prescriptions are very close for $`\alpha 0.1`$, and begin to differ at larger $`\alpha `$. This is connected, probably with larger deviation of the angular velocity $`\mathrm{\Omega }`$ from the Keplerian one, $`\mathrm{\Omega }_K`$, leading to larger difference between $`t_{r\varphi }`$ in both prescriptions. Unfortunately, our comparison of the prescriptions is not complete due to technical problems in calculation of the high viscosity models in the case of viscosity prescription (2).
The main difference of the present study from the previous ones is in using more sophisticated numerical technique which accurately treats the regularity conditions in the inner singular point of equations. We have performed an analytical expansion at the singular point to calculate the derivatives of physical quantities. These derivatives help us to find the proper integral curve passing through the singular point. The approach allow us to avoid numerical instabilities and inaccuracies, appearing when only variables at the singular point, but not its derivatives, are included into a numerical scheme.
We have found different behaviour of integral curves depending on used viscosity prescription. In the case of viscosity prescription (1) there are two singular points located at $`(r_s)_{in}`$ and $`(r_s)_{out}`$. The inner point, $`(r_s)_{in}`$, locates close to the last stable black hole orbit (see Figure 2), and is an analogy of the singular point in spherical flow, where the point divides the subsonic and supersonic regions. The location of the outer point, $`(r_s)_{out}`$, is determined by the accretion rate (see Figure 3). At low $`\alpha 0.1`$ both points are of a saddle-type. Only one integral curve (“separatrix”) simultaneously crosses two saddle-type points, as it is shown in Figure 1a, and corresponds to the global solution which smoothly connects the supersonic innermost region of the accretion disk and the subsonic outer (formally at $`r=\mathrm{}`$) parts. For larger $`\alpha 0.1`$ the inner singular point changes its type to a node. There was suggestion by Muchotrzeb-Czerny (1986) and Abramowicz et al. (1988) that there is no unique solution in this case, because of all integral curves cross the node. Existence of a unique separatrix crossed simultaneously both singular points preserves a uniqueness of the solution in this case (see Figure 1b). The conclusion of Muchotrzeb-Czerny (1986) and Abramowicz et al. (1988) is probably connected with their neglection of outer singular points inherited to the problem.
Matsumoto et al. (1984) used slightly different form of equations (6) and (8), and they found that in this case only one singular point exists and changes type from a saddle to node. The difference with respect to our results arises because of using different form of the pressure gradient force \[the first term on the right hand side of equation (6)\]. We use the pressure and density taken at the equatorial plane in this term, whereas Matsumoto et al. (1984) used the vertically averaged quantities in it. In the latter case the term has the following form,
$$\frac{1}{\mathrm{\Sigma }}\frac{𝒫}{dr},\mathrm{where}\mathrm{\Sigma }=2_0^h\rho 𝑑z\mathrm{and}𝒫=2_0^hP𝑑z.$$
$`(33)`$
The vertically integrated approach (33) introduces the difference because in this case the free terms with $`\alpha ^2`$ in (20) are absent. Formally, it corresponds to location of the outer singular point at infinity, where $`_2=0`$. Such a visible difference is not qualitatively important for the physical solution, because conditions at the outer singular point are only shifted to infinity, and the integral curve itself has little changes. Thus, similar to our results, in the approach by Matsumoto et al. (1984) the inner critical points of both types, a saddle or node, correspond to a unique solution.
In the case of viscosity prescription (2) we have found that one singular point exist. The point is always of a saddle-type and determines a unique solution. Another results had been obtaining by Chen & Taam (1993). They also found that equations have one singular point, but the point changes its type from a saddle to nodal one, depending on $`\alpha `$ and the accretion rate. It is not clear why such a result was obtained. There are two main differences in our equations (6) and (8), and those used by Chen & Taam (1993). First, they used the same vertical averaging for the equation of motion as Matsumoto et al. (1984). Second, they used the vertically averaged energy equation which corresponds to the inappropriate polytropic relation, $`𝒫\mathrm{\Sigma }^\gamma `$, when one neglects terms corresponding to the viscous heating and radiative cooling. Here $`\gamma `$ is the effective adiabatic index, and other notations are similar to those used in (33). Our equation (8) corresponds to the correct polytropic relation, $`P\rho ^\gamma `$. It could be that one of the mentioned differences results in the change of critical point type.
Acknowledgments. This work was supported in part by RFBR through grant 99-02-18180, the Royal Swedish Academy of Sciences, the Danish Natural Science Research Council through grant No 9701841, Danmarks Grundforskningsfond through its support for establishment of the Theoretical Astrophysics Center.
## Appendix A Numerical method
We use the finite-difference method to solve the systems of ordinary differential equations discussed in $`\mathrm{\S }`$3. The method has some resemblance to that used by Igumenshchev, Abramowicz & Novikov (1998). In this approach the problem is reduced to a solution of a system of non-linear algebraical equations written for a each pair of neighboring numerical grid points. The numerical grid $`\{r_i\}`$ extends over about three orders of magnitude in the radial direction. We look for a numerical solution in which the location of the inner grid point $`r_1`$ coincides with the location of the singular point $`(r_s)_{in}`$ near black hole. In our method we approximate the differential equation $`dy/dr=f(r,y)`$ by the following finite differences,
$$\frac{y_iy_{i1}}{r_ir_{i1}}=\epsilon f_{i1}(1\epsilon )f_i,i=1,2,\mathrm{},I,$$
$`(A1)`$
where the function $`y(r)`$ must be replaced by $`v(r)`$, $`c_s^2(r)`$, and additionally by $`\mathrm{\Omega }(r)`$ in the case of viscous prescription (2), $`\epsilon `$ is a parameter, $`I`$ is the number of grid points, and the lower indices indicate the corresponding grid point. The value of parameter $`\epsilon `$ is chosen to provide a stability of the numerical scheme. We use $`\epsilon =1`$ in most of the cases.
We use the Newton-Raphson iteration scheme to solve the set of equations (A1). In general formulation equations (A1) can be represented by $`K`$ functional relations, involving $`K`$ variables $`y_k`$,
$$F_k(y_1,y_2,\mathrm{},y_K)=0,k=1,2,\mathrm{},K,$$
$`(A2)`$
or in the vector notation, $`𝐅(𝐲)=0`$. The $`(n+1)`$-iteration improvement of an approximate solution $`𝐲^n`$ of (A2) has the form,
$$𝐲^{n+1}=𝐲^n+\omega ^n\delta 𝐲^n,$$
$`(A3)`$
where $`\omega ^n`$ is a parameter, $`\omega ^n1`$, and the correction $`\delta 𝐲^n`$ is the solution of the matrix equation
$$𝐉^n\delta 𝐲^n=𝐅^n.$$
$`(A4)`$
In (A4) $`𝐉`$ is the Jacobian matrix, $`J_{lm}F_l/y_m`$. The parameter $`\omega `$ should be chosen to optimize the convergency of the iteration process to a solution. We use the following form of $`\omega `$,
$$\omega =\frac{\eta }{max(\eta ,\mathrm{\Delta })},$$
$`(A5)`$
where the parameter $`\eta =0.03`$ and $`\mathrm{\Delta }`$ is the average relative correction,
$$\mathrm{\Delta }=\frac{1}{K}\underset{k=1}{\overset{K}{}}\left|\frac{\delta y_k}{y_k}\right|.$$
In the case of equations (17), (18) we have $`K=2(I1)`$ functional relations involved $`2I`$ variables, $`v_i`$ and $`(c_s^2)_i`$. Two variables, $`v_I`$ and $`(c_s)_I`$, must be fixed as boundary conditions, when the number of independent variables equals to the number of equations, and system (A2) can be solved. In the case of equations (26)-(28) we have $`K=3(I1)`$ and three boundary values, $`v_I`$, $`(c_s)_I`$ and $`\mathrm{\Omega }_I`$, which are fixed to provide a consistent solution of (A2).
The presence of the singular points $`(r_s)_{in}`$ which coincides with $`r_1`$ introduces additional complications to our method. To satisfy to two regularity conditions,
$$D=0\mathrm{and}N=0,$$
$`(A6)`$
at $`r_1`$ we modify the iteration procedure by the following way. We add the equation $`D=0`$ to the set of equations (A1) together with a new independent variable $`\mathrm{}_{in}`$. Applying the Newton-Raphson iteration scheme we obtain a solution in which $`D=0`$ at $`r_1`$, but $`N0`$ in general. To satisfy the condition $`N=0`$ at the point $`i=1`$ the appropriate choice of $`r_1`$ must be done. We find the correct value of $`r_1`$ using the bisection method in which $`r_1`$ is changed by displacing all the grid. We apply the Newton-Raphson iterations (A4) for each new grid location.
To correctly approximate the differential equations in the first grid interval ($`r_1,r_2`$), where $`r_1=(r_s)_{in}`$, we expand solution at $`r_1`$,
$$v(r)=v_1+v_s^{}(rr_1),$$
$`(A7)`$
$$c_s^2(r)=(c_s^2)_1+(c_s^2)_s^{}(rr_1).$$
$`(A8)`$
The procedure of calculation of the coefficients $`v_s^{}`$ and $`(c_s^2)_s^{}`$ are given in Appendix B. Using (A7) and (A8) we fix the values at the second grid point as follows,
$$v_2=v_1+v_s^{}(r_2r_1),$$
$`(A9)`$
$$(c_s^2)_2=(c_s^2)_1+(c_s^2)_s^{}(r_2r_1).$$
$`(A10)`$
Relations (A9) and (A10) are used instead of the correspondent difference equations in (A1). The modified system of difference equations avoids numerical instabilities connected with the presence of the inner singular point and allows us to obtain a solution crossed continuously this point.
We have found that described numerical procedure becomes unstable in the case of equations (26)-(28) at large values of $`\alpha `$ and $`\dot{m}`$. The numerical instability arises because of influence of equation (26) and results in small-scale oscillations of all quantities. The numerical instability of similar type was found by Beloborodov (1998), who suppressed the oscillations on the level of $`10^3`$ of relative amplitude by applying a smoothing procedure to $`\mathrm{\Omega }`$. Unfortunately, such a smoothing procedure can not be included into our method because of additional coupling of equations (A1) with the regularity conditions (A6).
## Appendix B Expansion at singular point
We consider first the case of equations (26)-(28). We expand the numerator $`N_2`$ and denominator $`D_2`$ at the singular point $`r_s`$, as follows
$$N_2(r)=\left(\frac{N_2}{r}+\frac{N_2}{v}v_s^{}+\frac{N_2}{c_s^2}(c_s^2)_s^{}+\frac{N_2}{\mathrm{\Omega }}\mathrm{\Omega }_s^{}\right)(rr_s),$$
$`(B1)`$
$$D_2(r)=\left(\frac{D_2}{r}+\frac{D_2}{v}v^{}+\frac{D_2}{c_s^2}(c_s^2)^{}\right)(rr_s).$$
$`(B2)`$
In (B1) and (B2) the partial derivatives are taken at $`r=r_s`$, $`v=v_s`$, $`c_s^2=(c_s^2)_s`$ and $`\mathrm{\Omega }=\mathrm{\Omega }_s`$, and we use notations $`v_s=v(r_s)`$, $`(c_s^2)_s=c_s^2(r_s)`$ and $`\mathrm{\Omega }_s=\mathrm{\Omega }(r_s)`$. The ‘prime’ mark means the radial derivative of the correspondent quantity in (B1) and (B2). In (B2) we take into account that $`D_2/\mathrm{\Omega }=0`$. We denote for convenience,
$$\xi _s=r_s\frac{v_s^{}}{v_s},\eta _s=r_s\frac{(c_s^2)_s^{}}{(c_s^2)_s},\chi _s=r_s\frac{\mathrm{\Omega }_s^{}}{\mathrm{\Omega }_s}.$$
The value of $`\chi _s`$ can be defined with help of (26). Substituting (B1) and (B2) into equations (27) and (28) we finally obtain a quadratic equation with respect to $`\xi _s`$:
$$\left(v_s\frac{D_2}{v}+2(c_s^2)_s\frac{D_2}{c_s^2}A\right)\xi _s^2+\left[r_s\frac{D_2}{r}+2(c_s^2)_s\left(\frac{D_2}{c_s^2}B\frac{N_2}{c_s^2}A\right)v_s\frac{N_2}{v}\right]\xi _s$$
$$\left(r_s\frac{N_2}{r}+\mathrm{\Omega }_s\frac{N_2}{\mathrm{\Omega }}\chi _s+2(c_s^2)_s\frac{N_2}{c_s^2}B\right)=0,$$
$`(B3)`$
where we denote
$$A=1_s^2,\mathrm{and}B=\frac{3}{2}+\frac{r_s}{r_s2}+\frac{1}{(c_s^2)_s}\left(\mathrm{\Omega }_s^2r_s^2\frac{r_s}{(r_s2)^2}\right).$$
Having $`\xi _s`$ from (B3) one obtain $`\eta _s=A\xi _s+B`$. The found values of $`\xi _s`$, $`\eta _s`$ and $`\chi _s`$ determine derivatives $`v_s^{}`$, $`(c_s^2)_s^{}`$ and $`\mathrm{\Omega }_s^{}`$ which we use in the numerical procedure discussed in Appendix A. Equation (B3) has two roots. Which root should be used is defined by convergency of the iterations. We have found that only one of the two roots corresponds to the convergent solution. The partial derivatives of $`N_2`$ and $`D_2`$ used in (B3) can be calculated analytically or using numerical differentiation from (29) and (30). The latter method is simpler, but less accurate than the analytic one. The analytic derivation requires some efforts and produces quite long expressions which we do not present here. In the numerical procedure we have used the analytical expressions for derivatives, but in addition, we have checked them using estimates from the numerical differentiation.
To judge the type of critical point we follow the procedure described by Kato et al. (1998). We introduce a new variable $`\tau `$ defined by
$$d\tau =\frac{dr}{rD_2}.$$
$`(B4)`$
From equation (26)-(28) one obtain
$$\frac{d\mathrm{\Omega }}{d\tau }=\chi \mathrm{\Omega }D_2,\frac{dv}{d\tau }=vN_2,\frac{dc_s^2}{d\tau }=2c_s^2(AN_2+BD_2).$$
$`(B5)`$
All of the variables are now expanded around those at $`r_s`$,
$$r=r_s+\mathrm{\Delta }r,v=v_s+\mathrm{\Delta }v,c_s^2=(c_s^2)_s+\mathrm{\Delta }c_s^2,\mathrm{\Omega }=\mathrm{\Omega }_s+\mathrm{\Delta }\mathrm{\Omega }.$$
$`(B6)`$
Substituting (B6) into (B4) and (B5) and retaining only linear terms, one obtain a system of linear differential equations with respect to $`\mathrm{\Delta }r`$, $`\mathrm{\Delta }v`$, $`\mathrm{\Delta }c_s^2`$ and $`\mathrm{\Delta }\mathrm{\Omega }`$. Assuming that these quantities depend on $`\tau `$ in the form $`\mathrm{exp}(\lambda \tau )`$, one obtain the following characteristic equation which determines the eigenvalues $`\lambda `$:
$$\lambda ^2\lambda \left[r_s\frac{D_2}{r}+v_s\frac{N_2}{v}+2(c_s^2)_s\left(A\frac{N_2}{c_s^2}+B\frac{D_2}{c_s^2}\right)\right]$$
$$2Ar_s(c_s^2)_s\left(\frac{N_2}{r}\frac{D_2}{c_s^2}\frac{D_2}{r}\frac{N_2}{c_s^2}+\frac{D_2}{c_s^2}\frac{\chi _s}{r_s}\mathrm{\Omega }_s\frac{N_2}{\mathrm{\Omega }}\right)$$
$$2Bv_s(c_s^2)_s\left(\frac{D_2}{v}\frac{N_2}{c_s^2}\frac{N_2}{v}\frac{D_2}{c_s^2}\right)+$$
$$r_sv_s\left(\frac{D_2}{r}\frac{N_2}{v}\frac{N_2}{r}\frac{D_2}{v}\right)v_s\frac{D_2}{v}\chi _s\mathrm{\Omega }_s\frac{N_2}{\mathrm{\Omega }}=0.$$
$`(B7)`$
Again, all the partial derivatives in (B7) are taken at $`r=r_s`$, $`v=v_s`$, $`c_s^2=(c_s^2)_s`$ and $`\mathrm{\Omega }=\mathrm{\Omega }_s`$. If two solutions of quadratic equation (B7) are real and have different signs, $`\lambda _1\lambda _2<0`$, then the singular point is a saddle. The point is a node if two real roots of (B7) have identical signs, $`\lambda _1\lambda _2>0`$. Complex conjugate roots of (B7) corresponds to the spiral singular point.
In the case of equations (17), (18) equations (B3) and (B7) must be modified by substituting $`N_1`$ and $`D_1`$ instead of $`N_2`$ and $`D_2`$, and assuming $`N_1/\mathrm{\Omega }=0`$. The notation for $`B`$ must be also changed to
$$B=\frac{3}{2}+\frac{r_s}{r_s2}+\frac{1}{(c_s^2)_s}\left[\left(\frac{\mathrm{}_{in}}{r_s}+\alpha \frac{(c_s^2)_s}{v_s}\right)^2\frac{r_s}{(r_s2)^2}\right].$$
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# The infimum, supremum and geodesic length of a braid conjugacy class
## 1. Introduction
The conjugacy problem in the $`n`$-string braid group $`B_n`$ is the following decision problem:
> Given two braids $`\alpha ,\alpha ^{}B_n`$, determine, in a finite number of steps, whether $`\alpha =\gamma \alpha ^{}\gamma ^1`$ for some $`\gamma B_n`$.
In the late sixties Garside solved the (word and) conjugacy problems in $`B_n`$. His solution to both problems was exponential in both word length and braid index. Subsequently, the efficiency of his algorithm was improved by Thurston and Elrifai-Morton to give a solution to the word problem which is polynomial in both word length and braid index.
All three papers , and work with the following well-known presentation of $`B_n`$, which we will call the *old presentation*:
> $`\begin{array}{cc}\text{generators:}\hfill & \sigma _1,\mathrm{},\sigma _{n1}\hfill \\ \text{relations:}\hfill & \sigma _i\sigma _j=\sigma _j\sigma _i,|ij|>1\hfill \\ & \sigma _i\sigma _j\sigma _i=\sigma _j\sigma _i\sigma _j,|ij|=1\hfill \end{array}`$
There is also a parallel and slightly more efficient solution to the word and conjugacy problems in , due to the authors of this paper. It uses a different presentation which we call the *new presentation*:
> $`\begin{array}{cc}\text{generators:}\hfill & a_{ts},nt>s1\hfill \\ \text{relations:}\hfill & a_{ts}a_{rq}=a_{rq}a_{ts},(tr)(tq)(sr)(sq)>0\hfill \\ & a_{ts}a_{sr}=a_{tr}a_{ts}=a_{sr}a_{tr},nt>s>r1\hfill \end{array}`$
The terms old and new are due to Krammer, who used the new presentation in . Both the old and new solutions to the word problem are polynomial in word length and braid index, but the best estimates obtained for the complexity of the solution to the conjugacy problem (see ) were rough exponential bounds. It was clear that better answers could not be obtained without more detailed information about the combinatorics, using either the old or new presentation.
Let $`|W|`$ denote the letter length of $`W`$, as a word in the given set of generators of $`B_n`$. The main result in this note is an algorithm which is polynomial in both $`|W|`$ and $`n`$ for computing two key integer invariants of the conjugacy class $`[W]`$ of $`W`$. The invariants in question are known as the infimum and supremum (or more informally inf and sup), using either presentation. See $`\mathrm{\S }`$2 below for precise definitions. We will also be able to compute the geodesic length (defined in $`\mathrm{\S }`$2,4 below) for the conjugacy class in polynomial time.
The reason we are able to do this requires some explanation. The method for finding inf (resp. sup) in both and rests on a procedure which is known as cycling (resp. decycling). While cycling and decycling are clearly finite processes, it had not been known how many times one must iterate them to either increase inf($`W^{}`$) (resp. decrease sup($`W^{}`$)) for a word $`W^{}[W]`$ or to guarantee that a maximum (resp. minimum) value, denoted by inf($`[W]`$) (resp. sup($`[W]`$)), for the conjugacy class has already been achieved. For the old presentation it had been claimed in that the bound is 1, however an example was given in for which 2 cyclings were needed to increase the infimum. Up to now, there were no published results which gave bounds, except for a very crude estimate in . Our main result in this note is to find upper and lower bounds for the number of times one must cycle (resp. decycle), using either presentation, in order to replace a given word $`W`$ with $`W^{}[W]`$, where $`\text{inf}(W^{})>\text{inf}(W)`$ (resp. $`\text{sup}(W^{})<\text{sup}(W))`$, or be sure that $`W`$ realizes inf($`[W]`$) (resp. sup($`[W]`$)). For the new presentation we will prove that our upper bound is the best possible one.
Here is an outline of this paper. In $`\mathrm{\S }`$2 we review the background and state our results in a precise way. See Theorem 1, Corollary 2 and Corollary 3. In §3, we prove these three results. In §4, we give examples which prove that the bound in Corollary 2 is sharp for the new presentation, with somewhat weaker results for the old. In $`\mathrm{\S }`$5 we discuss the open problem of whether the solutions which we know to the conjugacy problem are polynomial in word length and braid index, and state several conjectures relating to that matter and also to the ‘shortest word problem’ in $`B_n`$, defined in that section.
## 2. Statement of Results
In this section we state our results precisely. To do so we need to review what has already been done. Since almost all the machinery is identical in the two theories, it will be convenient to introduce unified notation, so that we may review both theories at the same time. The symbol $`W`$ will be used to indicate a word in the generators of $`B_n`$, using either presentation. The element and conjugacy class which $`W`$ represents will be denoted $`\{W\}`$ and $`[W]`$. The letter length of $`W`$ is $`|W|`$.
1. Note that the relations in the old and new presentations are equivalences between positive words with same word-length. So the word-length is easy to compute for positive words. Let $`B_n^+`$ be the semigroup defined by the same generators and relations in the given presentation. The natural map $`B_n^+B_n`$ is injective. .
2. There is a *fundamental braid* $`𝐃`$. In the old presentation, $`𝐃`$ has length $`((n^2n)/2)1`$ and is the half-twist
$$\mathrm{\Delta }=(\sigma _1\mathrm{}\sigma _{n1})(\sigma _1\mathrm{}\sigma _{n2})\mathrm{}(\sigma _1\sigma _2)\sigma _1.$$
In the new presentation it has length $`n1`$ and it is the $`(1/n)`$-twist
$$\delta =a_{n(n1)}a_{(n1)(n2)}\mathrm{}a_{32}a_{21}.$$
The fundamental braid admits many many braid transformations, in both the old and the new presentations, and so can be written in many ways as a positive word in the braid generators. As a result of this flexibility, it has two important properties:
1. For any generator $`a`$, there exist $`A,BB_n^+`$ such that:
$`𝐃=aA=Ba`$;
2. For each generator $`a`$ we have $`a𝐃=𝐃\tau (a)`$ and also $`𝐃a=\tau ^1(a)𝐃`$, where $`\tau `$ is the automorphism of $`B_n`$ which is defined by $`\tau (\sigma _i)=\sigma _{ni}`$ for the old presentation and $`\tau (a_{ts})=a_{(t+1)(s+1)}`$ for the new presentation.
3. $`\tau (\{𝐃\})=\{𝐃\}`$.
3. There are partial orderings ‘$``$’ and ‘$``$’ in $`B_n`$. For two words $`V`$ and $`W`$ in $`B_n`$ we say that $`VW`$ (resp. $`WV`$) if $`V=PW`$ (resp. $`V=WP`$) for some $`PB_n^+`$. Note that $`W`$ is a positive word if and only if $`We`$. We denote $`V<W`$ (resp. $`V>W`$) if $`VW`$ (resp. $`VW`$) and $`VW`$. In general $`VW`$ is not equivalent to $`WV`$, although if either $`W`$ or $`V`$ is a power of $`𝐃`$ the two ordering conditions are equivalent because powers of $`𝐃`$ commute with elements of $`B_n`$ up to powers of the index-shift automorphism $`\tau `$. Note that $`\tau `$ preserves the partial ordering.
4. The symbol $`𝐐`$ denotes the set of all initial subwords of $`𝐃`$, and $`𝐐^{}=𝐐\{e,𝐃\}.`$ The cardinality $`|𝐐_{old}|`$ is $`n!`$, whereas the cardinality $`|𝐐_{new}|`$ is the $`n^{th}`$ Catalan number. Note that $`|\delta |<|\mathrm{\Delta }|`$, also $`|𝐐_{new}|<|𝐐_{old}|`$. These are the main reasons why it is sometimes easier to work with the new presentation than the old.
5. The geodesic length $`l_Q(\{W\})`$ was introduced and investigated by Ruth Charney in . It is the smallest integer $`k`$ such that there is a word $`q_1q_2\mathrm{}q_k`$ representing $`\{W\}`$, with each $`q_i𝐐𝐐^1`$. Define the geodesic length of the conjugacy class $`l_Q([W])`$ to be the shortest such representation for words in the conjugacy class $`[W]`$.
6. For each positive word $`P`$, there is a decomposition, called the *left-greedy decomposition*, $`P=A_0P_0`$ for $`A_0𝐐`$ and $`P_0e`$, where $`A_0`$ has maximal length among all such decompositions, i.e. if $`P=A_0^{}P_0^{}`$, where $`A_0^{}𝐐`$ and $`P_0^{}B_n^+`$, then $`A_0^{}A_0`$. The term ‘greedy’ suggests that $`A_0`$ has absorbed as many letters from $`P_0`$ as it can without leaving $`𝐐`$. The canonical factor $`A_0`$ is called the *maximal head* of $`P`$. If $`P=A_0P_0=A_0^{}P_0^{}`$ in left greedy form, then $`\{A_0\}=\{A_0^{}\}`$ and $`\{P_0\}=\{P_0^{}\}`$. (Remark: The term left-canonical decomposition was used in and , however in recent years left-greedy decomposition has become the term of choice for the same concept in the literature, hence we now change our notation.)
7. Any word $`W`$ in the generators admits a unique normal form which solves the word problem in $`B_n`$. The normal form is:
$$W=𝐃^uA_1A_2\mathrm{}A_k,u\text{ZZ },A_i𝐐^{},$$
where for each $`1ik1`$, the product $`A_iA_{i+1}`$ is a left-greedy decomposition. The integer $`u`$ (resp. $`u+k`$) is called the *infimum* of $`W`$ (resp.supremum of $`W`$) and denoted by $`\text{inf}(W)`$ (resp. $`\text{sup}(W)`$).
8. To solve the conjugacy problem, we need to study the maximum and minimum values of inf and sup for the conjugacy class rather than for the word class. We consider the following two operations $`𝐜`$ and $`𝐝`$, called *cycling* and *decycling*, respectively. For a given braid in normal form $`W=𝐃^uA_1A_2\mathrm{}A_k`$, we define:
$`𝐜(W)`$ $`=`$ $`𝐃^uA_2A_3\mathrm{}A_k\tau ^u(A_1)`$
$`𝐝(W)`$ $`=`$ $`𝐃^u\tau ^u(A_k)A_1\mathrm{}A_{k1}.`$
In general the braids on the right hand side will not be in normal form, and must be rearranged into normal form before the operation can be repeated.
9. Theorem (see ):
1. If $`W`$ is conjugate to $`V`$ and if $`\text{inf}(V)>\text{inf}(W)`$, then repeated cycling will produce $`𝐜^{\mathrm{}}(W)`$ with $`\text{inf}(𝐜^{\mathrm{}}(W))>\text{inf}W`$.
2. If $`W`$ is conjugate to $`V`$ with $`\text{sup}(V)<\text{sup}(W)`$, then repeated decycling will produce $`𝐝^{\mathrm{}}(W)`$ with $`\text{sup}(𝐝^{\mathrm{}}(W))<\text{sup}(W)`$.
3. The maximum value of inf and the minimum value of sup can be achieved simultaneously.
10. The *super summit set* SSS($`[W]`$) () is the set of all conjugates of $`W`$ which have the maximal infimum and the minimal supremum in the conjugacy class $`[W]`$. It is a proper subset of the summit set SS$`([W])`$ which was introduced in by Garside in .
11. Theorem (see ): Let $`W,W^{}`$ be any two words in SSS($`[W]`$ = SSS$`([W^{}])`$. Then there is a sequence
$$W=W_0W_1\mathrm{}W_k=W^{}$$
such that each intermediate braid $`W_i\mathrm{SSS}([W])`$ and each $`W_{i+1}`$ is a conjugate of $`W_i`$ by a single member of $`𝐐`$.
12. By the theorems in $`\mathrm{\S }`$2.9 and $`\mathrm{\S }`$2.11 one can compute SSS($`[W]`$) as follows:
* Obtain an element $`W^{}`$ in the super summit set by iterating cyclings and decyclings, starting with any given word $`W`$.
* Compute the whole super summit set from $`W^{}`$ as follows: Compute $`AW^{}A^1`$ for all $`A𝐐`$ and collect the braids in the super summit set. Repeat the same process with each newly obtained element, until no new elements are obtained.
Therefore there is a finite time algorithm to generate SSS$`(W)`$. This algorithm solves the conjugacy problem in $`B_n`$. The integers $`\text{inf}([W])`$ and $`\text{sup}([W])`$ are the same for all members of SSS$`(W)`$ and so are partial invariants of the conjugacy class $`[W]`$.
In this article, we obtain an upper bound for the necessary number of cyclings and decyclings in the theorem in $`\mathrm{\S }`$2.9 above. for both the old presentation and the new presentation. We denote the word length of $`W`$ by $`|W|`$. Our main result is:
###### Theorem 1.
Let $`WB_n`$. If $`\text{inf}(W)`$ is not maximal for $`[W]`$, then inf$`(𝐜^{|𝐃|1}(W))>\text{inf}(W)`$. If $`\text{sup}(W)`$ is not minimal for $`[W]`$, then sup$`(𝐝^{|𝐃|1}(W))<\text{sup}(W)`$.
As immediate applications, we have:
###### Corollary 2.
Given any braid word $`WB_n`$, there is an algorithm which is polynomial in both word length and braid index for the computation of $`\text{inf}[W]`$ and $`\text{sup}[W]`$. Using the new presentation the complexity of the algorithm is $`O(|W|^2n^2)`$.
###### Corollary 3.
There is an algorithm which is polynomial in both word length and braid index for the computation of the geodesic length $`l_Q([W])`$ of the conjugacy class of $`W`$, using either presentation. Using the new presentation the complexity is $`O(|W|^2n^2)`$.
## 3. Proof of Theorem 1 and Corollaries 2 and 3.
Proof of Theorem 1: We focus on cycling because the proof and the difficulties are essentially identical for decycling.
Here is the plan of the proof. We begin with a word $`W=𝐃^uP`$ which is in normal form, so that $`u=\text{inf}(W)`$ and $`P>e`$. By hypothesis $`\text{inf}([W])>u`$, so there exists an integer $`m`$ such that $`u=\text{inf}(W)=\text{inf}(𝐜(W))=\mathrm{}=\text{inf}(𝐜^m(W))`$, but $`\text{inf}(𝐜^{m+1}(W))>u`$. Each instance of cycling can be realized by conjugation of $`W`$ by an element in $`𝐐^{}`$, so we know there are $`A_1^{},A_2^{}\mathrm{},A_m^{}𝐐^{}(𝐐^{})^1`$ such that after conjugating $`W`$, successively, by $`A_1^{},A_2^{}\mathrm{},A_m^{}`$ we obtain $`W^{}=R^{}𝐃^uP(R^{})^1`$ with $`\text{inf}(W^{})=u+1`$. (See Lemma 4 below.) Write $`R^{}`$ in normal form. (See Lemmas 8 and 9.) Our plan is to show that the sequence of lengths of the canonical factors $`H_m^{},\mathrm{},H_0^{}`$ for $`R^{}`$ satisfies $`|H_m^{}|<|H_{m1}^{}|<\mathrm{}<|H_0^{}|`$. Since each $`H_i^{}𝐐^{}`$, we have $`e<|H_i^{}|<|𝐃|`$. This places a limit on the length of the chain, i.e. $`m+1|𝐃|1`$ or $`m|𝐃|2`$, as claimed.
We used the symbols $`R^{},A_i^{},H_j^{}`$ in the description above, but in the actual proof we will use symbols $`R,A_i,H_j`$ which differ a little bit from $`R^{},A_i^{},H_j^{}`$ because we wish to focus on the changes in the positive part $`P`$ of $`W`$, rather than on changes in $`𝐃^uP`$:
###### Lemma 4.
Choose any $`WB_n`$. Let $`W=𝐃^uP`$, where $`u=\text{inf}(W)`$. Then $`\text{inf}([W])>\text{inf}(W)`$ if and only if there exists a positive word $`R`$ such that $`RP\tau ^u(R^1)𝐃`$.
###### Proof.
By hypothesis $`\text{inf}([W])>\text{inf}(W)`$, so there exists $`XB_n`$ with $`\text{inf}(XWX^1)>\text{inf}(W).`$ Let $`X=𝐃^vY,Ye`$, where $`v=\text{inf}(Y)`$. Then:
$$(𝐃^vY)(𝐃^uP)(Y^1𝐃^v)𝐃^{u+1},$$
which implies (via part (ii) of (3) above) that:
$$(\tau ^u(Y))(P)(Y^1)𝐃.$$
Set $`R=\tau ^u(Y)`$, so that $`Y=\tau ^u(R)`$. Then $`Re`$ and
$$(R)(P)(\tau ^u(R^1))𝐃,$$
as claimed. ∎
We will need to understand the structure of the positive word $`R`$ in Lemma 4, and to learn how the normal form of $`R`$ is related to that of $`W`$ and its images under repeated cycling. Once we understand all these issues, we will be able to extract information from $`R`$ about repeated cycling. We begin our work with several preparatory lemmas (i.e. Lemmas 5, 6 and 7.):
###### Lemma 5.
Suppose that $`Pe`$ and that $`RP𝐃`$ for some $`Re`$. Let $`P=A_0P_0`$ be in left-greedy form. Then $`RA_0𝐃`$.
###### Proof.
See Proposition 3.9 (IV) of for the new presentation and Proposition 2.10 of for the old presentation. ∎
###### Lemma 6.
If $`WB_n`$ and $`A𝐐`$, then $`\text{inf}(W)\text{inf}(WA)\text{inf}(W)+1.`$
###### Proof.
Since $`WWAW𝐃`$ and $`\text{inf}(W𝐃)=\text{inf}(W)+1`$ the assertion follows. ∎
For each $`A𝐐`$, let $`\overline{A}`$ denote the unique member of $`𝐐`$ which satisfies $`\overline{A}A=𝐃`$.
###### Lemma 7.
Let $`Z=B_{\mathrm{}}B_\mathrm{}1\mathrm{}B_1`$ be the normal form for $`Ze`$. Then $`Z^1𝐃^{\mathrm{}}e`$ and the normal form for $`Z^1𝐃^{\mathrm{}}`$ is
$$\tau (\overline{B}_1)\tau ^2(\overline{B}_2)\mathrm{}\tau ^{\mathrm{}}(\overline{B}_{\mathrm{}}).$$
###### Proof.
Observe that $`𝐃=\overline{B}_iB_i`$ implies that $`𝐃=\tau ^i(\overline{B}_i)\tau ^i(B_i)`$ for every $`i=1,\mathrm{},\mathrm{}`$. Therefore:
$`Z^1𝐃^{\mathrm{}}`$ $`=`$ $`B_1^1B_2^1\mathrm{}B_{\mathrm{}}^1𝐃^{\mathrm{}}`$
$`=`$ $`(𝐃\tau (B_1^1))(𝐃\tau ^2(B_2^1))\mathrm{}(𝐃\tau ^{\mathrm{}}(B_{\mathrm{}}^1))`$
$`=`$ $`(\tau (\overline{B}_1)\tau (B_1)\tau (B_1^1))(\tau ^2(\overline{B}_2)\tau ^2(B_2)\tau ^2(B_2^1))\mathrm{}(\tau ^{\mathrm{}}(\overline{B}_{\mathrm{}})\tau ^{\mathrm{}}(B_{\mathrm{}})\tau ^{\mathrm{}}(B_{\mathrm{}}^1))`$
$`=`$ $`\tau (\overline{B}_1)\tau ^2(\overline{B}_2)\mathrm{}\tau ^{\mathrm{}}(\overline{B}_{\mathrm{}})`$
because
$$\tau (B_i)\tau (B_i^1)=\tau (B_iB_i^1)=\tau (e)=e.$$
We continue the proof of Theorem 1. By $`\mathrm{\S }`$2.8 and Lemma 6 there exists a nonnegative integer $`m`$ such that $`\text{inf}(𝐜(W))=\text{inf}(𝐜^2(W))=\mathrm{}=\text{inf}(𝐜^m(W))`$ but $`\text{inf}(𝐜^{m+1}(W))=\text{inf}(W)+1`$. To prove Theorem 1, we must show that $`m+1`$ is bounded above by $`|𝐃|1`$. Let $`W=𝐃^uP`$, where $`Pe`$ and $`u=\text{inf}(W)`$. By Lemma 4 we know that $`\text{inf}([W])>\text{inf}(W)`$ if and only if there exists a positive word $`R`$ such that $`RP\tau ^u(R^1)𝐃`$. Assume that among all such words we have chosen $`R`$ so that $`|R|`$ is minimal. We wish to describe this shortest word $`R`$ as a specific product (in general not left-greedy) of canonical factors. Our first observation is:
###### Lemma 8.
$`\text{inf}(R)=0`$.
###### Proof.
If not, then $`R=𝐃R^{}`$ for some $`R^{}e`$ and
$$R^{}P\tau ^u(R^{})^1=𝐃^1RP\tau ^u(R)^1𝐃=\tau (RP\tau ^u(R)^1)𝐃,$$
which contradicts the minimality of $`|R|`$. ∎
###### Lemma 9.
Let $`𝐜^i(W)=𝐃^uA_iP_i`$, where $`A_i`$ is the maximal head of $`𝐜^i(W)`$. Then the positive word $`R`$ whose existence is guaranteed by Lemma 4 is related to the $`A_i`$’s as follows:
$$R=\tau ^m(\overline{A}_m)\mathrm{}\tau ^1(\overline{A}_1)\overline{A}_0.$$
###### Proof.
Our starting point is:
$$RP\tau ^u(R)^1𝐃,$$
which implies that $`RP𝐃`$. Since $`P=A_0P_0`$ is left-greedy, Lemma 5 then implies that $`RA_0𝐃`$ and so $`R=R_1\overline{A}_0`$ for some positive word $`R_1`$. Now:
$`RA_0P_0\tau ^u(R)^1`$ $`=`$ $`R_1\overline{A}_0A_0P_0\tau ^u(\overline{A}_0^1)\tau ^u(R_1^1)`$
$`=`$ $`R_1𝐃P_0\tau ^u(A_0)𝐃^1\tau ^u(R_1^1)`$
$`=`$ $`R_1\tau ^1\left(A_1P_1\right)\tau ^u(R_1^1).`$
Since $`RA_0P_0\tau ^u(R^1)𝐃`$, we conclude that:
$$R_1\tau ^1(A_1P_1)\tau ^u(R_1^1)𝐃.$$
Iterating the construction, we obtain $`R_1=R_2\tau ^1(\overline{A}_1)`$ for some positive word $`R_2`$, also $`R_2=R_3\tau ^2(\overline{A}_2),\mathrm{},R_m=R_{m+1}\tau ^m(\overline{A}_m)`$. Putting all of these together we learn that:
$$R=R_{m+1}\tau ^m(\overline{A}_m)\mathrm{}\tau ^1(\overline{A}_1)\overline{A}_0$$
for some positive word $`R_{m+1}`$. Let $`S=\tau ^m(\overline{A}_m)\mathrm{}\tau ^1(\overline{A}_1)\overline{A}_0,`$ so that $`R=R_{m+1}S.`$ A straightforward calculation shows that
$$\tau ^{(m+1)}(SP\tau ^u(S^1))=A_{m+1}P_{m+1}.$$
Since $`\text{inf}(𝐜^{m+1}(W))=\text{inf}(W)+1`$, we have
$$1=\text{inf}(\tau ^{(m+1)}(SP\tau ^u(S^1)))=\text{inf}(SP\tau ^u(S^1)).$$
By the minimality of $`|R|`$, we must have $`R=S`$. Lemma 9 is proved. ∎
The expression given for $`R`$ in the statement of Lemma 9 is in general not in normal form. We now study the maximal head $`H_0`$ of $`R=\tau ^m(\overline{A}_m)\mathrm{}\tau ^1(\overline{A}_1)\overline{A}_0`$, and related canonical factors $`H_1,\mathrm{},H_m`$. To define them, let $`H_k`$ be the maximal head of $`\tau ^m(\overline{A}_m)\mathrm{}\tau ^k(\overline{A}_k)R.`$
###### Lemma 10.
$`e<H_m<H_{m1}<\mathrm{}<H_1<H_0<𝐃.`$
###### Proof.
Our first observation is that $`\text{inf}(R)=0`$ (see Lemma 8). Since $`H_0`$ is the maximal head of $`R`$, it follows that
$$H_0<𝐃.$$
Our second observation is that by hypothesis $`\text{inf}(𝐜^m(W))=u`$ and $`𝐜^m(W)=𝐃^uA_mP_m`$ is left-greedy, so that $`A_m<𝐃`$, which implies that $`e<\overline{A}_m`$ and so:
$$e<H_m.$$
Our third observation is that by the definition of $`H_k`$ we must have:
$$H_mH_{m1}\mathrm{}H_1H_0.$$
Therefore the only thing that we need to prove is that $`H_{k+1}H_k`$ for $`k=0,1,\mathrm{},m1`$.
We first prove the assertion for $`k=0`$. Assume that $`H_1=H_0`$. We will show that this leads to a contradiction to our choice of $`R`$.
We are given that:
$$R=\tau ^m(\overline{A}_m)\mathrm{}\tau ^1(\overline{A}_1)\overline{A}_0=B_{\mathrm{}}B_\mathrm{}1\mathrm{}B_1,$$
where the decomposition on the left comes from Lemma 9 and the one on the right is the normal form for $`R`$. By Lemma 7 the normal form for $`R^1𝐃^{\mathrm{}}`$ is $`\tau (\overline{B}_1)\tau ^2(\overline{B}_2)\mathrm{}\tau ^{\mathrm{}}(\overline{B}_{\mathrm{}}).`$
By hypothesis $`H_1=H_0=B_{\mathrm{}}`$, so
$$\tau ^m(\overline{A}_m)\mathrm{}\tau ^1(\overline{A}_1)=B_{\mathrm{}}R_1$$
for some $`R_1e`$. Since $`B_{\mathrm{}}R_1\overline{A}_0=R=B_{\mathrm{}}B_\mathrm{}1\mathrm{}B_1`$, it follows that $`B_\mathrm{}1\mathrm{}B_1=R_1\overline{A}_0`$ and so
$$B_\mathrm{}1\mathrm{}B_1P𝐃.$$
Let $`a_i`$ be the infimum of $`B_\mathrm{}1\mathrm{}B_1P\tau ^u\left(\tau (\overline{B}_1)\tau ^2(\overline{B}_2)\mathrm{}\tau ^i(\overline{B}_i)\right)`$. Then
1. $`a_01`$ by the above discussion,
2. $`a_ia_{i+1}a_i+1`$ by Lemma 5, and
3. if $`a_i=a_{i+1}`$, then $`a_i=a_{i+1}=\mathrm{}=a_{\mathrm{}}`$ since $`\tau ^i(\overline{B}_i)\tau ^{i+1}(\overline{B}_{i+1})`$ is left-greedy.
If $`a_\mathrm{}1\mathrm{}`$, then
$$(B_\mathrm{}1\mathrm{}B_1)P\tau ^u\left(B_\mathrm{}1\mathrm{}B_1\right)^1𝐃,$$
which contradicts the minimality of $`|R|`$. So $`a_\mathrm{}1\mathrm{}1`$. Then $`a_i=a_{i+1}`$ for some $`i\mathrm{}2`$ and so $`a_i=a_{i+1}=\mathrm{}=a_{\mathrm{}}\mathrm{}1`$ so that $`\text{inf}(RP\tau ^u(R^1))0`$. However, by our choice of $`R`$, we know that $`\text{inf}(RP\tau ^u(R^1))>0`$. Retracing our steps we conclude that the assumption $`H_1=H_0`$ is impossible, so $`H_1<H_0<𝐃`$.
It remains to attack the cases $`k>0`$. The method is identical to the case $`k=0`$. Set $`V=𝐜^k(W)=𝐃^uA_kP_k`$ and let $`V`$ play the role of $`W=𝐃^uA_0P_0`$. ∎
The proof of Theorem 1 is almost complete. We have learned that $`e<H_m<H_{m1}<\mathrm{}<H_1<H_0<𝐃`$. This implies that:
$$0<|H_m|<|H_{m1}|<\mathrm{}<|H_1|<|H_0|<|𝐃|.$$
Thus the length $`m+1`$ of the chain must be smaller than $`|𝐃|`$, that is $`m+1|𝐃|1`$. The proof of Theorem 1 is complete. ∎
Proof of Corollary 2: The proof follows directly from Theorem 1 and the estimates in . In Theorem 4.4 of it is shown that for the new presentation there is an algorithm rewriting a word into its left greedy form that is a $`O(|W|^2n)`$ solution to the word problem. The initial preparation of our algorithm puts a given word $`W`$ into its left greedy form and takes $`O(|W|^2n)`$. Notice that the number of factors is proportional to $`|W|`$ in the worst case. In order to compute inf we need to cycle at most $`n2`$ times. After each cycling the new word so-obtained must be put into left greedy form but this time it takes only $`O(|W|n)`$ by Corollary 3.14 of . Thus the test to determine whether inf is maximal takes $`O(|W|n^2)`$. If it is not the entire process must be repeated, but the number of such repeats, i.e., the total increase of inf, is clearly bounded by the number of factors so the entire calculation is $`O(|W|^2n^2)`$. We note that if $`W`$ is a positive word, the total increase of inf is the maximum number of powers of $`\delta `$ formed cyclically from $`W`$ but this number is clearly bounded by $`|W|/(n1)`$, so the entire calculation is $`O(|W|^2n)`$. The discussion for the old presentation is similar and is left to the reader. ∎
Proof of Corollary 3: Let $`W=𝐃^uA_1A_2\mathrm{}A_k`$ be a word which is in normal form and which realizes the maximum value $`u`$ of inf and the minimum value $`k`$ of sup for the word class $`\{W\}`$. The geodesic length $`l_Q(\{W\})`$ of $`\{W\}`$ is computed in (or see ) as follows:
* If $`u0`$ then $`W`$ is a positive word of geodesic length $`l_Q(\{W\})=u+k`$.
* If $`ku<0`$, then we may use the fact that for every $`X_i𝐐`$ there exists $`Y_i𝐐`$ with $`X_iY_i=𝐃`$. From this it follows that $`𝐃^1X_i=Y_i^1`$. Using the additional fact that if $`\tau `$ is the index shift automorphism of $`\mathrm{\S }`$2.2, then $`\tau (𝐐)=𝐐`$, it follows that we may eliminate all of the powers of $`𝐃`$ and replace $`u`$ of the factors $`A_1,A_2,\mathrm{}A_u𝐐`$ with appropriate elements of $`𝐐^1`$, thereby achieving a shorter word. So in this case $`l_Q(\{W\})=k`$.
* If $`u<k`$ then every factor $`A_1,A_2,\mathrm{}A_u𝐐`$ is replaced by an appropriate element of $`𝐐^1`$. After all of these reductions the new word will be entirely negative. Its geodesic length is $`l_Q(\{W\})=u`$.
* The three cases may be combined into a single formula:
$`l_Q(\{W\})=max(k+u,u,k).`$
The above considerations relate to the length of a word class $`\{W\}`$. However, observe that the normal form for elements in the conjugacy class $`[W]`$ is identical to that for the word class, moreover if $`Y,Z`$ are in the super summit set of $`[W]`$ then $`\text{inf}(Y)=\text{inf}(Z)`$ and $`\text{sup}(Y)=\text{sup}(Z)`$. Since the complexity of computing $`l_Q([W])`$ is identical to the complexity of computing $`\text{inf}([W])`$ and $`\text{sup}([W])`$, the assertion then follows from Corollary 2. ∎
## 4. Are the cycling-decycling bounds sharp?
Note that the bound we obtained for the number of cyclings and decyclings in Theorem 1 is $`n2`$ for the new presentation and $`1+(n1)(n2)/2`$ for the old presentation. In this section we investigate whether these bounds are sharp.
We first give an example of $`n`$-braid written in the new generators for which $`n2`$ cyclings are required to increase the infimum. This shows that the bound given in Theorem 1 is sharp for the new presentation. To simplify notation, use $`[t,t1,\mathrm{},s]`$ instead of $`a_{t(t1)}a_{(t1)(t2)}\mathrm{}a_{(s+1)s}.`$ Consider the example $`W=([2,1][5,4,3])([3,2])`$ in normal form. Then
$`𝐜(W)`$ $`=`$ $`([3,2])([2,1][5,4,3])=([3,2,1][5,4])([4,3])`$
$`𝐜^2(W)`$ $`=`$ $`([4,3])([3,2,1][5,4])=([4,3,2,1])([5,4])`$
$`𝐜^3(W)`$ $`=`$ $`([5,4])([4,3,2,1])=[5,4,3,2,1]=\delta `$
So $`\text{inf}(W)=\text{inf}(𝐜(W))=\text{inf}(𝐜^2(W))=0`$ but $`\text{inf}(𝐜^3(W))=1`$. More generally, if
$$W=[2,1][n,n1,\mathrm{},3,2]=([2,1][n,n1,\mathrm{},3])([3,2]),$$
then $`\text{inf}(W)=\text{inf}(𝐜^{n3}(W))=0`$ but $`\text{inf}(𝐜^{n2}(W))=1`$. See Figure 1(a) for a sketch of the braid $`W`$ in the case $`n=7`$.
In the old presentation, the example in shows that $`\text{inf}(W)=\text{inf}(𝐜(W))=0`$ but $`\text{inf}(𝐜^2(W))=1`$. There are plenty of examples for which more than 2 cyclings are required to increase the infimum. Let $`(a_1,\mathrm{},a_n)`$ denote the permutation braid corresponding to the permutation $`\pi `$ on $`\{1,\mathrm{},n\}`$ defined by $`\pi (i)=a_i`$, Consider the following example, with $`WB_{2k+1}`$.
$$W=(\underset{2k1}{\underset{}{2k+1,2k,\mathrm{},3}},1,2)(\underset{k}{\underset{}{1,2,\mathrm{},k}},\underset{k}{\underset{}{k+2,\mathrm{},2k+1}},k+1)$$
Then $`\text{inf}(W)=\text{inf}(𝐜^{2k1}(W))=0`$ but $`\text{inf}(𝐜^{2k}(W))=1`$. See Figure 1(b) for a sketch of this example in the case $`n=7`$.
For another example let $`WB_{2k+1}`$ be such that
$`W`$ $`=`$ $`(\underset{k1}{\underset{}{2k+1,\mathrm{},k+3}},k+1,k+2,\underset{k}{\underset{}{k,k1,\mathrm{},1}})`$
$`(\underset{k1}{\underset{}{3,4,\mathrm{},k+1}},1,\underset{k1}{\underset{}{k+2,\mathrm{},2k}},2,2k+1)`$
See Figure 1(c). Then $`4k5`$ cyclings are needed to increase the infimum. So if $`n`$ is odd, there is an example for which $`2n7`$ cyclings are needed. Therefore the lower bound for the old presentation is at least linear in $`n`$.
We do not know an exact bound that works for every $`n`$-braid written in the old generators. It is easy to see that the upper bound in $`B_3`$ is 1, that is, if $`\text{inf}(W)=\text{inf}(𝐜(W))`$ for $`WB_3`$, then the infimum is already maximized. In an exhaustive search, we learned that the upper bound for $`B_4`$, using the old presentation, is 2 for positive words whose normal form contains up to 5 canonical factors.
## 5. Complexity issues and the conjugacy problem
In this section we consider implications of the work in the preceding sections for the complexity of the conjugacy problem in $`B_n`$.
### 5.1. The special cases $`n=3`$ and $`4`$:
Before discussing the problem, it will be helpful to review what is known about the cases $`n=3`$ and $`4`$, since well-chosen examples always help one to arrive at a better understanding of a problem. In the manuscript P.J. Xu introduced the new presentation for $`B_3`$ and used it to solve the word and conjugacy problems in $`B_3`$ and to study the letter lengths of shortest words in a word and conjugacy class in $`B_3`$, using the new presentation. Her main result in this regard was that words of shortest ‘geodesic length’ (she doesn’t use the term geodesic length, which was introduced after she completed her work) are, without further work, also words of shortest letter length in the new generators. She also found growth functions for $`B_3`$, both for word classes and conjugacy classes, proved that they were rational, and computed the rational functions which described them. Her algorithm for the conjugacy problem was clearly polynomial in $`|W|`$.
In the word and conjugacy problems were solved in $`B_4`$, using the new presentation and following the methods of . The authors also solved the shortest word problem in conjugacy classes. In a forthcoming paper the second and third author of this paper will prove that the algorithm for the conjugacy problem in is polynomial in $`|W|`$.
### 5.2. The conjugacy problem:
In $`\mathrm{\S }`$2.10 above the super summit set SSS$`([W])`$ of the conjugacy class of $`WB_n`$ is defined. It is a finite set and it can be computed in a systematic manner in a finite number of steps from any braid word $`W`$ which realizes $`\text{inf}([W])`$ and $`\text{sup}([W])`$. The Theorem which is quoted in $`\mathrm{\S }`$2.11 above asserts that $`W`$ is conjugate to $`V`$ in $`B_n`$ if and only if $`\text{inf}([W])=\text{inf}([V]),\text{sup}([W])=\text{sup}([V])`$ and SSS$`([W])=\mathrm{SSS}([V])`$.
The super summit set has a fairly transparent structure when $`n=3`$, the main reason being that words in $`𝐐^{}=𝐐\{\delta ,e\}`$ all have length 1. In $`B_4`$ the situation is a little bit more complicated, but still within reach. Let $`W=\delta ^uA_1A_2\mathrm{}A_k`$ be in the super summit set of $`[W]`$ and be in normal form, and let $`A=A_1A_2\mathrm{}A_k`$ be the ‘positive part’ of $`W`$. Notice that $`\delta `$ has letter length 3 and the $`A_j^{}s`$ are elements of $`𝐐^{}`$, and so have letter length 1 or 2. Let $`k_1`$ (resp. $`k_2`$) be the number of factors in $`A`$ which have length 1 (resp. 2). Let $`e`$ be the exponent sum of $`W`$. Clearly $`e`$ is a class invariant. Since $`k=k_1+k_2`$ and since $`e=3u+k_1+2k_2`$ it follows that $`k_1`$ and $`k_2`$ are determined by the triplet $`(u=\text{inf},k=\text{sup},e)`$. This makes the SSS somewhat easier to understand in the case $`n=4`$ than in the general case.
In the general case the super summit set SSS$`([W])`$ splits into orbits under cycling and decycling. Clearly the number of such orbits and their sizes are class invariants, but unfortunately we have examples to show that they are not complete invariants. The orbits are complicated by the fact that $`𝐐^{}`$ contains elements of letter length $`1,2,\mathrm{},n2`$, and the number $`k_i`$ of elements of letter length $`i`$ of a member of SSS$`([W])`$ is no longer controlled by $`(u,k,e)`$. We don’t know whether $`k_1,\mathrm{},k_{n2}`$ are orbit invariants, and if they can vary from one orbit to another. Also, while it is known that one can pass from any orbit to any other orbit by conjugating by an appropriate product of elements of $`𝐐`$, it is difficult to understand which products do the job. While the super summit set is a great improvement over the summit set of , it is still too big to make it possible to study many examples. For all these reasons the complexity of the conjugacy problem remains open at this time. Nevertheless, based on what we know, we conjecture:
###### Conjecture 11.
There is an algorithmic solution to the conjugacy problem in $`B_n`$, using the combinatorial approach which is described in this paper, which is polynomial in word length $`|W|`$ for each fixed braid index $`n`$.
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# Magnetic and electrical resistance behaviour of the oxides, Ca3-xYxLiRuO6 (x= 0.0, 0.5 and 1.0)
## I Introduction
There is a considerable interest in the current literature in identifying Ru-based magnetic and/or superconducting oxides. The magnetism in a Ru-based material for the first time was reported nearly forty years ago, namely, in SrRuO<sub>3</sub> exhibiting ferromagnetism below (T<sub>C</sub>=) 160 K . Subsequently there has been very little work in this direction until the discovery of superconductivity at (T<sub>c</sub>=) 1.5 K in Sr<sub>2</sub>RuO<sub>4</sub>, which turned out to be the first superconductor not containing Cu with the same structure as by-now-well-known K<sub>2</sub>NiF<sub>4</sub>-related high-T<sub>c</sub> oxides; this work naturally rekindled the interest on Ru-based oxides. The investigations in recent years led to the discovery of the following Ru based materials with the magnetism arising from Ru: Sr<sub>4</sub>Ru<sub>3</sub>O<sub>10</sub>, T<sub>C</sub>= 148 K ; Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub>, T<sub>C</sub>= 104 K , Eu(Gd)Sr<sub>2</sub>RuCu<sub>2</sub>O<sub>8</sub>, T<sub>N</sub>= 168 (185)K , YSr<sub>2</sub>RuO<sub>6</sub>, T<sub>N</sub>= 26 K , and A<sub>3</sub>A’RuO<sub>6</sub> (A= Ca, Sr; A’= Li, Na), T<sub>N</sub>= 70-120 K . While the ferromagnets of this class are generally metallic, the antiferromagnets are found to be insulators. However, there is no distinct evidence for high-T<sub>c</sub> superconductivity (from Ru d band) in any Ru-based oxides. At this point, it is worth noting that superconductivity could be induced in the temperature (T) range 30-50 K by small Cu doping in the antiferromagnetic insulator, YSr<sub>2</sub>RuO<sub>6</sub>, and in fact superconductivity and magnetism seem to compete in such substituted oxides ; however, considering that Cu doping is essential for superconductivity, one is not sure whether this phenomenon arises from Ru. Also, the compounds, R<sub>1.4</sub>Ce<sub>0.6</sub>RuSr<sub>2</sub>Cu<sub>2</sub>O<sub>10</sub> (R= Eu and Gd) are found to be superconducting below about 40 K arising from Cu layers, whereas Ru exhibits magnetic ordering below about 180 K . Thus, one may conclude that there is no clearcut evidence for high-T<sub>c</sub> in Ru oxides not containing Cu and that there is an upsurge of activity in identifying novel Ru based oxides. In this context, we report here the results of magnetization (M) and electrical resistivity ($`\rho `$) on a compound, Ca<sub>3</sub>LiRuO<sub>4</sub> , hitherto not extensively studied in the literature; we have also probed the influence of partial Y substitution for Ca on its properties in order to explore whether metallicity (and possibly superconductivity) could be induced by electron doping, considering that such doping effects have profound influnce on the properties of many other oxides, e.g., by now well-known, cuprates and manganates.
The compound under investigation has been synthesized and reported to adopt K<sub>4</sub>CdCl<sub>6</sub>-type rhombohedral structure . The crystallographic details can be found in Ref. 7. The structure consists of infinite chains of alternating face-sharing LiO<sub>6</sub> trigonal prisms and RuO<sub>6</sub> octahedra (antitrigonal prisms), which are separated by Ca ions.
## II Experimental
The compounds have been prepared by solid state method as described by Darriet et al . To start with, CaRuO<sub>3</sub> was prepared by heating in air of appropriate amounts of high purity CaCO<sub>3</sub> (Leico Industries, 99.995%) and RuO<sub>2</sub> (Cerac, 99.9%) at 750 <sup>o</sup>C for 24 hours and then for 8 days at 1100 <sup>o</sup>C in air. The samples were then prepared by reaction of CaRuO<sub>3</sub>, Li<sub>2</sub>CO<sub>3</sub> (Koch Chemicals, 99.999 %) and Y<sub>2</sub>O<sub>3</sub> (Johnson-Matthey, 99.9%) at 550 <sup>o</sup>C for 24 hours, at 800 <sup>o</sup>C for 24 hours and at 950 <sup>o</sup>C for two weeks with intermediate grindings. All the preparations were carried out under a flow of oxygen. The samples were subsequently characterized by x-ray diffraction (Cu K<sub>α</sub>) and scanning electron microscope. The magnetization measurements were performed by a commercial (Quantum Design) Superconducting Quantum Interference Device (SQUID) Magnetometer in the T interval 5 - 300 K. We have obtained the magnetization data for all the compositions in different ways (isothermal magnetization (M), T-dependence of magnetic susceptibility ($`\chi `$) at different magnetic fields (H), field-cooled (FC) and zero-field-cooled (ZFC) $`\chi `$ behaviour) and we show only those data which are relevant to highlight the main findings. The temperature (T= 77-300 K)) dependent electrical resistivity ($`\rho `$) behaviour was probed by a conventional four-probe method, employing silver paint to make electrical contacts.
## III Results and Discussion
The results of x-ray diffraction measurements are shown in Fig. 1. The diffraction pattern for the parent compound could be indexed on the basis of K<sub>4</sub>CdCl<sub>6</sub>-type structure except perhaps for the presence of few very weak unidentified lines ($`<`$5%); the diffraction pattern seen by us is identical to the one reported by Darriet et al and the lattice constants (a= 9.221 Åand c= 10.798 Å) are also in agreement with their report. A piece of the parent compound was subjected to the 950 <sup>o</sup>C (24 hrs) heat treatment in air and brought to room temperature by quenching and this heat treatment apparently resulted in purer sample, as indicated by the disappearance of the weak extra line (marked by asterisk in Fig. 1). The Y substituted samples also are characterized by the same x-ray diffraction pattern as that of that of the parent compound, though few, weak additional lines start appearing in the range 2$`\mathrm{\Theta }`$ = 30 to 40 degrees, the origin of which is at present unclear. The x-dependence of unit-cell volume clearly reveals that there is a contraction of the unit-cell (along c-direction, see Fig. 1) with Y substitution, expected for replacement of (bigger) Ca ions. The homogeniety of the samples were further checked by scanning electron spectroscopy and we do not find evidence for segregation of any other phase. The energy dispersive x-ray analysis in addition established that the proportion of the metallic elements are in good agreement with the starting compositions and uniform throughout the sample. We therefore conclude that the properties reported here are characteristic of the pure phases.
In figure 2a we show the T dependent $`\chi `$ behavior recorded in the presence of a field of 2 kOe for the ZFC as well as FC state of a specimen of the parent compound. It is distinctly clear that there is a sharp rise at 113 K as T is lowered with a peak at about 107 K for the ZFC state, however without any such peak for FC state; thus ZFC-FC curves deviate below 113 K in perfect agreement with the experimental observations of Darriet et al . This observation establishes that we have prepared a sample with the same magnetic characteristics as those of these authors. Thus, these data provide evidence for the fact that there is a magnetic ordering setting in below (T<sub>N</sub>=) 113 K in this compound. In order to see how T<sub>N</sub> is influenced by Y substitution, we have recorded the FC data very carefully in a field of 100 Oe and the data (see Fig. 2b) suggest that the value of T<sub>N</sub> undergoes only a marginal reduction with increasing x (108 and 106 K for x= 0.5 and 1.0 respectively). With respect to $`\chi `$ behaviour in the paramagnetic state, the plot of inverse $`\chi `$ versus T is found to be linear in the range 160-300 K (see Fig. 3) and the value of the effective moment ($`\mu `$<sub>eff</sub>) and paramagnetic Curie temperature ($`\theta `$<sub>p</sub>) obtained from the region 150-300 K are nearly independent of x (3.96, 3.65 and 3.77 $`\mu `$<sub>B</sub> and -250, -180 and -207 K for x= 0.0, 0.5 and 1.0 respectively). It may be noted that the value of $`\mu `$<sub>eff</sub> is indicative of pentavalent state of Ru ion, assuming spin-only contribution to magnetic moment. It is important to note that the sign of $`\theta `$<sub>p</sub> is negative, which suggests that the exchange interaction is of an antiferromagnetic type. However, the value of T<sub>N</sub> is far below the magnitude of $`\theta `$<sub>p</sub>, which may indicate complex nature of the magnetism of this compound. In order to get better insight on the nature of magnetic ordering, we have also performed isothermal M measurements at various temperatures for the ZFC state of the samples at many temperatures and typical behaviour are shown in Fig. 4. It is obvious that, at 120 K, M varies linearly with H as expected for a paramagnetic state; for T$`<`$100 K, M is a non-linear function of H for initial applications of H (below about 20 kOe), however undergoing linear variation for higher values of H, without any indication for saturation. This behaviour of M establishes that these compounds are better classified as antiferromagnets, though we observe some hysteresis loops as shown in the inset of Fig. 4, for instance for x= 0.0 and 0.5.
Thus, all these results establish that the magnetism of Ru ions is fairly insensitive to Y substitution at the Ca site. Thus magnetic Ru chains are decoupled by intervening Ca ions. This finding may support the idea of quasi-one-dimensional behaviour of magnetism in this compound, a question of debate in this class of compounds .
We now present the results of $`\rho `$ measurements (Fig. 5). It is obvious that the Y substituted compounds remain insulating like the parent compound, exhibiting activated behaviour with the activation energy marginally decreasing by about 10 meV by Y substitution (for both x= 0.5 and 1.0) from 90 meV for x= 0.0. Interestingly, the $`\rho `$ of quenched specimen of the parent compound is significantly lower and the activation energy is also considerably reduced to 1.5 meV (Fig. 6).
## IV conclusions
To conclude, the magnetic and electrical transport behaviour of the oxides Ca<sub>3-x</sub>Y<sub>x</sub>LiRuO<sub>6</sub> (x= 0.0, 0.5 and 1.0) have been investigated. These are one of the very few Ru-based oxides exhibiting magnetic ordering at rather high temperatures, presumably of a complex type. The carrier doping induced by Y substitution for Ca apparently does not bring about any significant change in the magnetic and transport behaviour of the parent compound. On the basis of the present results, we infer that the Ru-magnetism could be of quasi-one-dimensional character. Finally, we would like to point out an observation in the low-field magnetization data for the ZFC-state of the specimen (see Fig. 4): In the magnetically ordered state, below about 400 Oe, the sign of M is negative (which is however positive for the FC specimen). Though a possible source of this negative M could be that the specimen may not be in true ZFC state due to a small negative remanent field in the SQUID magnetometer we employed, this view can be challenged by the observation that the data at 120 K (in the paramagnetic state) do not show negative M at low fields. Alternatively, the magnetic moments (below T<sub>N</sub>) in the ZFC state may have got locked in the direction of the negative remanence field requiring larger applications of H in the positive direction to reorient them. More work is required to understand this aspect of these compounds.
## V References
1. J.T. Randall and R. Ward, J. Am. Chem. Soc., 81 (1959) 2629.
2. Y. Maeno, H. Hashimoto, I.K. Yoshida, S. Ishizaki, T. Fujita, J.G. Bednorz, and F. Lichtenberg, Nature (London), 372 (1994) 532.
3. G. Cao, S.K. McCall, J.E. Crow, and R.P. Guertin, Phys. Rev. B, 56 (1997) R5740.
4. G. Cao, S. McCall, and J.E. Crow, Phys. Rev. B, 55 (1997) R672.
5. I. Felner, U. Asaf, S. Reich, and Y. Tsabba, Physica C,311 (1999) 163.
6. I. Felner, U. Asaf, Y. Levi, and O. Millo, Phys. Rev. B, 55 (1997) R3374.
7. M.F. Wu, D.Y. Chen, F.Z. Chien, S.R. Sheen, D.C. Ling, C.Y. Tai, G.Y. Tseng, D.H. Chen, and F.C. Zhang, Z. Phys. B, 102 (1997) 37.
8. J. Darriet, F. Grasset, and P.D. Battle, Mater. Res. Bull., 32 (1997) 139.
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# Clustering statistics on a light-cone in the cosmological redshift space
## Abstract
We summarize a series of our recent work concerning the cosmological redshift-space distortion and light-cone effects. After briefly describing the theoretical formalism, we show how those effects are sensitive to the cosmological parameters. Then we apply this formalism to predict the two-point correlation functions and power spectra for the X-ray clusters, galaxies and quasars in future surveys, and discuss their cosmological implications.
## 1. Cosmological effects in the high-z universe
Redshift surveys of galaxies definitely serve as the central database for observational cosmology. In addition to the existing catalogues including CfA1, CfA2, SSRS, and the Las Campanas survey, upcoming galaxy surveys such as 2dF (2-degree Field Survey) and SDSS (Sloan Digital Sky Survey) are expected to provide important clues to the evolution of structure in the universe. In addition to those shallower surveys, clustering in the universe in the range $`z=13`$ has been partially revealed by, for instance, the Lyman-break galaxies and X-ray selected AGNs. In particular, the 2dF and SDSS QSO redshift surveys promise to extend the observable scale of the universe by an order of magnitude, up to a few Gpc. A proper interpretation of such redshift surveys in terms of the clustering evolution, however, requires an understanding of many cosmological effects which can be neglected for $`z1`$ and thus have not been considered seriously so far. These cosmological contaminations include linear redshift-space (velocity) distortion (Kaiser 1987), nonlinear redshift-space (velocity) distortion (e.g., Suto & Suginohara 1991; Cole, Fisher, & Weinberg 1994), cosmological redshift-space (geometrical) distortion (Alcock & Paczyński 1979; Ballinger, Peacock, & Heavens 1996; Matsubara & Suto 1996), and cosmological light-cone effect (Yamamoto & Suto 1999; Suto et al. 1999; Yamamoto, Nishioka & Suto 1999).
We describe a theoretical formalism to incorporate those effects, in particular the cosmological redshift-distortion and light-cone effects, and present several specific predictions in cold dark matter (CDM) models.
## 2. Cosmological redshift-space distortion
Due to a general-relativistic effect through the geometry of the universe, the observable separations perpendicular and parallel to the line-of-sight direction, $`x_\mathrm{s}=(c/H_0)z\delta \theta `$ and $`x_\mathrm{s}=(c/H_0)\delta z`$, are mapped differently to the corresponding comoving separations in real space $`x_{}`$ and $`x_{}`$:
$`x_\mathrm{s}(z)`$ $`=`$ $`x_{}cz/[H_0(1+z)d_\mathrm{A}(z)]x_{}/c_{}(z),`$ (1)
$`x_\mathrm{s}(z)`$ $`=`$ $`x_{}H(z)/H_0x_{}/c_{}(z),`$ (2)
with $`d_\mathrm{A}(z)`$ being the angular diameter distance. The difference between $`c_{}(z)`$ and $`c_{}(z)`$ generates an apparent anisotropy in the clustering statistics, which should be isotropic in the comoving space. Then the power spectrum in cosmological redshift space, $`P^{(\mathrm{CRD})}`$, is related to $`P^{(\mathrm{S})}`$ defined in the comoving redshift space as
$$P^{(\mathrm{CRD})}(k_\mathrm{s},k_\mathrm{s};z)=\frac{1}{c_{}(z)^2c_{}(z)}P^{(\mathrm{S})}(\frac{k_\mathrm{s}}{c_{}(z)},\frac{k_\mathrm{s}}{c_{}(z)};z),$$
(3)
where the first factor comes from the Jacobian of the volume element $`dk_\mathrm{s}^2dk_\mathrm{s}`$, and $`k_\mathrm{s}=c_{}(z)k_{}`$ and $`k_\mathrm{s}=c_{}(z)k_{}`$ are the wavenumber perpendicular and parallel to the line-of-sight direction. If one assumes a scale-independent deterministic linear bias, the power spectrum distorted by the peculiar velocity field, $`P^{(\mathrm{S})}(k;z)`$, is known to be well approximated by the following expression (Cole et al. 1995; Peacock & Dodds 1996):
$$P^{(\mathrm{S})}(k_{},k_{};z)=b^2(z)P_{\mathrm{mass}}^{(\mathrm{R})}(k;z)\left[1+\beta (z)\left(\frac{k_{}}{k}\right)^2\right]^2D\left[k_{}\sigma _\mathrm{P}(z)\right],$$
(4)
where $`k_{}`$ and $`k_{}`$ are the comoving wavenumber perpendicular and parallel to the line-of-sight of an observer, and $`P_{\mathrm{mass}}^{(\mathrm{R})}(k;z)`$ is the mass power spectrum in real space. The finger-of-god effect is modeled by the damping function, $`D\left[k_{}\sigma _\mathrm{P}(z)\right]`$, for which we assume a Lorentzian. Then equation (3) reduces to
$`P^{(\mathrm{CRD})}(k_\mathrm{s},\mu _k;z)`$ $`=`$ $`{\displaystyle \frac{b^2(z)}{c_{}(z)^2c_{}(z)}}P_{\mathrm{mass}}^{(\mathrm{R})}({\displaystyle \frac{k_\mathrm{s}}{c_{}(z)}}\sqrt{1+\left[{\displaystyle \frac{1}{\eta (z)^2}}1\right]\mu _k^2};z)`$ (6)
$`\times \left\{1+\left[{\displaystyle \frac{1}{\eta (z)^2}}1\right]\mu _k^2\right\}^2\left\{1+\left[{\displaystyle \frac{1+\beta (z)}{\eta (z)^2}}1\right]\mu _k^2\right\}^2\left[1+{\displaystyle \frac{k_\mathrm{s}^2\mu _k^2\sigma _{\mathrm{P}}^{}{}_{}{}^{2}}{2c_{}^2(z)}}\right]^1,`$
where we introduce
$`k_\mathrm{s}\sqrt{k_\mathrm{s}^2+k_\mathrm{s}^2},\mu _kk_\mathrm{s}/k_\mathrm{s},\eta c_{}/c_{},`$ (7)
following Ballinger et al. (1996) and Magira et al. (2000).
Figure 1. shows anisotropic power spectra $`P^{(\mathrm{CRD})}(k_\mathrm{s},\mu _k;z=2.2)`$. As specific examples, we consider SCDM (standard CDM), LCDM (Lambda CDM), and OCDM (Open CDM) models, which have $`(\mathrm{\Omega }_0,\lambda _0,h,\sigma _8)`$ $`=(1.0,0.0,0.5,0.6)`$, $`(0.3,0.7,0.7,1.0)`$, and $`(0.3,0.0,0.7,1.0)`$, respectively. These sets of cosmological parameters are chosen so as to reproduce the observed cluster abundance (Kitayama & Suto 1997). Our theoretical predictions use the fitting formulae of Peacock & Dodds (1996; PD) for the nonlinear power spectrum, $`P_{\mathrm{mass}}^{(\mathrm{R})}(k;z)=2\pi ^2\mathrm{\Delta }_{\mathrm{NL}}^2(k,z)/k^3`$, of Mo, Jing, & Börner (1997) for the pair-wise peculiar velocity dispersions:
$`\sigma _{\mathrm{P},\mathrm{MJB}}^2`$ $``$ $`\mathrm{\Omega }(z)H_0^2\left[1{\displaystyle \frac{1+z}{D_+^2(z)}}{\displaystyle _z^{\mathrm{}}}{\displaystyle \frac{D_+^2(z^{})}{(1+z^{})^2}}𝑑z^{}\right]{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dk}{k}}{\displaystyle \frac{\mathrm{\Delta }_{\mathrm{NL}}^2(k,z)}{k^2}},`$ (8)
with $`D_+(z)`$ being the linear growth rate. Clearly the linear theory predictions ($`\sigma _\mathrm{P}=0`$; top panels) are quite different from the results of N-body simulations (bottom panels), indicating the importance of the nonlinear velocity effects ($`\sigma _\mathrm{P}=\sigma _{\mathrm{P},\mathrm{MJB}}`$; middle panels).
Next we decompose the power spectrum into harmonics:
$`P(k,\mu _k;z)={\displaystyle \underset{l:\mathrm{even}}{}}P_l(k)L_l(\mu _k),P_l(k;z){\displaystyle \frac{2l+1}{2}}{\displaystyle _1^1}𝑑\mu _kP(k,\mu _k;z)L_l(\mu _k),`$ (9)
where $`L_l(\mu _k)`$ are the $`l`$-th order Legendre polynomials. Similarly, the two-point correlation function is decomposed as
$`\xi (x,\mu _x;z)={\displaystyle \underset{l:\mathrm{even}}{}}\xi _l(x)L_l(\mu _x),\xi _l(x;z){\displaystyle \frac{2l+1}{2}}{\displaystyle _1^1}𝑑\mu _x\xi (x,\mu _x;z)L_l(\mu _x),`$ (10)
using the direction cosine, $`\mu _x`$, between the separation vector and the line-of-sight. The above multipole moments satisfy the following relations:
$`\xi _l(x;z)={\displaystyle \frac{1}{2\pi ^2i^l}}{\displaystyle _0^{\mathrm{}}}P_l(k;z)j_l(kx)k^2𝑑k,`$ (11)
with $`j_l(kx)`$ being spherical Bessel functions. Substituting $`P^{(\mathrm{CRD})}(k_\mathrm{s},\mu _k;z)`$ in equation (9) yields $`P_l^{(\mathrm{CRD})}(k_\mathrm{s};z)`$, and then $`\xi ^{(\mathrm{CRD})}(𝐱_𝐬;z)`$ can be computed from equation (11).
Comparison of the monopoles and quadrupoles from simulations and model predictions exhibits how the results are sensitive to the cosmological parameters, which in turn may put potentially useful constraints on $`(\mathrm{\Omega }_0,\lambda _0)`$. Figure 2. indicates the feasibility, which interestingly results in a constraint fairly orthogonal to that from the Supernovae Ia Hubble diagram.
## 3. Cosmological light-cone effect
Observing a distant patch of the universe is equivalent to observing the past. Due to the finite light velocity, a line-of-sight direction of a redshift survey is along the time, as well as spatial, coordinate axis. Therefore the entire sample does not consist of objects on a constant-time hypersurface, but rather on a light-cone, i.e., a null hypersurface defined by observers at $`z=0`$. This implies that many properties of the objects change across the depth of the survey volume, including the mean density, the amplitude of spatial clustering of dark matter, the bias of luminous objects with respect to mass, and the intrinsic evolution of the absolute magnitude and spectral energy distribution. These aspects should be properly taken into account in order to extract cosmological information from observed samples of redshift surveys. We apply the formulation on the light-cone originally developed by Yamamoto & Suto (1999) to X-ray selected clusters and on-going SDSS galaxy and QSO catalogues.
### 3.1. Two-point correlation functions of X-ray selected clusters
Provided an X-ray flux-limited sample of clusters ($`S>S_{\mathrm{lim}}`$), it is fairly straightforward to compute its two-point correlation function $`\xi _{\mathrm{cl}}^\mathrm{S}(R,z;>S_{\mathrm{lim}})`$ at a given $`z`$; a fairly accurate empirical expression for the bias parameter $`b(z)`$ as a function of the halo mass is known (e.g., Jing 1998), and the mass is translated to the X-ray temperature assuming the virial equilibrium, and then to the X-ray luminosity from the observed luminosity-temperature relation (e.g., Kitayama & Suto 1996). The corresponding correlation function on the light-cone is given by
$`\xi _{\mathrm{X}\mathrm{cl}}^{\mathrm{LC}}(R;>S_{\mathrm{lim}})={\displaystyle \frac{{\displaystyle _{z_{\mathrm{min}}}^{z_{\mathrm{max}}}}dz{\displaystyle \frac{dV_\mathrm{c}}{dz}}n_0^2(z)\xi _{\mathrm{cl}}^\mathrm{S}(R,z(r);>S_{\mathrm{lim}})}{{\displaystyle _{z_{\mathrm{min}}}^{z_{\mathrm{max}}}}𝑑z{\displaystyle \frac{dV_\mathrm{c}}{dz}}n_0^2(z)}},`$ (12)
where $`R`$ is the comoving separation of a pair of clusters, $`z_{\mathrm{max}}`$ and $`z_{\mathrm{min}}`$ denote the redshift range of the survey, and $`dV_\mathrm{c}/dz`$ is the comoving volume element per unit solid angle (Suto et al. 2000; Moscardini et al. 2000). The comoving number density of clusters in the flux-limited survey, $`n_0(z;>S_{\mathrm{lim}})`$, is computed by integrating the Press – Schechter mass function.
Figure 3. plots several predictions for two-point correlation functions under different assumptions; linear and nonlinear mass correlations in real space at $`z=0`$ using the Bardeen et al. (1986; BBKS) and PD formulae for mass power spectra, cluster correlations with linear redshift-space distortion (Kaiser 1987) and with full redshift-space distortion at $`z=0`$ using $`\sigma _{\mathrm{P},\mathrm{MJB}}`$. These should be compared with our final predictions on the light-cone in redshift space (with $`S_{\mathrm{lim}}=10^{14}`$erg/s/cm<sup>2</sup>; thick solid lines). Figure 4. shows our predictions for $`\xi _{cl}^{\mathrm{LC}}(R)`$ for cluster samples selected with different flux-limit $`S_{\mathrm{lim}}`$ (left panels), and with additional temperature and absolute bolometric luminosity limits, $`T_{\mathrm{lim}}`$ and $`L_{\mathrm{lim}}`$ (middle and right panels). For the latter two cases, $`S_{\mathrm{lim}}=10^{14}`$erg/s/cm<sup>2</sup> is assumed for definiteness. The results are insensitive to $`S_{\mathrm{lim}}`$, but very sensitive to $`T_{\mathrm{lim}}`$ and $`L_{\mathrm{lim}}`$, reflecting the strong dependence of the bias on the latter quantities.
For a cosmological application of the present result, it is interesting to examine how $`r_{c,0}(S_{\mathrm{lim}})`$ defined through
$`\xi _{\mathrm{cl}}^{\mathrm{LC}}(r_{c0};>S_{\mathrm{lim}})=1`$ (13)
depends on $`\mathrm{\Omega }_0`$. This is summarized in Figure 5., where we fix the value of the fluctuation amplitude $`\sigma _8`$ adopting the cluster abundance constraint (Kitayama & Suto 1997). Again the results are not sensitive to the flux limit $`S_{\mathrm{lim}}`$. The dependence on $`\mathrm{\Omega }_0`$ is rather strong, and these predictions combined with the future observational results will be able to break the degeneracy of the cosmological parameters.
### 3.2. Power spectra of SDSS galaxy and QSO samples
Finally we present theoretical predictions of power spectra relevant for SDSS galaxy and QSO samples, fully taking account of the cosmological redshift-space distortion and light-cone effects. Denoting the comoving number density and the selection function of the objects by $`n_0^{\mathrm{com}}(z)`$, and $`\varphi (z)`$, Suto, Magira & Yamamoto (2000) obtain
$`P_l^{(\mathrm{LC},\mathrm{CRD})}(k_\mathrm{s})`$ $`=`$ $`{\displaystyle \frac{{\displaystyle _{z_{\mathrm{min}}}^{z_{\mathrm{max}}}}𝑑z{\displaystyle \frac{dV_\mathrm{c}}{dz}}[\varphi (z)n_0^{\mathrm{com}}(z)]^2c_{}(z)^2c_{}(z)P_l^{(\mathrm{CRD})}(k_\mathrm{s};z)}{{\displaystyle _{z_{\mathrm{min}}}^{z_{\mathrm{max}}}}𝑑z{\displaystyle \frac{dV_\mathrm{c}}{dz}}[\varphi (z)n_0^{\mathrm{com}}(z)]^2c_{}(z)^2c_{}(z)}}.`$ (14)
Figure 6. compares several predictions for the angle-averaged (monopole) power spectra normalized by the real-space counterpart in linear theory. The upper and lower panels adopt the selection functions appropriate for galaxies in $`0<z<z_{\mathrm{max}}=0.2`$ and QSOs in $`0<z<z_{\mathrm{max}}=5`$, respectively. The left and right panels present the results in SCDM and LCDM models. For simplicity we adopt a scale-independent linear bias model of Fry (1996),
$$b(z)=1+\frac{1}{D_+(z)}[b(k,z=0)1],$$
(15)
with $`b(k,z=0)=1`$ and $`1.5`$ for galaxies and quasars, respectively. It is clear that the cosmological redshift-space distortion and the light-cone effect substantially change the predicted shape and amplitude of the power spectra, even for the SDSS galaxy sample.
## 4. Summary and conclusions
We have presented a theoretical formalism to predict the two-point clustering statistics on a light-cone in the cosmological redshift space. The present methodology will find two completely different applications. For relatively shallower catalogues like galaxy samples, the evolution of bias is not supposed to be so strong. Thus, one may estimate the cosmological parameters from the observed degree of the redshift distortion, as has been conducted conventionally. Most importantly, one can now correct for the systematics due to the light-cone and geometrical distortion effects, which affect the estimate of the parameters by $`10`$%. Alternatively, for deeper catalogues like high-redshift quasar samples, one can extract information on the nonlinearity, scale-dependence and stochasticity of the object-dependent bias only by correcting the observed data on the basis of our formulae. In this case, although one should adopt a set of cosmological parameters a priori, those will be provided both from the low-redshift analysis described above and from precision data of the cosmic microwave background and supernovae Ia. In a sense, the former approach uses the light-cone and geometrical distortion effects as real cosmological signals, while the latter regards them as inevitable, but physically removable, noise. In both cases, the present methodology is essential in properly interpreting the observations of the universe at high redshifts.
I thank Y.P.Jing, Tetsu Kitayama, Hiromitsu Magira, Takahiko Matsubara, Hiroaki Nishioka, and Kazuhiro Yamamoto for enjoyable collaborations on which the present talk is based. Numerical computations were carried out on VPP300/16R and VX/4R at the Astronomical Data Analysis Center of the National Astronomical Observatory, Japan, as well as at RESCEU and KEK (National Laboratory for High Energy Physics, Japan). This research was supported in part by the Grants-in-Aid by the Ministry of Education, Science, Sports and Culture of Japan to RESCEU (07CE2002).
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# SU(2) WZW D-branes and their noncommutative geometry from DBI action Work supported in part by Polish State Committee for Scientific Research (KBN) under contract 2P 03 B03 715 (1998-2000) and the Alexander-von-Humboldt Foundation.
## Abstract:
Using properties of the DBI action we find D-branes on $`S^3`$ of the radius $`Q_5`$ corresponding to the conjugacy classes of SU(2). The branes are stable due to nonzero 2-form NSNS background. In the limit of large $`Q_5`$ the dynamics of branes is governed by the non-commutative Yang-Mills theory. The results partially overlap with those obtained in the recent paper hep-th/0003037.
Recently it has been discovered that in a special limit the dynamics of the matrix model is described by the non-commutative Yang-Mills theory . This has sparked new interest in NCG for strings propagating in NSNS antisymmetric tensor field background . Introduction of D-branes has provided deeper understanding of the role of non-commutativity and it has allowed to derive conditions under which NCG starts to play dominant role in the dynamics of strings. It has also led to new understanding of the connection between quantum groups and WZW models . These two papers have used the standard CFT language what is a drastic bound on possible applications. In particular it does not allow to analyze RR backgrounds so much studied in the context of Maldacena’s conjecture .
The purpose of this paper is to provide understanding of the results of in the more universal language then that of WZW models. The hope is that after taking this lesson one would be able to derive interesting results for more general string/M-theory backgrounds. Thus we shall describe various branes on the background of SU(2) WZW model using the D-brane effective action (DBI action) only. We shall also show how the non-commutativity appears in this approach. Methods applied here are limited to the case of large level of the SU(2) WZW model what in gravity language means large radius of the $`S^3`$.
Let us recall some of the results of and . D-branes in the level $`k`$ SU(2) WZW model are in one-to-one correspondence with special integer conjugacy classes $`ghg^1`$ for some fixed $`h`$ . There are $`k+1`$ of them: two D-particles ($`h=\pm 1`$) and $`k1`$ D2-branes corresponding to two-spheres. The n-th sphere passes through the point $`\mathrm{exp}(i\pi n\sigma ^3/k)`$SU(2), $`n=1\mathrm{}k1`$. We must also stress that D3-branes and D1-branes are excluded from this list. For large $`k`$ the 2-spheres are in fact so-called fuzzy spheres .
The example of string theory background which involves the level $`k`$ SU(2) WZW model is the near horizon limit of the F1, NS5 system (see e.q. ). Below we write only the relevant terms
$`ds^2/\alpha ^{}`$ $`=`$ $`Q_5d\mathrm{\Omega }_3^2`$
$`H^{NSNS}/\alpha ^{}`$ $`=`$ $`2Q_5ϵ_3`$
$`e^{2\varphi }`$ $`=`$ $`const.`$ (1)
where $`Q_5`$ denotes the number of NS5-branes and it is equal to the level $`k`$ of the SU(2) WZW model, $`ϵ_3`$ is the volume element of the unit 3-sphere.
The effective action of the D-branes is given by DBI expression
$$S_{DBI}=T_p_{\mathrm{Vol}}e^\varphi \sqrt{\mathrm{det}[(X^{}G+2\pi \alpha ^{}F+X^{}B)_{ab}]}$$
(2)
In the following we shall discuss classical configurations of branes embedded in $`S^3`$ of (1). Before we start to analyze equation of motion resulting from (2) we state several assumptions we make which seems to be natural here. We require the string coupling constant to be small and $`Q_5`$ to be large in which case (1) is the part of an exact string background as at this limit supergravity is the perfect description of string theory. Moreover D-branes can be described completely classically by the DBI action. We also assume that the higher order correction to the DBI action are negligible. We shall be interested here only in the $`S^3`$ part of the configuration thus it is even irrelevant if we consider IIB (as above) or IIA string. Thus some of the arguments given in our paper could be easily generalized to branes embedded in a 10d manifold of the form $`S^3\times M^7`$ under the condition that the embeddings are of the product structure i.e. the induced metric, pull-back of $`B`$ and $`F`$ fields are of block diagonal form.
Recall that D-branes are defined to be the ends of the open strings. The string couple to the external sources (gauge $`A`$ and $`B`$ field) as follows
$$\mathrm{exp}[\frac{i}{2\pi \alpha ^{}}(_\mathrm{\Sigma }2\pi \alpha ^{}X^{}A+_\mathrm{\Sigma }X^{}B)]$$
(3)
The example of the WZW model shows that the above formula can not be well defined globally for topologically non-trivial $`A`$ and $`B`$ fields. It is known that for closed strings $`\mathrm{\Sigma }=0`$ the proper formula is $`\mathrm{exp}[\frac{i}{2\pi \alpha ^{}}_{\widehat{\mathrm{\Sigma }}}\widehat{X}^{}H)]`$ for $`H`$ being locally $`dB`$ and $`\widehat{X}`$ is an extension of $`X`$ to a 3-manifold $`\widehat{\mathrm{\Sigma }}`$ such that $`\widehat{\mathrm{\Sigma }}=\mathrm{\Sigma }`$. Now we consider configuration of D-brane embedded into submanifold $`M_D`$ of the target space manifold $`M`$. One must repeat the above construction for the open string case . For the world sheet with one boundary the appropriate 3-manifold must respect $`\widehat{\mathrm{\Sigma }}=\mathrm{\Sigma }+D^2`$. Rewriting the WZW model with boundary we get the proper global form of (3) <sup>1</sup><sup>1</sup>1$``$ in front of the first term is due to different orientation of boundary $`\mathrm{\Sigma }=D^2`$ in $`D^2`$ compare to $`\mathrm{\Sigma }`$.
$$\mathrm{exp}\left[\frac{i}{2\pi \alpha ^{}}\left(_{D^2}\widehat{X}^{}(2\pi \alpha ^{}F+B)+_{\widehat{\mathrm{\Sigma }}}\widehat{X}^{}H\right)\right]$$
(4)
where $`\widehat{X}`$ is an extension of $`X(\mathrm{\Sigma })`$ to a full 2-disc $`D^2`$ such that $`\widehat{X}(D^2)M_D`$ ($`F0`$ only on the D-brane manifold $`M_D`$). We stress that (4) has proper gauge invariance and for topologically trivial $`H`$ it reduces to
$$\mathrm{exp}\left[\frac{i}{2\pi \alpha ^{}}\left(_{D^2}\widehat{X}^{}(2\pi \alpha ^{}F)+_\mathrm{\Sigma }X^{}B\right)\right]$$
(5)
i.e. to (3). Notice that one must be able to define $`B`$ on any $`\widehat{X}(D^2)`$ thus we must have $`[H]_{M_D}=0`$. The value of the integral (4) should not depend on the way one make the extension. This forces to put
$`{\displaystyle \frac{i}{2\pi \alpha ^{}}}\left[{\displaystyle _{C^2}}(2\pi \alpha ^{}F+B){\displaystyle \frac{i}{2\pi \alpha ^{}}}{\displaystyle _{C^3}}H\right]=2i\pi m`$ (6)
where $`C^2H_2(M_D)`$ and $`C^3H_3(M,M_D)`$. Here a note is necessary concerning topology of the problem. For the argument we need an exact sequence of homologies
$$\mathrm{}H_3(M_D)H_3(M)H_3(M,M_D)H_2(M_D)H_2(M)\mathrm{}$$
(7)
If we assume that $`H_3(M_D)=H_2(M)=0`$ then all cycles of $`H_3(M)`$ are in $`H_3(M,M_D)`$. Then one can write $`_{C^3}H=_{C^2}B`$ mod $`\frac{i}{2\pi }_{C_M^3}H=2\pi Q_5m`$ for $`C_M^3H_3(M)`$. Thus the quantization condition reads
$$_{C^2}F=2\pi n,n\text{ }\text{Z}$$
(8)
and $`n`$ is defined modulo $`Q_5`$. This is the same as postulated in . In the above we have disregarded the difference between the cycles in $`M`$ and $`M_D`$ and their image given by $`\widehat{X}`$.
D3 brane. Here we discuss the Dp-branes wrapped on the entire $`S^3`$. According to the condition $`[H]_{M_D}=0`$ we see that such a wrapping is impossible. We would like to provider here a different argument based on DBI action. First one must notice that due to $`[H]_{M_D}=0`$ the DBI action (2) is not well defined as $`B`$ is not well defined on $`S^3`$.
In order to be more specific we concentrate on D3 brane in the background (1) and change the brane description to the dual form of the DBI action discussed e.g. in . It has the same classical solutions as (2) what is the property we are interested in.
$$S_{DBI}_{\mathrm{Vol}}e^\varphi \sqrt{\mathrm{det}[(X^{}G+2\pi \alpha ^{}\stackrel{~}{F})_{ab}]}\pi \alpha ^{}\stackrel{~}{F}B$$
(9)
The last term come form CS part of the DBI action. Integrating it by parts we get
$$\pi \alpha ^{}\stackrel{~}{A}H^{NSNS}$$
(10)
thus the action contains only the well defined $`B`$ field strength. With the $`H^{NSNS}`$ background given by (1) we see that there is a U(1) charge generated on the D3 world-volume. The charge can not stay on $`S^3`$ as it is a compact space, thus it forces the brane to partially unwrap the sphere. In the case of $`AdS_3\times S^3`$the brane runs into the boundary of the AdS space. The above argument follows the baryon construction of .
D2 brane. Here we concentrate upon D2-brane case totally wrapped on $`S^3`$. It can be also a e.g. partially wrapped D3 brane. We analyze its equation of motion and find that contrary to the naive expectation the static brane it stable. As the indication of stability we invoke the lack of the tachyonic mode for the fluctuation of the brane.
In order to analyze the classical equations of motion we must find out the pull back of $`B`$ field to the brane world-volume. On any 2-d submanifold of $`S^3`$ the $`B`$ filed is well (but not uniquely) defined. We have the freedom of changing $`B`$ by an exact 2-form - in our case this will realized by choice of the solution for $`F_{12}`$. In the coordinates in which the metric on $`S^3`$ is $`ds^2=Q_5[d\varphi ^2+\mathrm{sin}^2(\varphi )d\mathrm{\Omega }_2^2]`$ we have for the chart covering $`\varphi =0`$
$$B=Q_5\alpha ^{}(\varphi \nu \frac{1}{2}\mathrm{sin}(2\varphi ))ϵ_2$$
(11)
where $`ϵ_2`$ is the volume form of the unit $`S^2`$.
We shall find extrema of the DBI action corresponding to branes wrapped on $`S^2`$ given by some constant angle $`\varphi `$. The Euler-Lagrange equations are respected by
$$2\pi \alpha ^{}F=Q_5\alpha ^{}(\varphi \nu )ϵ_2,\varphi (x)=\varphi =const.$$
(12)
with all the other components of $`F`$ equals zero. We also set $`\nu =0`$ requiring that the charge and the tension of the $`\varphi =0`$ brane be zero. It is worth to note that the classical solution exists for any angle $`\varphi `$. When we apply the quantization condition (8) we get
$$Q_5\varphi =2\pi n$$
(13)
We remind that $`n`$ is defined only modulo $`Q_5`$.
We can compare (13) with results one gets assuming that the brane couple to some RR fields i.e. carry RR charge
$$+T_pe^{2\pi \alpha ^{}F+X^{}B}\underset{q}{}C_q$$
(14)
where, $`T_{Dp}=1/((2\pi )^p\alpha ^{(p+1)/2}g_s)`$. The background $`2\pi \alpha ^{}F+X^{}B`$ generates RR charge of the D(p-2)-brane equals to
$$T_p_{S^2}(2\pi \alpha ^{}F+X^{}B)=T_p(4\pi \alpha ^{}Q_5)\frac{1}{2}\mathrm{sin}(2\varphi ).$$
One expects that this charge is integer multiple of $`T_{(p2)}`$ i.e.
$$\frac{1}{2\pi \alpha ^{}}_{S^2}(2\pi \alpha ^{}F+X^{}B)=2\pi n$$
(15)
but this is in contradiction with (13) for finite $`Q_5`$. If one takes the $`Q_5\mathrm{}`$ limit then both formulae agree. <sup>2</sup><sup>2</sup>2The gap between (13) and (15) has been filled recently in .
The second derivative of the DBI action with respect to the gauge fields and $`\varphi (x)`$ gives kinematics of fluctuations. One easily finds that fluctuations of $`\varphi (x)`$ only are massive but $`\varphi (x)`$ mixes with gauge field $`F`$ leading to some massless modes . Thus there is no tachyon in the spectrum and the brane configuration is stable.
Non-commutative geometry. We can also claim that at the $`Q_5\mathrm{}`$limit some of the branes are described by the non-commutative geometry. Here we follow the route of . First we notice that at the $`Q_5\mathrm{}`$limit we have
$`ds^2`$ $`=`$ $`\alpha ^{}{\displaystyle \frac{(n\pi )^2}{Q_5}}d\mathrm{\Omega }_2^20`$
$`2\pi \alpha ^{}F+X^{}B`$ $`=`$ $`\alpha ^{}(n\pi )ϵ_2`$ (16)
Thus the closed string metric goes to zero while induced $`2\pi \alpha ^{}F+X^{}B`$ is constant on the D2-branes world-volume what is a good sign for the non-commutativity. Next we calculate the open string metric and the Poisson structure inverting $`2\pi \alpha ^{}F+X^{}(B+g)`$. The inverse matrix is
$$\left(\frac{1}{X^{}(g+B)+2\pi \alpha ^{}F}\right)=\frac{1}{Q_5\mathrm{sin}\varphi \alpha ^{}}\left(\begin{array}{cc}\mathrm{sin}\varphi & \mathrm{cos}\varphi \\ \mathrm{cos}\varphi & \mathrm{sin}\varphi \end{array}\right)$$
(17)
Inverse of its symmetric part is the open string metric. We have $`G_{ab}=\alpha ^{}Q_5\delta _{ab}`$. Hence from the open string point of view all spheres have the same area! This, of course, is directly related to the flat direction $`\varphi =const`$ in the solution (12). The Poisson structure on $`S^2`$ (also called deformation parameter) is
$$\mathrm{\Theta }^{12}=\frac{2\pi }{Q_5}\mathrm{cot}\varphi \frac{2}{n}$$
(18)
The symplectic structure is the inverse of the Poisson structure and it is $`\omega _{12}=(n/2)`$. One can check that this parameter precisely corresponds to the symplectic structure used by Berezin in order to quantize $`S^2`$ . The non-commutative version on this $`S^2`$ is called the fuzzy spheres . From one may claim that the Y-M theory on this sphere is a theory of $`(n+1)\times (n+1)`$ hermitian matrices. Such a Y-M theory has $`(n+1)^2`$ degrees of freedom. Here we must stress that these results are in full agreement with . It would be interesting to make explicit comparison of the brane dynamics and the above matrix model.
We conclude that the branes dynamics is described by the non-commutative Y-M theory. The branes world-volume are 2-spheres which are non-commutative manifolds with the non-commutativity parameter $`\mathrm{\Theta }^{12}=\frac{2}{n}`$.
A note added. Some of the results of this paper have been independently obtained in the recent paper .
###### Acknowledgments.
I am grateful S.Theisen for illuminating discussions and reading the manuscript. I also thank K.Gawȩdzki, A.Alekseev, G. Arutyunov and A. Recknagel for comments and interest in this work.
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# Importance of Coupling between the Charge and Spin Degrees of Freedom in High 𝑇_𝑐 Superconductivity
## Abstract
Based on an improved SU(2) slave-boson approach showing coupling between the charge and spin degrees of freedom, we derive a phase diagram of high $`T_c`$ cuprates which displays both the superconducting and pseudogap phases in the plane of temperature vs. hole doping rate. It is shown that phase fluctuations in the order parameters results in a closer agreement with the observed phase diagram of an arch shape, by manifesting the presence of an optimal doping rate.
High $`T_c`$ superconductivity arises as a consequence of hole(or electron) doping in the parent cuprate oxides which are Mott insulators with antiferromagnetic long-range order. The observed phase diagram in the plane of temperature $`T`$ vs. hole doping rate $`\delta `$ shows the bose condensation(superconducting temperature) curve of an ‘arch’ shape rather than the often predicted linear increase, by manifesting the presence of the optimal doping rate of $`\delta =0.16`$ to $`0.2`$. On the other hand, the observed pseudogap temperature displays nearly a linear decrease with $`\delta `$. The high $`T_c`$ cuprate of $`Bi_2Sr_2CaCu_2O_{8+\delta }`$ with a higher pseudogap(spin gap) temperature $`T^{}`$ is observed to have a higher superconducting transition temperature $`T_c`$ than the cuprate of $`La_{2x}Sr_xCuO_4`$ with a lower $`T^{}`$. Further we find from the observed phase diagrams of both cuprates above that the two different high $`T_c`$ cuprates, $`La_{2x}Sr_xCuO_4`$ and $`Bi_2Sr_2CaCu_2O_{8+\delta }`$ display an universal behavior of $`T^{}/T_c`$ as a function of hole(positive charge) doping $`\delta /\delta _o`$ with $`\delta _o`$, the optimal doping rate, as is shown in Fig.1. Earlier Nakano et al found another type of universality by pointing out that the observed spin gap temperature $`T^{}`$ scaled by the maximum superconducting transition temperature $`T_c^{max}`$ at optimal doping, $`T^{}/T_c^{max}`$ as a function of doping scaled by the optimal doping, $`\delta /\delta _o`$ falls on the same curve for the two different high $`T_c`$ cuprates above. The two observations manifests the presence of a relationship between the spin gap (relevant to the spin degree of freedom) and the superconductivity (related to the charge degree of freedom). Thus, the spinon pairing(spin singlet pairing) for pseudogap phase and the charge pairing(holon pairing) for superconductivity are not independent owing to the manifest presence of coupling between the charge and spin degrees of freedom.
Various U(1) slave-boson approaches to the t-J Hamiltonian were able to predict such a linear decrease in the pseudogap temperature as a function of $`\delta `$-. In our earlier U(1) slave-boson study, we presented a phase diagram based on the allowance of holon pairing channel, thus showing the feature of the holon-pair bose condensation temperature rather than the single-holon bose condensation temperature-. On the other hand, all of these theories failed to predict the experimentally observed bose condensation temperature $`T_c`$ of the arch shape as a function of $`\delta `$. Instead a linear increase of $`T_c`$ with $`\delta `$ was predicted. Further the pseudogap phase was shown to disappear when the gauge fluctuations are introduced into the U(1) slave-boson mean field theory. Most recently Wen and Lee proposed an SU(2) theory to readily estimate the low energy phase fluctuations of order parameters and made a brief discussion on the possibility of holon(boson) pair condensation. In view of failure in the correct prediction of the bose condensation temperature $`T_c`$ in the phase diagram with earlier theories, in the present study we examine the variation of the holon-pair condensation temperature with the hole doping rate, by treating the phase fluctuations of the order parameters in the SU(2) slave boson theory. We realize from the aforementioned observation of the universality in $`T^{}/T_c`$ vs. $`\delta /\delta _o`$ for high $`T_c`$ cuprates that coupling between the spin(spinon) and charge(holon) degrees of freedom is essential for superconductivity. Our theoretical derivation from the use of the slave-boson theory for t-J Hamiltonian manifests this feature as is shown in the fourth term of the effective Hamiltonian in Eq.(32). In addition, comparison between the two approaches will be made to reveal the importance of the low energy phase fluctuations of the order parameters. The present work differs from our previous U(1) slave-boson study(of the phase diagram involving the holon-pair bose condensation) and other earlier studies-(involving the single holon condensation) in that coupling between the holon and spinon degrees of freedom in the slave-boson representation of the Heisenberg term of the t-J Hamiltonian is no longer neglected. We find from the treatment of the coupling that the predicted phase diagram displays the arch-shaped bose condensation curve(temperature $`T_c`$) as a function of hole doping rate in both treatments of the U(1) and SU(2) slave-boson approaches.
We write the t-J Hamiltonian,
$`H`$ $`=`$ $`t{\displaystyle \underset{<i,j>}{}}(c_{i\sigma }^{}c_{j\sigma }+c.c.)+J{\displaystyle \underset{<i,j>}{}}(𝐒_i𝐒_j{\displaystyle \frac{1}{4}}n_in_j).`$ (1)
Here $`𝐒_i`$ is the electron spin operator at site $`i`$, $`𝐒_i=\frac{1}{2}c_{i\alpha }^{}𝝈_{\alpha \beta }c_{i\beta }`$ with $`𝝈_{\alpha \beta }`$, the Pauli spin matrix element and $`n_i`$, the electron number operator at site $`i`$, $`n_i=c_{i\sigma }^{}c_{i\sigma }`$. We note that $`𝐒_i𝐒_j\frac{1}{4}n_in_j=\frac{1}{2}(c_{i2}^{}c_{j1}^{}c_{i1}^{}c_{j2}^{})(c_{j1}c_{i2}c_{j2}c_{i1})`$ leads to $`\frac{1}{2}b_ib_jb_j^{}b_i^{}(f_i^{}f_j^{}f_i^{}f_j^{})(f_jf_if_jf_i)`$ in the U(1) slave boson representation. It should be emphasized that the derivation of $`P(𝐒_i𝐒_j\frac{1}{4}n_in_j)P=\frac{1}{2}b_ib_jb_j^{}b_i^{}(f_i^{}f_j^{}f_i^{}f_j^{})(f_jf_if_jf_i)`$ is correct based on the U(1) slave-boson constraint and that most importantly physics demands the presence of this four boson operator $`b_ib_jb_j^{}b_i^{}`$ to represent the charge degree of freedom, in addition to the four fermion operator which represents the spin degree of freedom. Each one of $`n_in_j`$ and $`𝐒_i𝐒_j`$ in the Heisenberg term self-evidently reveals such physics(e.g., note that $`n_i`$ here represents the electron(or physically the negative charge) number operator). In earlier studies of the slave-boson theory, it is often assumed that $`b_ib_jb_j^{}b_i^{}=1`$. Strictly speaking, this is precise only at half-filling(or no hole doping). This is because charge fluctuations can not occur owing to the prohibition of actual electron hopping from site to site.
By admitting the coupling between the charge and spin degrees of freedom in the SU(2) slave-boson representation, the t-J Hamiltonian above can be written,
$`H`$ $`=`$ $`{\displaystyle \frac{t}{2}}{\displaystyle \underset{<i,j>\sigma }{}}[(f_{\sigma i}^{}f_{\sigma j})(b_{1j}^{}b_{1i}b_{2i}^{}b_{2j})`$ (5)
$`+(f_{\sigma j}^{}f_{\sigma i})(b_{1i}^{}b_{1j}b_{2j}^{}b_{2i})`$
$`+(f_{2i}f_{1j}f_{1i}f_{2j})(b_{1j}^{}b_{2i}+b_{1i}^{}b_{2j})`$
$`+(f_{1j}^{}f_{2i}^{}f_{2j}^{}f_{1i}^{})(b_{2i}^{}b_{1j}+b_{2j}^{}b_{1i})]`$
$``$ $`{\displaystyle \frac{J}{2}}{\displaystyle \underset{<i,j>}{}}(1h_i^{}h_i)(1h_j^{}h_j)\times `$ (10)
$`(f_{2i}^{}f_{1j}^{}f_{1i}^{}f_{2j}^{})(f_{1j}f_{2i}f_{2j}f_{1i})`$
$`\mu _0{\displaystyle \underset{i}{}}(h_i^{}h_i\delta ){\displaystyle \underset{i}{}}[i\lambda _i^{(1)}(f_{1i}^{}f_{2i}^{}+b_{1i}^{}b_{2i})`$
$`+i\lambda _i^{(2)}(f_{2i}f_{1i}+b_{2i}^{}b_{1i})`$
$`+i\lambda _i^{(3)}(f_{1i}^{}f_{1i}f_{2i}f_{2i}^{}+b_{1i}^{}b_{1i}b_{2i}^{}b_{2i})].`$
Here $`f_{\alpha i}`$ ( $`f_{\alpha i}^{}`$ ) is the spinon annihilation(creation) operator and $`h_i\left(\begin{array}{c}b_{1i}\\ b_{2i}\end{array}\right)`$ $`\left(h_i^{}=(b_{1i}^{},b_{2i}^{})\right)`$, the doublet of holon annihilation(creation) operators. $`\lambda _i^{(1),(2),(3)}`$ are the real Lagrangian multipliers to enforce the local single occupancy constraint in the SU(2) slave-boson representation.
The Heisenberg interaction term(the second term in Eq.(10)) above can be decomposed into terms involving mean fields and fluctuations respectively,
$`{\displaystyle \frac{J}{2}}(1h_i^{}h_i)(1h_j^{}h_j)(f_{2i}^{}f_{1j}^{}f_{1i}^{}f_{2j}^{})(f_{1j}f_{2i}f_{2j}f_{1i})`$ (11)
$`=`$ $`{\displaystyle \frac{J}{2}}<(1h_i^{}h_i)(1h_j^{}h_j)>\times `$ (20)
$`(f_{2i}^{}f_{1j}^{}f_{1i}^{}f_{2j}^{})(f_{1j}f_{2i}f_{2j}f_{1i})`$
$`{\displaystyle \frac{J}{2}}<(f_{2i}^{}f_{1j}^{}f_{1i}^{}f_{2j}^{})(f_{1j}f_{2i}f_{2j}f_{1i})>\times `$
$`(1h_i^{}h_i)(1h_j^{}h_j)`$
$`+{\displaystyle \frac{J}{2}}<(1h_i^{}h_i)(1h_j^{}h_j)>\times `$
$`<(f_{2i}^{}f_{1j}^{}f_{1i}^{}f_{2j}^{})(f_{1j}f_{2i}f_{2j}f_{1i})>`$
$`{\displaystyle \frac{J}{2}}((1h_i^{}h_i)(1h_j^{}h_j)<(1h_i^{}h_i)(1h_j^{}h_j)>)\times `$
$`((f_{2i}^{}f_{1j}^{}f_{1i}^{}f_{2j}^{})(f_{1j}f_{2i}f_{2j}f_{1i})`$
$`<(f_{2i}^{}f_{1j}^{}f_{1i}^{}f_{2j}^{})(f_{1j}f_{2i}f_{2j}f_{1i})>).`$
By introducing the Hubbard-Stratonovich fields, $`\rho _i^k`$, $`\chi _{ij}`$ and $`\mathrm{\Delta }_{ij}`$ in association with the direct, exchange and pairing channels of the spinon, we obtain the effective Hamiltonian from Eq.(10),
$`H_{eff}=`$ (22)
$`{\displaystyle \frac{J(1\delta )^2}{2}}{\displaystyle \underset{<i,j>}{}}{\displaystyle \underset{l=0}{\overset{3}{}}}\left((\rho _{ij}^l)^2\rho _{ij}^l(f_i^{}\sigma ^lf_i)\right)`$
$`+`$ $`{\displaystyle \frac{J(1\delta )^2}{4}}{\displaystyle \underset{<i,j>}{}}[|\chi _{ij}|^2\{f_{\sigma i}^{}f_{\sigma j}`$ (24)
$`+{\displaystyle \frac{2t}{J(1\delta )^2}}(b_{1i}^{}b_{1j}b_{2j}^{}b_{2i})\}\chi _{ij}c.c.]`$
$`+`$ $`{\displaystyle \frac{J(1\delta )^2}{2}}{\displaystyle \underset{<i,j>}{}}[|\mathrm{\Delta }_{ij}|^2\{(f_{2i}^{}f_{1j}^{}f_{1i}^{}f_{2j}^{})`$ (26)
$`{\displaystyle \frac{t}{J(1\delta )^2}}(b_{1j}^{}b_{2i}+b_{1i}^{}b_{2j})\}\mathrm{\Delta }_{ij}c.c.]`$
$``$ $`{\displaystyle \frac{J}{2}}{\displaystyle \underset{<i,j>}{}}|\mathrm{\Delta }_{ij}^f|^2\left[{\displaystyle \underset{\alpha ,\beta }{}}b_{\alpha i}^{}b_{\beta j}^{}b_{\beta j}b_{\alpha i}(h_j^{}h_j+h_i^{}h_i2\delta )\delta ^2\right]`$ (27)
$`+`$ $`{\displaystyle \frac{t^2}{J(1\delta )^2}}{\displaystyle \underset{<i,j>}{}}[(b_{1i}^{}b_{1j}b_{2j}^{}b_{2i})(b_{1j}^{}b_{1i}b_{2i}^{}b_{2j})`$ (29)
$`+{\displaystyle \frac{1}{2}}(b_{1j}^{}b_{2i}+b_{1i}^{}b_{2j})(b_{2i}^{}b_{1j}+b_{2j}^{}b_{1i})]`$
$`+`$ $`{\displaystyle \frac{J(1\delta )^2}{2}}{\displaystyle \underset{i,\sigma }{}}(f_{\sigma i}^{}f_{\sigma i})\mu _0{\displaystyle \underset{i}{}}(h_i^{}h_i\delta )`$ (30)
$``$ $`{\displaystyle \underset{i}{}}[i\lambda _i^1(f_{1i}^{}f_{2i}^{}+b_{1i}^{}b_{2i})+i\lambda _i^2(f_{2i}f_{1i}+b_{2i}^{}b_{1i})`$ (32)
$`+i\lambda _i^3(f_{1i}^{}f_{1i}f_{2i}f_{2i}^{}+b_{1i}^{}b_{1i}b_{2i}^{}b_{2i})],`$
where $`\mathrm{\Delta }_{ij}=<(f_{1i}f_{2j}f_{2i}f_{1j})\frac{t}{J(1\delta )^2}(b_{2i}^{}b_{1j}+b_{2j}^{}b_{1i})>=\mathrm{\Delta }_{ij}^f\frac{t}{J(1\delta )}\chi _{ij;12}^b`$, with $`\chi _{ij;12}^b=<b_{2i}^{}b_{1j}+b_{2j}^{}b_{1i}>`$ with $`\delta `$, hole doping rate. In Eq.(32) above we introduced $`<(f_{2i}^{}f_{1j}^{}f_{1i}^{}f_{2j}^{})(f_{1j}f_{2i}f_{2j}f_{1i})><(f_{2i}^{}f_{1j}^{}f_{1i}^{}f_{2j}^{})><(f_{1j}f_{2i}f_{2j}f_{1i})>=|\mathrm{\Delta }_{ij}^f|^2`$ and $`<(1h_i^{}h_i)(1h_j^{}h_j)><(1h_i^{}h_i)><(1h_j^{}h_j)>=(1\delta )^2`$ and neglected the last term in Eq.(20) above.
The four boson term in the fourth term of Eq.(32) allows holon pairing and a scalar boson field, $`\mathrm{\Delta }_{ij;\alpha \beta }^b`$ is introduced for the holon pairing between the nearest neighbor $`b_\alpha `$ and $`b_\beta `$single bosons with the boson index, $`\alpha ,\beta =`$ $`1`$ or $`2`$. Using the saddle point approximation, we obtain from Eq.(32) the mean field Hamiltonian,
$`H^{MF}={\displaystyle \frac{J(1\delta )^2}{2}}{\displaystyle \underset{<i,j>}{}}\left[|\mathrm{\Delta }_{ij}^f|^2+{\displaystyle \frac{1}{2}}|\chi _{ij}|^2+{\displaystyle \frac{1}{4}}\right]`$ (43)
$`+{\displaystyle \frac{J}{2}}{\displaystyle \underset{<i,j>}{}}|\mathrm{\Delta }_{ij}^f|^2\left[{\displaystyle \underset{\alpha ,\beta }{}}|\mathrm{\Delta }_{ij;\alpha \beta }^b|^2+\delta ^2\right]`$
$`{\displaystyle \frac{J(1\delta )^2}{2}}{\displaystyle \underset{<i,j>}{}}[\mathrm{\Delta }_{ij}^f(f_{1j}f_{2i}f_{2j}f_{1i})+c.c.]`$
$`{\displaystyle \frac{J(1\delta )^2}{4}}{\displaystyle \underset{<i,j>}{}}[\chi _{ij}(f_{\sigma i}^{}f_{\sigma j})+c.c.]+`$
$`{\displaystyle \frac{t}{2}}{\displaystyle \underset{<i,j>}{}}\left[\chi _{ij}(b_{1i}^{}b_{1j}b_{2j}^{}b_{2i})\mathrm{\Delta }_{ij}^f(b_{1j}^{}b_{2i}+b_{1i}^{}b_{2j})\right]c.c.`$
$`{\displaystyle \underset{<i,j>,\alpha ,\beta }{}}{\displaystyle \frac{J}{2}}|\mathrm{\Delta }_{ij}^f|^2[\mathrm{\Delta }_{ij;\alpha \beta }^b(b_{\alpha i}b_{\beta j})+c.c.]`$
$`{\displaystyle \underset{i}{}}[\mu _i(h_i^{}h_i\delta )+i\lambda _i^1(f_{1i}^{}f_{2i}^{}+b_{1i}^{}b_{2i}^{})`$
$`+i\lambda _i^2(f_{2i}f_{1i}+b_{2i}b_{1i})+i\lambda _i^3(f_{1i}^{}f_{1i}f_{2i}f_{2i}^{}+b_{1i}^{}b_{1i}+b_{2i}^{}b_{2i})]`$
$`{\displaystyle \frac{t}{2}}{\displaystyle \underset{<i,j>}{}}\left(\mathrm{\Delta }_{ij}^f(f_{1j}f_{2i}f_{2j}f_{1i})\right)\chi _{ij;12}^bc.c.`$
$`+{\displaystyle \frac{t^2}{2J(1\delta )^2}}{\displaystyle \underset{<i,j>}{}}\left|\chi _{ij;12}^b(b_{2i}^{}b_{1j}+b_{2j}^{}b_{1i})\right|^2`$
$`+{\displaystyle \frac{t^2}{J(1\delta )^2}}{\displaystyle \underset{<i,j>}{}}(b_{1i}^{}b_{1j}b_{2j}^{}b_{2i})(b_{1j}^{}b_{1i}b_{2i}^{}b_{2j}),`$
where $`\chi _{ij}=<f_{\sigma j}^{}f_{\sigma i}+\frac{2t}{J(1\delta )^2}(b_{1j}^{}b_{1i}b_{2i}^{}b_{2j})>`$, $`\mathrm{\Delta }_{ij}^f=<f_{1j}f_{2i}f_{2j}f_{1i}>`$, $`\mathrm{\Delta }_{ij;\alpha \beta }^b=<b_{i\alpha }b_{\beta j}>`$ and $`\mu _i=\mu _0\frac{J}{2}_{j=i\pm \widehat{x},i\pm \widehat{y}}|\mathrm{\Delta }_{ij}^f|^2`$. The Hubbard Stratonovich field $`\rho _i^{k=1,2,3}=<\frac{1}{2}f_i^{}\sigma ^kf_i>`$ for direct channel is taken to be $`0`$ and $`\rho _i^{k=0}=\frac{1}{2}`$. Owing to the energy cost the exchange interaction terms(the last two positive energy terms in Eq.(43)) is usually ignored-.
We now introduce the uniform hopping order parameter, $`\chi _{ij}=\chi `$, the d-wave spinon pairing order parameter, $`\mathrm{\Delta }_{ij}^f=\pm \mathrm{\Delta }_f`$ with the sign $`+()`$ for the nearest neighbor link parallel to $`\widehat{x}`$ ($`\widehat{y}`$) and the s-wave holon pairing order parameter, $`\mathrm{\Delta }_{ij;\alpha \beta }^b=\mathrm{\Delta }_{\alpha \beta }^b`$ with the boson indices $`\alpha `$ and $`\beta `$. For the case of $`\mathrm{\Delta }_{\alpha \beta }^b=0`$, $`\lambda ^{(k)}=0`$ and $`\mathrm{\Delta }^f\chi `$, the $`b_1`$-bosons are populated at and near $`k=(0,0)`$ in the momentum space and the $`b_2`$-bosons, at and near $`k=(\pi ,\pi )`$. Pairing of two different($`\alpha \beta `$) bosons(holons) gives rise to the non-zero center of mass momentum. On the other hand, the center of mass momentum is zero only for pairing between identical($`\alpha =\beta `$) bosons. Thus writing $`\mathrm{\Delta }_{\alpha \beta }^b=\mathrm{\Delta }_b(\delta _{\alpha ,1}\delta _{\beta ,1}\delta _{\alpha ,2}\delta _{\beta ,2})`$ for pairing between the identical holons and allowing the uniform chemical potential, $`\mu _i=\mu `$, the mean field Hamiltonian from Eq.(43) is derived to be,
$`H^{MF}`$ $`=`$ $`NJ(1\delta )^2\left({\displaystyle \frac{1}{2}}\chi ^2+\mathrm{\Delta }_f^2+{\displaystyle \frac{1}{4}}\right)+NJ\mathrm{\Delta }_f^2(2\mathrm{\Delta }_b^2+\delta ^2)`$ (44)
$`+`$ $`{\displaystyle \underset{k}{}}E_k^f(\alpha _{k1}^{}\alpha _{k1}\alpha _{k2}\alpha _{k2}^{})`$ (45)
$`+`$ $`{\displaystyle \underset{k,s=1,2}{}}\left[E_{ks}^b\beta _{ks}^{}\beta _{ks}+{\displaystyle \frac{1}{2}}(E_{ks}^b+\mu )\right]+\mu N\delta .`$ (46)
Here $`E_k^f`$ and $`E_{ks}^b`$ are the quasiparticle energies of spinon and holon respectively. $`\alpha _{ks}(\alpha _{ks}^{})`$ and $`\beta _{ks}(\beta _{ks}^{})`$ are the annihilation(creation) operators of the spinon quasiparticles and the holon quasiparticles respectively.
From the diagonalized Hamiltonian Eq.(46), we readily obtain the total free energy,
$`F`$ $`=`$ $`NJ(1\delta )^2\left({\displaystyle \frac{1}{4}}+\mathrm{\Delta }_f^2+{\displaystyle \frac{1}{2}}\chi ^2\right)`$ (50)
$`2k_BT{\displaystyle \underset{k}{}}ln[\mathrm{cosh}(\beta E_k^f/2)]`$
$`+NJ\mathrm{\Delta }_f^2(\mathrm{\Delta }_b^2+\delta ^2)+k_BT{\displaystyle \underset{k,s}{}}ln[1e^{\beta E_{ks}^b}]`$
$`+{\displaystyle \underset{k,s}{}}{\displaystyle \frac{E_{ks}^b+\mu }{2}}+\mu N\delta .`$
The chemical potential is determined from the number constraint of doped holes,
$`{\displaystyle \frac{F}{\mu }}={\displaystyle \underset{k}{}}[{\displaystyle \frac{1}{e^{\beta E_{k1}^b}1}}{\displaystyle \frac{ϵ_k^b\mu }{E_{k1}^b}}+{\displaystyle \frac{1}{2}}({\displaystyle \frac{ϵ_k^b\mu }{E_{k1}^b}}1)`$ (52)
$`+{\displaystyle \frac{1}{e^{\beta E_{k2}^b}1}}{\displaystyle \frac{ϵ_k^b\mu }{E_{k2}^b}}+{\displaystyle \frac{1}{2}}({\displaystyle \frac{ϵ_k^b\mu }{E_{k2}^b}}1)]N\delta =0,`$
and the Lagrangian multipliers are determined by the following three constraints imposed by the SU(2) slave-boson theory,
$`{\displaystyle \frac{F}{\lambda ^{(k)}}}={\displaystyle \underset{k}{}}\mathrm{tanh}{\displaystyle \frac{\beta E_k^f}{2}}{\displaystyle \frac{E_k^f}{\lambda ^{(k)}}}`$ (54)
$`+{\displaystyle \underset{k,s}{}}{\displaystyle \frac{e^{\beta E_{ks}^b}+1}{2(e^{\beta E_{ks}^b}1)}}{\displaystyle \frac{E_{ks}^b}{\lambda ^{(k)}}}=0,\text{ }k=1,2,3.`$
It can be readily proven from Eq.(54) above that $`\lambda ^{(k)}=0`$ satisfies the three constraints above.
By minimizing the free energy, the order parameters $`\chi `$, $`\mathrm{\Delta }_f`$ and $`\mathrm{\Delta }_b`$ are numerically determined as a function of temperature and doping rate. In Fig.2 the mean field results of the U(1) slave-boson theory(dotted lines) are displayed for $`J=0.2`$ $`t`$, $`J=0.3`$ $`t`$ and $`J=0.4`$ $`t`$ for comparison with the predicted phase diagrams(solid lines). The predicted pseudogap(spin gap) temperature, $`T_{SU(2)}^f`$ is consistently higher than $`T_{U(1)}^f`$, the U(1) value. $`T_{SU(2)}^b`$ at optimal doping is predicted to be lower than the value of $`T_{U(1)}^b`$ predicted by the U(1) theory. The predicted optimal doping rate is shifted to a larger value, showing closer agreement with observation than the U(1) mean field treatment. Such discrepancies are attributed to the phase fluctuations of order parameters, which were not treated in the U(1) mean field theory. We note from the four boson operator $`\frac{J}{2}|\mathrm{\Delta }_{ij}^f|^2b_{\alpha i}^{}b_{\beta j}^{}b_{\beta j}b_{\alpha i}`$ in the fourth term of Eq.(32) that the strength of holon pairing depends on the spinon pairing amplitude(order parameter) $`\mathrm{\Delta }_{ij}^f`$. Accordingly the predicted holon pair condensation temperature(superconducting transition temperature) $`T_{SU(2)}^b`$ depends on the spin gap(pseudogap) temperature $`T^{}`$; $`T_{SU(2)}^b`$ decreases with $`T^{}`$ in the overdoped region. Indeed it is shown in Eq.(10) that the predicted holon pair bose condensation at $`T_c`$($`=T_{SU(2)}^b`$) is not independent of the spin gap(pseudogap) formation at $`T^{}`$, by exhibiting the diminishing trend of superconducting temperature $`T_c`$ as the spin gap temperature $`T^{}`$ decreases in the overdoped region. This is consistent with an experimental observation of the universal behavior of $`T^{}/T_c`$ as a function of hole doping rate $`\delta /\delta _o`$ for different high $`T_c`$ cuprates, as is shown in Fig.1.
In summary, based on the SU(2) slave-boson symmetry conserving t-J Hamiltonian which shows coupling between the charge and spin degrees of freedom, we derived a phase diagram of high $`T_c`$ cuprates which displays the bose condensation temperature of an arch shape as a function of hole doping rate. Unlike other previous studies which predicted a linear increase with the hole doping rate, this result is consistent with observation. We showed that the low energy fluctuations cause a shift of the optimal doping rate to a larger value and a suppression of the holon pair bose condensation temperature, thus allowing a closer agreement with observation compared to the U(1) case.
One(SHSS) of us acknowledges the generous supports of Korea Ministry of Education(BSRI-99). We thank Tae-Hyoung Gimm for helpful discussions.
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# Correlations of triggering noise in driven magnetic clusters
## I Introduction
Field driven disordered ferromagnets at low temperatures exhibit Barkhausen (BH) noise, a very important physical phenomenon, which is used for noninvasive characterization technique in commercial alloys. It has been recognized that measured BH noise exhibits scale free behavior when only the external field is varied for all values of disorder and driving fields in real experiments (a short summary of the experimental data can be found in Ref. ). This is in sharp contrast to some theoretical conclusions emphasizing fine tuning of the strength of disorder to a single critical value. It has been also attempted to understand the dynamics of domain walls which results in BH avalanches in terms of the dynamics of a sandpile model (for a recent review of sandpile models see ) and of the models of interface depinning . It should be stressed that in BH noise the presence of disorder plays an important role via pinning of domain walls. It remains to be understood how the motion of domain walls is affected by the pinning.
Different scaling behaviors of BH noise can be attributed to the domain structure, which is related to annealing and type of impurities, and nucleation and coalescence of domains, as well sa varying driving conditions. Annealing the samples in the applied anisotropic stress or in the magnetic field leads to a characteristic structure of extended system-size domains with a $`180^{}`$ domain walls parallel to the anisotropy axis . Numerical simulations with extended domain wall in two dimensions led to the conclusion that different scaling behavior can be expected in two limiting cases: (a) When disorder is weak BH response is dominated by motion of a single (extended) domain wall; (b) For strong disorder a multidomain structure occurs with many competing domain walls. This picture is in a qualitative agreement with experimental results in stress annealed Fe-B-Si and Fe-Co-B alloys. Therefore the two different universality classes can be related to surface and bulk criticality, respectively. Numerical simulations starting from an uniform ground state , on the other hand, do not take into account extended domain walls, and thus correspond to the behavior at high-disorder. In both cases, however, the origin of scaling has not been fully understood.
Recently the exact results of the random-field Ising model on the Bethe lattice show that the avalanche distributions at a fixed driving field have finite cutoffs. However, an infinite cutoff appears for a range of disorder strengths $`\mathrm{\Delta }<\mathrm{\Delta }_c`$ if a distribution is integrated over the hysteresis loop. The integration thus involves an infinite jump in magnetization (when system size $`L\mathrm{}`$) at a critical value of driving field. On the other hand, for strong disorder $`\mathrm{\Delta }`$ above $`\mathrm{\Delta }_c`$ the avalanche distributions remain exponentially damped for all fields . These results encourage further study of the field-integrated distributions to elucidate the role of driving field in the appearance of the scale-free BH noise.
In this work we study properties of the threshold driving field which triggers Barkhausen avalanches in the multidomain structure, which is generated in the linear part of the ascending branch of hysteresis. When the system is slowly driven by increasing the external magnetic field $`H`$ with time, even by continuous field changes, there exists a threshold field that equals weakest pinning force of a domain wall in the system and thus starts an avalanche. In the numerical experiment we can adjust the driving field updates to the weakest local field in a system (so-called infinitely slow driving), and thus we can examine precisely the properties of the triggering noise. Surprisingly, we find that time series of such field updates are long-range correlated. Both Fourier spectrum of the threshold field fluctuations $`\mathrm{\Delta }H(t)H(t+1)H(t)`$ and distribution of distances of the successive avalanches which that field triggers appear to decay with a power law. The exponents are (weakly) universal in a range of values of disorder where BH avalanches have a power-law distributions. In addition, we find that the distribution of trapping times of domain wall at a given point in space exhibits scaling behavior. We show how the long-range noise correlations and fractal properties of the trapping time distribution are related to the observed scaling behavior of BH avalanches.
The present study is motivated by recently observed pattern formation and activity correlations along a driven interface in $`1+1`$ dimensions stacked by random defects. In contrast to these models, here we have a 2-dimensional system with many interacting interfaces, and two separate time scales (compared to extremal driving): slow time scale of field updates, and avalanche propagation time scale between two field updates. Despite of these differences, in both cases the intermittent avalanche-like dynamics occurs, which we believe is essential for the observed scaling behavior.
## II Model and simulations
We consider a simple model with disorder represented by local random fields $`h_x`$, which appears from an original disorder via coarse-graining
$$=\underset{<x,x^{}>}{}J_{x,x^{}}S_xS_x^{}\underset{x}{}(h_x+H)S_x,$$
(1)
where $`x(x_{},x_{})`$ and $`J_{x,x^{}}=1`$ is a constant interaction between nearest-neighbor spins $`S_x=\pm 1`$. A Gaussian distribution of $`h_x`$ is assumed with zero mean and width $`f`$. (Other types of disorder have been also considered ). A domain wall of reversed spins is created along $`<11>`$ direction on the square lattice rotated by $`\pi /4`$. This model is motivated by the extended domain walls in stress-annealed samples, as discussed in Ref. . Periodic boundaries are applied in the direction of wall and an open boundary at the opposite side of the wall. The system is driven globally—updated value of the external field is applied to all spins in the system. The dynamics consists of spin flips when the local field exceeds zero (see below).
As discussed in Ref. , the fact that motion of the $`11`$ domain wall has no energy threshold at vanishing disorder has several advantages. In particular, the wall depinning occurs along a line $`H_c(f)\kappa f`$ in the $`(f,H)`$ plane, where $`\kappa `$ is a constant (see Fig. 2 in Ref. ). Therefore for the domain wall along $`11`$ direction an infinitesimally small field $`H`$ is sufficient to move and depin the domain wall for small disorder, $`H_c(f)0`$ for $`f0`$. This important property of the model allows us to use much smaller lattice sizes compared to earlier studies , where a large system has to be used in order to find a random field large enough to surmount the energy barrier $`2J`$ in the case of wall along $`10`$ direction, or the nucleation energy $`4J`$ in the case of uniform ground states. On the other hand, by increasing disorder the distance between pinning centers decreases, and thus a critical disorder $`f^{}`$ exists at which depinning is no longer possible. Instead, it becomes energetically easier to nucleate new domains in the bulk. The exact value of $`f^{}`$ depends on the type of distribution , and is still not known exactly. In Ref. it was found using a finite-size scaling analysis of the averaged interface velocity and avalanche distributions that for the present model $`f^{}=0.62`$ within numerical error bars.
The scaling properties of BH avalanches in low disorder regime were discussed in Ref. . Here we are interested in the high disorder region, where the $`11`$ interface remains pinned at all fields. Nevertheless, the initial state with the $`11`$ interface ensures that the ratio of the characteristic lengths $`\xi _1/L1`$ holds for the applied disorder values (thus the avalanche size cut-off $`s_0L^2`$). (The characteristic length $`\xi _1`$ represents distance between strong pinning centers for given strength of disorder.) On the contrary, when $`\xi _1/L1`$, the $`11`$ interface moves, meaning that the selected system size $`L`$ is small for given strength of pinning, and thus system is not in the high disorder regime. Thus, having the condition $`\xi _1/L1`$ satisfied for each applied value of disorder, we can concentrate on the effects of disorder fluctuations by applying many configurations of random fields at fixed $`f`$ and $`L`$ values. We use up to $`L=768`$ and up to $`10^3`$ disorder configurations (i.e., up to $`36\times 10^6`$ spins).
Since we are interested in the properties of the triggering noise, which are related to distribution of disorder and the actual dynamics of the system, we avoid any “short-cuts” which can speed the algorithm . We apply the natural slow algorithm, which consists of the following steps: Random fields are stored at each site of an $`L\times L`$ lattice; The system is searched for the minimum local field $`h_{min}=\mathrm{min}(h_x^{loc})`$, where $`h_x^{loc}_x^{}J_{x,x^{}}S_x^{}+h_x+Hh_{ir,x}+H`$, and then the driving field is set to exactly $`h_{min}+\eta `$ (we use $`\eta =10^{10}`$), which thus triggers an avalanche at that site; A list is made consisting of the sites which may flip at the next time step (neighbor sites to flipped spins); The spins on the list are examined for flip and a new list is made; The process is continued until no more spins can flip, then next minimum local field is found. It should be stressed that all spins in a single list (spin shell) are updated in parallel, i.e., in a single time step, in analogy to parallel update in cellular automata models. In this way we have well defined time scale of avalanche evolution (internal time scale) as the number of steps that the updating procedure goes before the avalanche stops. On the other hand, the slow (external) time scale is set by the number of driving field updates.
It was shown in Ref. that for low disorder $`f<f^{}=0.62`$ the extended domain wall moves through the system and depinns at a critical field $`H_c(f)`$, however, above the critical disorder $`f^{}`$ the built-in domain wall remains pinned and many new domains of reversed spins are nucleated inside the system. In this region, corresponding to high disorder (or low tensile stress), a typical structure of clusters which occurs in the linear part of the hysteresis is shown in Fig. 1. The theoretical value $`f^{}`$ can, in principle, be related to a critical value of tensile stress $`\sigma ^{}f_{}^{\mu }{}_{}{}^{}`$, below which the domain structure changes, as also observed in the experiment in Ref. . More precise experimental data would be necessary in order to determine the exponent $`\mu `$. A rough estimate is that $`\mu `$ is given by the correlation length exponent at the transition $`f^{}`$, $`\mu \nu 2.3`$ .
## III Fractal properties of triggering noise and domain-wall trapping
Time series of the magnetization jumps corresponding to the individual avalanches are known as BH noise. In the high disorder region either a new domain is nucleated and grown, or already existing domain extended, a BH pulse is associated with motion of a domain wall from one position to a new one. The area between two consecutive domain wall positions corresponds to the size of BH avalanche (measured as the area covered by a single BH pulse). In principle, a domain growth process is not “linear” but fractal, leading to the dynamic exponent $`z=1.23`$, and fractal dimension of avalanche $`D=1.88`$ . As usual, the dynamic exponent $`z`$ and the fractal dimension $`D`$ are defined via the scaling of characteristic duration and size of avalanches with the change of length scale, $`t_LL^z`$, and $`s_LL^D`$, respectively. Scaling properties are studied in detail (see set of scaling exponents in Ref. ). In particular, in the critical region above $`f^{}`$ the distribution of size of avalanches $`D(s)s^{\tau _s}`$ and duration $`P(t)t^{\tau _t}`$ scale with the exponents $`\tau _s=1.30`$ and $`\tau _t=1.47`$, respectively. The avalanche distributions integrated over hysteresis loop for $`f>f^{}`$ can be scaled according to the scaling form
$$P(t,f,L)=t^{\tau _t}𝒫(t(\delta f)^{z\nu },tL^z),$$
(2)
where $`\delta f(f/f^{}1)`$ and $`(\delta f)^\nu `$ measures the correlation length due to the critical point at $`f^{}`$. For instance, when the condition $`L(\delta f)^\nu `$ is satisfied, the avalanche cut-offs can be scaled with disorder strength (in the entire region of power-law behavior), leading to $`f^{}=0.62`$ and $`\nu =2.3`$ within statistical error bars (see for details). It should be stressed that in this region of disorder the spatial extension of an avalanche is determined not only by strength of pinning but also by the blocking by previous clusters (see Fig. 1). The blocking effects dynamically alter the strength of pinning and thus influence the scaling properties. This illustrates a complex interplay between the disorder and dynamics, in contrast to, for instance, the equilibrium clusters in the random-field Ising model. In order to get an insight into these complex phenomena, we study next the distributions of wall trapping times and distances between the points with weakest pinning, from which an avalanche is released. We find that both of these distributions exhibit a scaling behavior.
The distribution of distances between initial points of the successive avalanches is shown in Fig. 2, where we distinguish between distances measured in the parallel ($`x_{}`$) and transverse() direction relative to the direction of the wall. In fact, due to anisotropy of domain wall motion in this case, these distributions show different scaling exponents $`\tau _{}=1.04\pm 0.03`$ and $`\tau _{}=0.60\pm 0.03`$ according to
$$G(x_{},x_{})x_{}^\tau _{}𝒢(x_{}/x_{}^{\zeta _G}),$$
(3)
where $`\zeta _G=\tau _{}/\tau _{}`$. Notice that for small distances correlations in both directions are almost equivalent. However, a crossover to a distinct power-law behavior for $`G(x_{})`$ occurs at a distance ($`x30`$ in Fig. 2) which is presumably related to the characteristic size of avalanche at given disorder. The open boundary in the direction opposite to wall leads to the cut-off at large $`x_{}`$. It should be stressed that spatial correlations are sensitive to distribution of disorder. Therefore, here we used a large number of samples in order to minimize scatter of the data due to disorder fluctuations and to clearly distinguish between correlations in parallel and perpendicular directions.
In the inset to Fig. 2 we show the Fourier spectrum of the threshold field fluctuations $`\mathrm{\Delta }H(t)`$ measured on the external time scale. It exhibits two correlated regions corresponding to earlier and later times, respectively. Steeper slope in the inset to Fig. 2 vary from 0.6-1.03, and flatter slope from 0.25 to 0.36, depending on $`f`$. The long-range correlations exhibited in Fig. 1 show that the system evolves in such way that its next relaxation event depends on the history of the present state of the system. It is interesting to note that a similar time series of the fluctuation of magnetization $`\mathrm{\Delta }M(t)M(t+1)M(t)`$ represents the Barkhausen noise itself (looked at the external time scale). An example of the BH noise signal is shown in Fig. 3. Note that on the external time scale duration of each elementary signal is equal to one, whereas the height of each elementary signal represents corresponding avalanche size. Sizes of the successive avalanches are only weakly correlated in time. However, a time derivative of the signal, namely $`a(t)d[\mathrm{\Delta }M(t)]/dt`$, representing acceleration of a domain wall at each field update, shows certain correlation properties. In the inset to Fig. 3 we show the Fourier spectrum of the numerical derivative of the the signal, which is shown in main Fig. 3. Average slope of the curve is close to one.
It should be noted that the average number of avalanches occurring at fixed strength of disorder increases with system size as $`n_aL^2/sL^{[2D(2\tau _s)]}\varphi (L(\delta f)^\nu )`$, where $`\varphi (L(\delta f)^\nu )`$ is unknown scaling function. This implies that in the theoretical case of infinitely slow driving the average size of the field jump $`\mathrm{\Delta }H(t)`$ decreases with $`L`$ as $`\mathrm{\Delta }H(t)H_{sat}(L)/n_aL^{0.68}H_{sat}(L)/\varphi (L(\delta f)^\nu )`$, where $`H_{sat}`$ is the saturation field. In real experimets the size of field jumps are restricted by the driving frequency, however, number of avalanches detected varies with the size of pick-up coil. This problem requires more detailed theoretical and experimental investigation .
Further understanding of the role of disorder in the dynamics can be achieved by considering the time intervals that a domain wall resides at a given point in space (trapping time $`T_{trap}`$). We calculate the distribution of trapping times $`T_{trap}`$ of a domain wall at a given site, which is determined as the time interval since the domain wall is pinned at a site $`(x_{},x_{})`$ until it eventually moves away from that site. In this way, $`T_{trap}`$ measures the time interval between two successive activities at that site. Notice that $`T_{trap}`$ in this case is somewhat different from so called first-return time in 1-dimensional interface . Here we have 2-dimensional system with many interacting interfaces. Another similar example is trapping of grains in rice-pile models . A reasonable time scale to measure $`T_{trap}`$ is the number of field updates (external time scale), since in the zero-temperature dynamics during an avalanche a spin at a given site is either fixed or flips only once.
The trapping-time distribution is shown in the inset to Fig. 4 for various values of disorder $`f>f^{}`$. Between the lower cut-off $`T_0`$ (below which all trapping times are equally probable), and an upper cut-off due to lattice size $`L`$, the trapping time distribution shows a short region with correlations. In particular, for large trapping times we can determine the slope as $`\tau _{trap}=2.3\pm 0.1`$. One can understand that a metastable configuration of domain walls corresponds to the system residing in a “pocket” of the fractal free energy landscape. A domain wall may move from a site $`(x_{},x_{})`$ as a part of an avalanche which started in the neighborhood of that point, corresponding to a local reorganization of the landscape near a shallow minimum (for a given value of the external field). When the system resides in a deeper minimum, however, it waits for larger driving fields (i.e., trapping time increases) or for a more global reconstruction of the landscape, which occurs with smaller probability compared to the one discussed above. This leads to the large slope of the distribution $`P(T_{trap})T_{trap}^{\tau _{trap}}`$ at large $`T_{trap}`$. It is interesting to note that a similar slope for large trapping times was found in the case of rice-pile model , where trapping of grains are considered, and it was argued that the slope is related to the roughness exponent of the rice-pile surface (see detailed analysis in Ref. ). Analogies between an interface motion and rice-pile model have been established in the literature . Considering the avalanche exponents the analogy also applies for the BH noise with an extended domain wall in the limit of low disorder . Notice that, in contrast to ricepile model, the range of power-law behavior in BH noise depends not only on the system size, but also on the strength of disorder, which restricts spatial extension of avalanches to $`ss_{max}L^2`$ and their durations to $`TT_{max}`$. Therefore, in this region of disorder we have rather small range of the power-law behavior. However, the scale-free behavior of the distribution of trapping times can be demonstrated via finite size scaling analysis when different system sizes are used. We find that the following scaling form
$$P(T_{trap},L)=L^\alpha 𝒫(T_{trap}L^{z_T}),$$
(4)
applies with the exponents which are weakly dependent on disorder $`\alpha =0.4\pm 0.05`$ and $`z_T=1.66\pm 0.05`$ (see Fig. 4). By increasing the system size $`L`$ the trapping times increase, leading to larger cut-offs of the distribution. We notice that the lower cut-offs $`T_0`$ also increase, leaving a limited correlation range. However, the cut-offs scale nicely with $`L`$, as shown in Fig. 4. In the inset to Fig. 4 we show how the trapping times distribution varies with disorder. By increasing disorder the cut-off $`T_0`$ moves towards larger values and the correlation region shrinks, corresponding to a lesser correlations in the dynamics, which also manifests in the decrease of the characteristic avalanche size. In the limit of infinitely strong disorder the dynamics becomes completely random, consisting of individual spin flips which align along a local random field.
It is interesting to note that unlike the avalanche distributions in Eq. (2), the distribution of trapping times in Eq. (4) as well as $`G(x)`$ (3) do not show explicit disorder dependence apart from a weak dependence in the exponents. In fact, the disorder effects in this case are included in the gradient of driving force $`\mathrm{\Delta }F`$ (i.e., the increments of the external field $`\mathrm{\Delta }H`$, which are adjusted to the minimum local fields). For example, for the distribution of trapping times we have in general
$$P(T_{trap},L,\mathrm{\Delta }F)=L^\alpha 𝒫(T_{trap}(\mathrm{\Delta }F)^{z_T/\lambda _F},T_{trap}L^{z_T}),$$
(5)
where the increase in the driving force $`\mathrm{\Delta }FH(t+T_{trap})H(t)`$ contains all the contributions due to interaction and random pinning which occurred during the time interval $`(t,t+T_{trap})`$, i.e., $`\mathrm{\Delta }F=_t^{t+T_{trap}}h_{ir,x}`$, which steam from the different sites in the system. According to the above results in Fig. 2, sum of these contributions has no characteristic scale. Therefore we may conclude that $`(\mathrm{\Delta }F)^{1/\lambda _F}\mathrm{}`$. Hence the first argument in the right hand side of Eq. (5) can be neglected compared to $`T_{trap}/L^{z_{trap}}`$, leading to the scaling form (4), which is in agreement with the scaling plot in Fig. 4.
Another interesting observation regards the dynamic exponent $`z_T`$, which scales the tail of the trapping time distribution with the system size $`L`$. It may be related to the scaling of the length of the optimal path in strongly disordered medium, $`z_T=D_{OP}`$. In fact, the optimal path between two points can be constructed from the most persistent sections of the domain walls in the dynamics of BH avalanches. The length of the optimal path scales with the linear distance between the end points with the exponent $`D_{OP}=\tau _{OP}=1.66`$ .
## IV Transport equation with noise correlations
The long-range noise correlations are shown to be relevant for the scaling properties of the interface motion. In the literature the noise correlations in the interface depinning are viewed as originating from another (external) dynamic processes . In the case of Barkhausen avalanches we see that the noise correlations appear as an intrinsic property of the dynamics when the system is driven infinitely slowly. The active section of domain wall can be represented by a surface $`h(x,\tau )`$ which is pinned by quenched defects. The transport at the site $`(x,\tau )`$ is described by the equation
$$dh/d\tau =\nu _{}_{}^2h+\nu _{}_{}^2h+\eta (x,t,h).$$
(6)
Here we distinguish total elapsed time $`\tau `$, which is defined as the sum of evolution times of all individual avalanches, on one hand, and $`t`$, which represents time scale of field updates, on the other. The quenched noise $`\eta (x,t,h)`$ is generated by varying the driving field in the steps which are adjusted to the minimum local strength of pinning, as discussed above. We expect that a dominate $`h`$-dependence of the noise $`\eta `$ can be expressed as $`\eta p(x)h`$, where $`p(x)`$ represents the (anisotropic) local velocity of domain wall per field rate. In addition, time variation of the noise are related to the updates of the external field, and thus $`\eta `$ varies on the external time scale only. Therefore, we can write
$$\eta (x,t,h)p(x)(a_0+\mathrm{\Delta }H(t))(h),$$
(7)
where $`a_0`$ is a constant which depends on initial configuration and the hysteresis loop properties. The local interface velocity per field rate, $`p(x)`$, is a random variable which is determined by the spatial distribution and strength of pinning associated with a given value of the driving field $`H`$, i.e., history of the system. Therefore $`p(x)`$ is governed by a probability distribution, which can be deduced from the properties of the probability that an avalanche starts at distance $`x`$ from the preceding avalanche, i.e.,
$$<p(x)p(x^{})>=\gamma G(xx^{}).$$
(8)
It is assumed that the same type of correlations apply for the successive activities during the evolution of an avalanche, while the external field is constant. In the previous section we found that the probability $`G(x)`$ exhibits long-range correlations in the perpendicular and parallel direction as $`G(x_{},x_{})x_{}^1x_{}^{0.60}`$. We also notice that $`\mathrm{\Delta }H(t)t^\theta `$, where $`\theta `$ is finite in the case of infinitely slow driving discussed above, whereas, $`\theta =0`$ in the case of finite driving rate in small steps , where $`\mathrm{\Delta }H(t)=const`$. Notice that Eq. (6) together with Eqs. (7) and (8) leads to an effective nonlinear term of the form $`\gamma G(x_{},x_{})(a_0+t^\theta )^2(h)^2`$, which differs from the Kardar-Parisi-Zhang nonlinearity by the power-law correlations in the coefficient. As discussed above, the origin of these correlations lies in the the dynamically varying pinning and blocking effects when a multidomain structure is slowly driven through the hysteresis loop.
The relevance of the noise correlations for the universality class of the interface depinning has been discussed in the literature using dynamic renormalization group (RG) . Another example where correlations of random noise play an important role is represented by the scaling properties of river networks . In that case the correlations of the type $`Gx_{}^{2+\delta }x_{}^{1\zeta }`$ were assumed, where $`\delta `$ is an expansion parameter and the anisotropy exponent $`\zeta `$ has to be determined self-consistently at a fixed point of the RG . In the result, the scaling exponents do depend on the range of correlations $`\delta `$. Similarly, we can expect that the long-range correlations of the form given in Eqs. (7) and (8) can be related to the universal scaling exponents of the BH noise. The exponents $`z`$, $`\zeta _G`$, which are defined above, and the roughness exponent $`\chi `$, which governs behavior of the dynamic variable $`h`$ with the change of scale, can in principle be obtained by the dynamic RG applied to Eq. (6) with noise properties in Eqs. (7-8). Then the avalanche exponents $`\tau _t`$ and $`\tau _s`$ can be deduced using scaling relations in the critical region. The RG analysis of the transport equation (6) with noise properties specified by Eqs. (7-8) requires additional work and is left out of the scope of this paper. Here we discuss only the scaling relations between the exponents at RG fixed point, e.g., the dynamic exponent $`z`$, and the avalanche exponents.
It should be stressed that such scaling relations are not universal and that they depend on the nature of the dynamic process (see for instance for the case of river networks). We argue that the scaling relation derived below are valid for BH avalanches in the high-disorder (or multidomain) region. The rational behind these scaling relations is found in the directed nature of the avalanche propagation. The evolution of an avalanche (see Fig. 1) can be visualized as growth of a directed percolation (DP) cluster projected to $`1+1`$ dimensions, with extra dimension representing time axis. The fractal dimension of the equivalent DP cluster measured with respect to the time axis is then $`D_{}D/z`$, where $`D`$ and $`z`$ are fractal dimension of BH avalanche and dynamic exponent, respectively, as defined above. Then $`D_{}=1+\zeta _{DP}\delta _{DP}`$, where $`\zeta _{DP}=z/2`$ is the anisotropy exponent measured with respect to time axis, and $`\delta _{DP}`$ is the survival probability exponent of the equivalent DP clusters. Notice also that for the directed dynamic processes $`\tau _t=D_{}`$ and thus $`\tau _s=21/\tau _t`$ , which completes the set of avalanche exponents. In addition, the roughness exponent $`\chi `$ can be related to the trapping time distribution as $`\tau _{trap}=2+\chi `$, as noticed in Refs. . For instance, using the numerical values for $`D`$ and $`z`$ from Ref. , the following values of the avalanche exponents are predicted by the above scaling relations: $`\tau _t=1.52`$, $`\tau _s=1.34`$, compared with $`1.47`$ and $`1.30`$, and $`\chi =0.23`$, obtained by direct numerical simulations.
## V Conclusions and discussions
We have demonstrated numerically the existence of long-range correlations in triggering noise which is intrinsic to the domain wall dynamics in slowly driven disordered ferromagnets. Although the distribution of disorder is initially uncorrelated, the spatio-temporal correlations develop in time. The appearance of these correlations can be attributed to the applied global driving in which always next weakest pinning force in the system is selected, on one hand, and to a finite extension of avalanches occurring between two consecutive field updates, on the other. Hence, attempts to reduce the domain wall dynamics to the Glauber spin flips in the presence of local random fields appears to be inadequate for the range of disorder where the cooperative avalanche dynamics occurs. We suggest an alternative transport equation which incorporates the observed noise correlations. The long-range correlations of triggering noise can be related to the universal scaling behavior of Barkhausen avalanches and to the transport properties via the fractal distribution of the domain wall trapping times. We believe that an analysis of the transport equation by the dynamic RG will contribute to understanding of the universality classes of Barkhausen avalanches by the infinitely slow driving. By varying the driving conditions, however, these correlations are changed. This might be the origin of different scaling behavior of BH avalanches at finite driving rates, as observed in experiments . Our results also suggest that technical procedures which speed the numeric algorithms in driven disordered systems should be taken with care. In particular, an algorithm which alters the properties of triggering noise may lead to a different scaling behavior, which is unrelated with the original problem.
###### Acknowledgements.
This work was supported by the Ministry of Science and Technology of the Republic of Slovenia. I am grateful to Deepak Dhar for many helpful comments and suggestions.
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# Spectral properties and pseudogaps in a model with d-wave pairing symmetry
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## Abstract
A model with d-wave pairing symmetry is studied by employing a non-perturbative sum rule approach. At low temperature the magnitude of a normal state pseudogap shows strong $`\stackrel{}{k}`$ or angle dependence well fitted by $`\mathrm{cos}2\varphi `$ form. With increasing temperature, the pseudogap closes at some critical angle $`\varphi _c`$ and beyond this angle a single quasiparticle-like peak appears. The resulting Fermi surface is strongly temperature dependent. Both in the spectral function and the density of states, the pseudogap disappears in a manner that the spectral weight fills in the pseudogap instead of closing it with increasing temperature. All these features are qualitatively consistent with ARPES for underdoped cuprates.
\]
The nature of an excitation gap in high-temperature superconductors has been one of the puzzling issues in the community of condensed matter physicists. Through several years of extensive experimental work, general consensus regarding the superconducting gap symmetry seems to be reached that the superconducting gap has mainly d-wave character with possibility of a small mixture of other angular momentum states, in contrast to conventional BCS superconductors with an isotropic s-wave gap. Recent discovery of a normal state pseudogap in underdoped cuprates has shed another side of the anomalous behaviors in the copper oxide superconductors. For these materials the low frequency spectral weight begins to be strongly suppressed below some characteristic temperature $`T^{}`$ higher than $`T_c`$. This behavior has been observed through various experimental probes such as photoemission, specific heat, tunneling, NMR, and optical conductivity. In particular recent angle resolved photoemission spectroscopy (ARPES) and tunneling experiments indicate that the pseudogap phenomenon is closely related to pairing fluctuations. These measurements clearly exhibit that the normal state pseudogap has the same angular dependence and magnitude as the superconducting gap and that often the only difference between the spectra in the pseudogap state and the superconducting state is in their linewidths. Typically $`T^{}`$ is much higher than $`T_c`$ and their doping dependence is qualitatively different. While $`T_c`$ decreases with underdoping, $`T^{}`$ increases in contrast. This feature suggests that $`T^{}`$ does not follow $`T_c`$ characterized by long-range phase coherence, but instead some kind of a mean-field critical temperature $`T_{MF}`$. In spite of several possible scenarios such as the spinon pair formation without the Bose-Einstein condensation of holons, strong superconducting phase fluctuations and a magnetic scenario near the antiferromagnetic instability and so on, at present there is no consensus in the origin of the pseudogap.
For the past several years extensive theoretical effort has been also made by several groups to understand this anomalous pseudogap behavior in the context of short range effective (attractive) interaction between electrons. This may be divided into two different classes of approach. In the first class, quantum Monte Carlo (QMC) simulations were made for the attractive Hubbard model. Although this Hamiltonian is not a realistic model for understanding the complex physical behaviors in the underdoped cuprates, it is believed to capture the important ingredient of the paring fluctuations in those materials. In the absence of a small parameter this numerical method has played an important role in understanding of the model, in spite of some uncertainties due to finite size effect and numerical analytical continuation. In this approach the one-particle spectral function as well as the density of states clearly show the precursor of the superconducting gap in the normal state. In the second class, various assumptions and approximations are used. This class includes the effect of vortex phase fluctuations on the single-particle properties and the ‘paring approximation’ theory, and T-matrix and self-consistent T-matrix approximations for the two-dimensional attractive Hubbard model. More recently the self-consistent T-matrix approach was also applied to a model with $`d_{x^2y^2}`$ pairing. The present approach is a variant of the T-matrix approximation applied to a model with d-wave pairing symmetry.
In a previous study, a non-perturbative sum rule approach was developed for the attractive Hubbard model by extending previous work on the repulsive Hubbard model. It is found that in two dimensions, the mean-field transition temperature is replaced by a crossover temperature where the correlation length starts to grow exponentially. At sufficiently low temperature, a Kosterlitz-Thouless O(2) transition should occur, but it is not reproduced by that approach since it is in the $`O(n=\mathrm{})`$ universality class. Nevertheless, the agreement with Monte Carlo calculations is quantitative for both one and two-particle correlation functions over the whole range of parameters accessible by Monte Carlo calculations where it is found that in two dimensions, a pseudogap appears in the single-particle spectral weight as well as in the density of states. Recent dimensional crossover study by Preosti et al. shows that the pseudogap effect is basically absent in three dimensions. In weak to intermediate coupling, the appearance of a pseudogap is traced back to the growing critical pairing fluctuations in the low-temperature renormalized classical regime of the low-dimensional system. With increasing temperature, the spectral weight fills in the pseudogap instead of closing it. Furthermore, the pseudogap appears earlier in the density of states than in the spectral function. It was also noted that the qualitative features found in this study should apply to the d-wave case. In this paper we study in detail spectral properties and pseudogaps in a model Hamiltonian with d-wave pairing symmetry which is more appropriate for high-temperature superconductors.
We consider on a two dimensional square lattice a simple model Hamiltonian which has a superconducting ground state with d-wave symmetry
$`H`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{k},\sigma }{}}\epsilon _\stackrel{}{k}c_{\stackrel{}{k},\sigma }^+c_{\stackrel{}{k},\sigma }`$ (1)
$`+`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{\stackrel{}{k},\stackrel{}{k}^{},\stackrel{}{q}}{}}V_{\stackrel{}{k},\stackrel{}{k}^{}}c_{\stackrel{}{k},}^+c_{\stackrel{}{k}+\stackrel{}{q},}^+c_{\stackrel{}{k}^{}+\stackrel{}{q},}c_{\stackrel{}{k}^{},},`$ (2)
where $`\epsilon _\stackrel{}{k}=2t(\mathrm{cos}k_x+\mathrm{cos}k_y)`$ and $`N`$ is the number of lattice sites. We choose $`V_{\stackrel{}{k},\stackrel{}{k}^{}}`$ as a d-wave separable potential given as $`V\mathrm{\Pi }(\stackrel{}{k})\mathrm{\Pi }(\stackrel{}{k}^{})`$ where $`\mathrm{\Pi }(\stackrel{}{k})=\mathrm{cos}k_x\mathrm{cos}k_y`$, or as $`1.5264\mathrm{cos}2\varphi =1.5264\frac{k_x^2k_y^2}{k_x^2+k_y^2}`$ where $`\varphi =\mathrm{arctan}\frac{k_y}{k_x}`$. In the real space notation the interaction term may be written as $`H_\text{I}=V_i\mathrm{\Delta }_i^+\mathrm{\Delta }_i`$ where $`\mathrm{\Delta }_i=_\delta g(\delta )c_{i,}c_{i+\delta ,}`$. $`g(\delta )`$ is the Fourier transform of $`\mathrm{\Pi }(\stackrel{}{k})`$ given as $`g(\delta )=\frac{1}{N}_\stackrel{}{k}e^{i\stackrel{}{k}\stackrel{}{\delta }}\mathrm{\Pi }(\stackrel{}{k})`$. For example, for the following pair structure,
$`g(\delta )=\{\begin{array}{ccc}1/2\hfill & \text{if }\delta =(\pm 1,0),\hfill & \\ 1/2\hfill & \text{if }\delta =(0,\pm 1),\hfill & \\ 0\hfill & \text{if otherwise},\hfill & \end{array}`$
it leads to $`\mathrm{\Pi }(\stackrel{}{k})=\mathrm{cos}k_x\mathrm{cos}k_y`$. Obviously $`\mathrm{\Pi }(\stackrel{}{k})=1.5264\mathrm{cos}2\varphi `$ requires in $`g(\delta )`$ further neighbors beyond the four nearest neighbors. Here a numerical factor of 1.5264 is introduced to normalize $`g(\delta )`$ with $`_\delta g^2(\delta )=1`$. In terms of $`\mathrm{\Delta }_i`$ the d-wave pairing susceptibility is defined as
$`\chi _{pp}(i,\tau )=T_\tau \mathrm{\Delta }_i(\tau )\mathrm{\Delta }_0^+,`$ (3)
where $`T_\tau `$ is the imaginary time ordering operator.
It is well-known that for many strongly correlated electronic models there is no obvious small parameter with which a systematic perturbative approximation can be made. In such a situation, a non-perturbative approach may be a good alternative for at least qualitative understanding of the physics in those strongly correlated model Hamiltonians. In this paper we apply a non-perturbative sum rule approach to the model Hamiltonian (2), which is still not a fully controlled approximation in nature. Because of an extended nature of $`\mathrm{\Delta }_i`$ for a d-wave pair and the increasing difficulty of the derivation in the real space representation, here we use the analogy from our previous studies for the repulsive and attractive Hubbard models. These studies show the important many-body modification with respect to the standard RPA and T-matrix approaches. The modification comes in two different places. First, the paring susceptibility is calculated through vertex function $`U_{pp}`$ instead of bare interaction strength $`U`$, which is constant in our approximation. Second, the pairing fluctuations are coupled to electrons, leading to the self-energy with $`UU_{pp}`$ form instead of $`U^2`$. In particular, the latter structure is in agreement with the fact that there is no Migdal theorem for this problem, contrary to the case of electron-phonon interactions. And the renormalized interaction constant $`U_{pp}`$ is determined by the exact sum rule for the pairing susceptibility. By using the sum rule and the renormalized constant $`U_{pp}`$, we determine in effect the Ginzburg-Landau parameters due to mode-mode coupling. Thus this approach is similar to the one-loop renormalization group approach within the Gaussian approximation and as a result the Mermin-Wager theorem is formally satisfied in two dimensions. This sum rule approach was systematically compared with the QMC simulations for the repulsive and attractive Hubbard models and the agreement was in a quantitative level both in one- and two-particle functions. We believe that this generic feature of the many-body modification should also carry over to a Hamiltonian with d-wave pairing symmetry. In this modification, the pairing susceptibility and the self-energy can be written in Fourier space in terms of renormalized interaction strength $`V_{pp}`$
$`\chi _{pp}(q)`$ $`=`$ $`{\displaystyle \frac{\chi _{pp}^0(q)}{1+V_{pp}\chi _{pp}^0(q)}},`$ (4)
$`\mathrm{\Sigma }(k)`$ $`=`$ $`VV_{pp}\mathrm{\Pi }^2(\stackrel{}{k}){\displaystyle \frac{T}{N}}{\displaystyle \underset{q}{}}\chi _{pp}(q)G^0(qk),`$ (5)
where the irreducible susceptibility is defined as
$`\chi _{pp}^0(q)={\displaystyle \frac{T}{N}}{\displaystyle \underset{k}{}}\mathrm{\Pi }^2(\stackrel{}{k})G^0(qk)G^0(k).`$ (6)
The above expressions are reduced to the standard T-matrix approximation when $`V_{pp}`$ is replaced by bare $`V`$. We determine this renormalized constant by employing the exact sum rule for the d-wave pairing susceptibility
$`{\displaystyle \frac{T}{N}}{\displaystyle \underset{q}{}}\chi _{pp}(q)e^{i\nu _m0^{}}`$ (7)
$`=`$ $`{\displaystyle \underset{\delta ,\delta ^{}}{}}g(\delta )g(\delta ^{})c_{\delta ^{},}^+c_{\delta ,}c_{}^+c_{}.`$ (8)
Note that the convergence factor is necessary in the sum rule, because away from half-filling $`\chi _{pp}(i\nu _m)`$ decays like $`1/\nu _m`$ at large frequencies. In the previous studies for the repulsive and attractive Hubbard models, the right-hand side of the sum rule evaluated in the SDW and BCS mean-field ground states, respectively, was found to be in excellent agreement with the QMC values in the intermediate to strong coupling regimes. In this paper the right-hand side of Eq. (8) is also evaluated in the d-wave BCS mean-field ground state
$`{\displaystyle \frac{1}{N^2}}{\displaystyle \underset{\stackrel{}{k},\stackrel{}{p},\stackrel{}{k}^{},\stackrel{}{p}^{}}{}}\mathrm{\Pi }(\stackrel{}{k})\mathrm{\Pi }(\stackrel{}{p})c_{\stackrel{}{k},}^+c_{\stackrel{}{p},}c_{\stackrel{}{k}^{},}^+c_{\stackrel{}{p}^{},}\delta _{\stackrel{}{k}+\stackrel{}{k}^{},\stackrel{}{p}+\stackrel{}{p}^{}}`$ (9)
$`=`$ $`[{\displaystyle \frac{1}{N}}{\displaystyle \underset{\stackrel{}{k}}{}}\mathrm{\Pi }^2(\stackrel{}{k})v^2(\stackrel{}{k})][{\displaystyle \frac{1}{N}}{\displaystyle \underset{\stackrel{}{k}}{}}v^2(\stackrel{}{k})]`$ (10)
$`+`$ $`[{\displaystyle \frac{1}{N}}{\displaystyle \underset{\stackrel{}{k}}{}}\mathrm{\Pi }(\stackrel{}{k})u(\stackrel{}{k})v(\stackrel{}{k})]^2,`$ (11)
where
$`u_\stackrel{}{k}^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1+{\displaystyle \frac{\epsilon _\stackrel{}{k}\mu }{E_\stackrel{}{k}}}),`$ (12)
$`v_\stackrel{}{k}^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1{\displaystyle \frac{\epsilon _\stackrel{}{k}\mu }{E_\stackrel{}{k}}}),`$ (13)
$`E_\stackrel{}{k}`$ $`=`$ $`\sqrt{(\epsilon _\stackrel{}{k}\mu )^2+\mathrm{\Delta }^2(\stackrel{}{k})}.`$ (14)
Here $`\mathrm{\Delta }(\stackrel{}{k})`$ is the d-wave BCS mean-field gap. The chemical potential $`\mu `$ and the gap $`\mathrm{\Delta }(\stackrel{}{k})`$ are determined self-consistently through the number and gap equations for given $`V`$, $`T`$ and $`n`$.
Before starting we comment on some differences associated with $`\mathrm{\Pi }(\stackrel{}{k})=\mathrm{cos}k_x\mathrm{cos}k_y`$ and with $`\mathrm{\Pi }(\stackrel{}{k})=1.5264\mathrm{cos}2\varphi `$ structure. The interaction Hamiltonian with the first structure depends not only on the angle of $`\stackrel{}{k}`$ but also on its magnitude. Thus the pairing interaction is always strongest at $`\stackrel{}{k},\stackrel{}{k}^{}=(\pm \pi ,0)`$ or $`(0,\pm \pi )`$, although the noninteracting Fermi surface can be far from these points. As a result two features occur for a particle density far away from half-filling (Detailed calculations were performed but not shown in this paper). First, the peaks associated with the precursor of the Bogoliubov quasiparticles occur asymmetrically with respect to the Fermi energy. Second, the locus of $`\stackrel{}{k}`$ points satisfying $`\omega \epsilon _\stackrel{}{k}+\mu Re\mathrm{\Sigma }(\stackrel{}{k},\omega )=0`$ at $`\omega =0`$ can be substantially different from the noninteracting Fermi surface, thus strongly violating the Luttinger’s theorem. For $`\mathrm{\Pi }(\stackrel{}{k})=1.5264\mathrm{cos}2\varphi `$ which depends only on the angle, however, the above features disappear and the locus of $`\stackrel{}{k}`$ points satisfying the above equation is almost identical with the noninteracting Fermi surface. Near half-filling the differences are negligible. Throughout the calculations, lattice spacing, $`\mathrm{}`$, and $`k_B`$ are set to be unity, and all energies are measured in unit of $`t`$. We used a discrete lattice as large as $`128\times 128`$ in momentum space and performed the calculations by means of fast Fourier transforms (FFT). Equations (5) and (6) are computed in terms of Matsubara frequencies and the analytic continuation from Matsubara to real frequencies are made via Pade approximants. In order to detect any spurious features associated with numerical analytical continuation, we performed real frequency calculations in parallel. Except for some spiky features in the real frequency formulation coming from the Lorentzian approximation of the non-interacting Green’s function, the two results are almost identical.
We begin by presenting the spectral functions along $`(0,0)(\pi ,0)`$ direction for $`V=4t`$, $`n=0.5`$ and $`T=0.15t`$ in Fig. 1. Throughout the paper the density is fixed at $`n=0.5`$. The results for other densities are similar. For this parameter, $`V_{pp}`$ satisfying the sum rule is found to be $`1.87t`$ significantly different from the bare value ($`V=4t`$). This shows the importance of the mode-mode coupling effect already in the intermediate coupling regime. Let us label wave vectors by $`(m\pi /8,0)`$. Below the Fermi wave vector ($`m=5`$), the main peak stays below the Fermi energy and the secondary peak grows in strength as $`\stackrel{}{k}`$ approaches $`\stackrel{}{k}_F`$. At the Fermi wave vector two peaks reminiscent of the Bogoliubov quasiparticles appear almost symmetrically with respect to the Fermi energy. Since electrons are still in the normal state, this is the precursor of the superconducting gap, namely, a normal state pseudogap. Above the Fermi wave vector, the main peak stays above the Fermi energy and the secondary peak becomes stronger as $`\stackrel{}{k}`$ approaches $`\stackrel{}{k}_F`$. This result should be contrasted with that by Engelbrecht et al. In the self-consistent T-matrix approximation with $`\mathrm{\Pi }(\stackrel{}{k})=\mathrm{cos}k_x\mathrm{cos}k_y`$ form, these authors argued that along $`(0,0)(\pi ,0)`$ direction the dominant peak of $`A(\stackrel{}{k},\omega )`$ never crosses the Fermi energy and bounces back towards the negative frequency. Their overall finding, however, is qualitatively different from our results. In the present calculations, the dominant peak of $`A(\stackrel{}{k},\omega )`$ eventually passes through the Fermi energy for $`\stackrel{}{k}`$ far above the Fermi wave vector. At the noninteracting Fermi surface our spectral weight is strongly suppressed at the Fermi energy to become a local minimum, while the local minimum of $`A(\stackrel{}{k},\omega )`$ in their study is located significantly away from the Fermi energy. Presumably this is due to the features associated with $`\mathrm{\Pi }(\stackrel{}{k})=\mathrm{cos}k_x\mathrm{cos}k_y`$ structure mentioned in the previous paragraph and also due to the self-consistent approximation that does not take the vertex function and the Green’s function at the same level of approximation.
In Fig. 2 the spectral function is shown along the Fermi wave vectors at low temperature ($`T=0.15t`$). The angle $`\varphi `$ is defined as $`\mathrm{arctan}(k_y/k_x)`$ along the noninteracting Fermi surface. At $`\varphi =0^o`$ the magnitude of the pseudogap is largest and as $`\varphi `$ increases it progressively decreases, leading to strong momentum or angle dependence in the size of the pseudogap. Since along the diagonal directions the pairing interaction vanishes, the pseudogap completely closes and a sharp quasiparticle peak appears at $`\varphi =45^o`$. In Fig. 3 the magnitude of the pseudogap (circles) is plotted as a function of angle $`\varphi `$ for $`T=0.15t`$ along with the ground state gap symmetry $`\mathrm{\Delta }(\varphi =0^o)\mathrm{cos}2\varphi `$ (dashed curve). The angle dependence of the normal state pseudogap is well fitted by d-wave symmetry, consistent with ARPES experiments for underdoped cuprates.
At higher temperatures a drastic change is found in the spectral function. For $`T=0.225t`$ the pseudogap closes well below $`45^o`$, as shown in Fig. 4. The local minimum at the Fermi energy disappears at $`\varphi =31.0^o`$ ($`26.6^o`$ is just on the crossover). Beyond this angle, the quasiparticle-like peak appears and thereby the Fermi surface is partially restored. We define this angle as critical angle $`\varphi _c`$. In Fig. 5 $`\varphi _c`$ is plotted for different temperatures (open circles). Below $`T=0.175t`$ the pseudogap is found everywhere along the Fermi wave vectors except at $`\varphi =45^o`$, thus the Fermi surface is destroyed everywhere except along the diagonal directions. With increasing temperature (up to $`T=0.29t`$), however, $`\varphi _c`$ becomes smaller than $`45^o`$. Thus, the pseudogap region shrinks and at the same time the Fermi surface grows from $`\varphi =45^o`$ up to the critical angle. Above $`T=0.29t`$ the whole Fermi surface is completely restored in spite of some broadening of the spectral function due to interaction as well as finite temperature. This feature is qualitatively consistent with ARPES for underdoped cuprates. We can theoretically calculate the temperature dependence of $`\varphi _c`$. By using the Ornstein-Zernike form of the pairing correlation function and taking the classical fluctuations ($`iq_n=0`$), the scattering rate at the Fermi energy is found to be proportional to $`\mathrm{\Pi }^2(\stackrel{}{k})\xi /\xi _T`$, where $`\xi `$ and $`\xi _T`$ are pairing correlation length and thermal de Broglie wave length defined as $`\xi _T=v_F(\stackrel{}{k})/T`$, respectively. For the pseudogap to disappear, the scattering rates should be much smaller than unity $`\mathrm{\Pi }^2(\stackrel{}{k})\xi /\xi _T1`$, allowing us to define a critical angle $`\xi =v_F(\varphi _c)/T/\mathrm{\Pi }^2(\varphi _c)`$. In Fig. 5 this critical angle is also shown as stars. For the best fit near $`0^o`$ angle, a numerical factor of 1.06 multiplies $`\xi `$. Although a small deviation is found near $`45^o`$, the overall magnitude and shape are in reasonable agreement throughout the whole angle. For a d-wave model, the condition for the appearance of a pseudogap in a given momentum depends not only on the anisotropy of the Fermi velocity (which is the only relevant condition in the attractive Hubbard model) but also more importantly on an angle dependent form factor $`\mathrm{\Pi }^2(\stackrel{}{k})`$.
In Fig. 6 we show the spectral function at the Fermi wave vector and the density of states for different temperatures. As the temperature is increased, the spectral weight starts to fill in the pseudogap and at $`T=0.3t`$ the precursor of the superconducting gap completely disappears as shown in Fig. 6(a). At this temperature, however, the pseudogap still persists in the density of states and at higher temperature ($`T=0.4t`$) it finally disappears as shown in Fig. 6(b). Except close to half-filling, the pseudogap appears at higher temperature in the density of states than in the spectral function. Compared with the s-wave case, the density of states is suppressed linearly near the Fermi energy, a reminiscence of $`N(\omega )\omega `$ in the superconducting state. Both in the spectral function and the density of states, the pseudogap disappears in a manner that the spectral weight fills in the pseudogap instead of closing it with increasing temperature. This feature is also consistent with ARPES for underdoped cuprates. This may suggest that in our approach phase fluctuations (spin-wave type) rather than amplitude fluctuations are mainly responsible for the pseudogap formation, although the present approach includes both. Like in the s-wave case, a pseudogap also appears in the density of states when the characteristic pairing frequency $`\nu _c`$ is equal to or smaller than temperature. This corroborates the origin of the d-wave pseudogap, namely, growing d-wave paring fluctuations in the low temperature renormalized classical regime of the low dimensional system. The calculated $`T^{}`$ follows the same trend as the mean-field critical temperature $`T_{MF}`$. (To be more precise, $`T^{}`$ is approximately half of $`T_{MF}`$ for most of the densities.) As noted in Ref., near a point with high order parameter symmetry (half-filling in that paper), the transition temperature $`T_c`$ decreases while the pseudogap temperature increases along with $`T_{MF}`$. As a result it leads to a large pseudogap regime, consistent with the phase diagram in the underdoped side of cuprates.
There are several advantages in our formulation. First, the pairing fluctuation sum rule Eq. (8) is exactly satisfied (by construction). Through this sum rule, the Mermin-Wagner theorem is formally fulfilled and the strength of pairing fluctuations is properly constrained within the Gaussian approximation. This latter feature is crucial in our formulation, because an approximate treatment of pairing fluctuations without constraining the strength can easily overestimate or underestimate the magnitude of fluctuations particularly in low dimensional systems. The Mermin-Wagner theorem is also satisfied in some other approaches such as the self-consistent T-matrix (FLEX) and the “pairing” approximation schemes. Second, there is an exact relation between one-particle (self-energy, Green’s function) and two-particle (interaction term) functions:
$`\underset{\tau 0^{}}{lim}{\displaystyle \frac{T}{N}}{\displaystyle \underset{\stackrel{}{k},i\omega _n}{}}\mathrm{\Sigma }(\stackrel{}{k},i\omega _n)G(\stackrel{}{k},i\omega _n)e^{i\omega _n\tau }=V\mathrm{\Delta }_i^+\mathrm{\Delta }_i,`$ (15)
where here the self-energy includes the Hartree-Fock term. When the interacting Green’s function $`G(\stackrel{}{k},i\omega _n)`$ is replaced by the noninteracting one $`G^0(\stackrel{}{k},i\omega _n)`$ in Eq. (15), it is exactly satisfied. With $`G(\stackrel{}{k},i\omega _n)`$, the difference between the left- and right-hand sides is less than $`6\%`$ for all temperatures studied. In this paper the Kosterlitz-Thouless phase transition and its fluctuation effect near $`T_{KT}=T_c`$ have been neglected, since the present formulation is inadequate to describe the topological nature of the phase in two dimensions and vortex anti-vortex binding-unbinding physics. Since for similar parameters for the attractive Hubbard model the QMC results indicate that $`T_{KT}0.05t`$ three times smaller than the lowest temperature in our calculations, we do not expect in the present results any significant influence from vortex phase fluctuations. Our approach is valid for weak to intermediate coupling and for temperature not too deep in the pseudogap regime. At very low temperature, even the Ginzburg-Landau functional form itself may change, for instance, a possible crossover of the dynamical critical exponent from z=2 to other value, and eventually vortex phase fluctuations may come into play.
In summary, we have studied spectral properties and pseudogaps in a model with d-wave pairing symmetry by using a non-perturbative sum rule approach. We applied to this model our previous experience of many-body theory in the repulsive and attractive Hubbard models. The magnitude of the normal state pseudogap shows strong angle dependence well fitted by $`\mathrm{cos}2\varphi `$ form at low temperature. With increasing temperature the pseudogap closes at some critical angle $`\varphi _c`$ and beyond this angle a single quasiparticle-like peak appears. The resulting Fermi surface is strongly temperature dependent. With increasing temperature, the pseudogap region shrinks and at the same time the Fermi arc grows from $`\varphi =45^o`$ to $`\varphi _c`$. Both in the spectral function and density of states the pseudogap disappears in a manner that the spectral weight fills in the pseudogap instead of closing it with increasing temperature. All these features are qualitatively consistent with ARPES for underdoped cuprates. The pseudogap is caused by growing d-wave critical pairing fluctuations in the low-temperature classical renormalized regime of the low-dimensional system, as in the repulsive and attractive Hubbard models. We argue that although the real critical behaviors and critical exponents are governed by the vortex phase fluctuations close to the $`T_{KT}`$, the initial growth of pairing fluctuations can be driven by spin-wave phase fluctuations, leading to the normal state pseudogap formation. The behavior of the spin-wave type phase fluctuations belonging to the $`O(n=\mathrm{})`$ universality class can be qualitatively different in two dimensions (particularly near a point with high order parameter symmetry) from that of the $`O(2)`$ Kosterlitz-Thouless vortex phase fluctuations.
The author would like to thank A. M. Tremblay for numerous help and discussions throughout the work. This work was supported by a grant from the Natural Sciences and Engineering Research Council (NSERC) of Canada and the Fonds pour la formation de Chercheurs et d’Aide à la Recherche (FCAR) of the Québec government.
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# Generalised Hilbert Numerators II
## 1. Introduction
The ring $`K[X]`$ of formal polynomials was used by Halter-Koch to study polynomial functions on modules. The author used it to study initial ideals of generic ideals of the same type, for instance generated by a quadratic and a cubic generic form, but in ever more variables . In brief, it is the largest $``$-graded subring of the power series ring $`K[[X]]`$ on a countable set $`X`$ of indeterminates. There are *truncation maps* $`\rho _n:K[X]K[x_1,\mathrm{},x_n]`$, and inclusion maps going the other way, which relate ideals in $`K[X]`$ with sequences $`(I_n)_{n=1}^{\mathrm{}}`$ of ideals, where $`I_n`$ is an ideal in $`K[x_1,\mathrm{},x_n]`$, and where $`I_{n+1}`$ maps to $`I_n`$ under the map
$$K[x_1,\mathrm{},x_{n+1}]\frac{K[x_1,\mathrm{},x_{n+1}]}{(x_{n+1})}K[x_1,\mathrm{},x_n].$$
As an example, choose positive integers $`d_1,\mathrm{},d_r`$, and let for each positive integer $`n`$ $`f_1^{(n)},\mathrm{},f_r^{(n)}`$ be forms in $`n`$ variables of degree $`d_1`$ to $`d_r`$. Suppose furthermore that $`f_{\mathrm{}}^{(n+1)}f_{\mathrm{}}^{(n)}`$ is divisible by $`x_{n+1}`$, and that the coefficients of the forms are choosen randomly. Let $`I_n=(f_1^{(n)},\mathrm{},f_r^{(n)})`$, and let $`>`$ the the lexicographic term order on the various polynomial rings. Then the initial ideals $`\mathrm{in}(I_n)`$ will converge to a monomial ideal in infinitely many indeterminates, as $`n\mathrm{}`$, and this monomial ideal is the lex-initial ideal of $`I=(f_1,\mathrm{},f_r)K[X]`$, where $`f_{\mathrm{}}=limf_{\mathrm{}}^{(n)}`$ . In this case, we have that for $`nr`$, the Hilbert series of $`\frac{K[x_1,\mathrm{},x_n]}{I_n}`$ is $`(1t)^n_{j=1}^r(1t^{d_j})`$.
For arbitrary homogeneous, finitely generated ideals $`IK[X]`$, we conjecture that
$$(1t)^nH(\frac{K[x_1,\mathrm{},x_n]}{\rho _n\left(I\right)};t)$$
is eventually constant, where $`H(\frac{K[x_1,\mathrm{},x_n]}{\rho _n\left(I\right)};t)`$ denotes the Hilbert series of $`\frac{K[x_1,\mathrm{},x_n]}{\rho _n\left(I\right)}`$. What we have in fact proved is that for all homogeneous ideal which are countably generated and which have but finitely many minimal generators of each total degree (we call such an ideal *locally finitely generated*), the polynomials
$$(1t)^nH(\frac{K[x_1,\mathrm{},x_n]}{\rho _n\left(I\right)};t)p(t)[[t]],$$
we call the power series $`p(t)`$ the *generalised Hilbert numerator* of $`I`$. The outstanding question is thus whether for finitely generated ideals, $`p(t)`$ is always a polynomial.
In this article, we first show that for monomial ideals in $`K[X]`$, the polynomial ring on countably many indeterminates, the usual methods of calculating multigraded Hilbert series can be used, and that this Hilbert series can always be written $`\frac{p(X)}{_{i=1}^{\mathrm{}}(1x_i)}`$. We call $`p(X)`$ the $`[X]`$-multigraded Hilbert numerator of the ideal. For a locally finitely generated ideal in $`K[X]`$, we form its initial ideal, contract it to a monomial ideal in $`K[X]`$, calculate the $`[X]`$-multigraded Hilbert numerator of that, then collapse the grading to a $``$-grading to get the $``$-graded Hilbert numerator. Apart from providing a more attractive proof of the existence of Hilbert numerators, this methodology yields immediately $`^r`$-graded Hilbert numerators for $`^r`$-graded locally finitely generated ideals of $`K[X]`$.
We give exact (although opaque) descriptions of the set of all $`[X]`$-graded Hilbert numerators of monomial ideals, and the set of all $``$-graded Hilbert numerators of locally finitely generated ideals of $`K[X]`$. The latter result can be described briefly as follows: the set of *all polynomial* $``$-graded Hilbert numerators is the set of sufficiently high iterated differences of admissible $`H`$-vectors (in the sense of Macaulays characterisation of admissible Hilbert functions of finitely generated algebras), and the set of *all* $``$-graded Hilbert numerators is the closure in $`[[t]]`$ of the previous set.
## 2. Notation
Let $`,,^+`$ denote the set of integers, non-negative integers, and positive integers, respectively. For any set $`A`$ and any positive integer $`k`$, we denote by $`\left(\genfrac{}{}{0pt}{}{A}{k}\right)`$ the set of $`k`$-subsets of $`A`$, by $`[A]`$ the free abelian monoid on $`A`$, and by $`[A]_k`$ the subset of monomials of total degree $`k`$. If $`m[A]_k`$ we write $`|m|=k`$.
If $`B`$ is another set, then $`B^A`$ denotes the set of all functions $`AB`$. If $`A`$ is pointed, that is, has an extinguished zero element (for instance, $`[A]`$ is pointed), then $`B^{(A)}`$ denotes the set of all finitely supported maps $`AB`$.
If $`K`$ is a commutative ring, then $`K^{[A]}`$ becomes a commutative $`K`$-algebra under component-wise addition and multiplication of scalars, and with multiplication given by the *convolution product*
$$f\times g(m)=\underset{t|m}{}f(t)g(m/t).$$
(1)
We denote this ring by $`K[[A]]`$. The set $`B^{([A])}`$ is a subring, which we denote by $`K[A]`$.
## 3. Rings of formal power series and formal polynomials in countably many indeterminates
Let $`K`$ be a field of containing the rational numbers, and let $`X=\{x_1,x_2,x_3,\mathrm{}\}`$ be a set of indeterminates. Form the large power series ring $`K[[X]]`$ and the polynomial ring $`K[X]`$ as above. For $`K[[X]]f=_{m[X]}c_mm`$ we define
$$Supp(f)=\left\{m\text{ }c_m0\right\}$$
(2)
If $`t[X]`$, we define
$$[t]f=c_t.$$
(3)
The ring $`K[X]`$ is $`[X]`$-graded, and in particular $``$-graded, whereas $`K[[X]]`$ is not. The largest $`[X]`$-graded subring of $`K[[X]]`$ is $`K[X]`$, whereas the largest $``$-graded subring is the ring $`K[X]`$ generated by all *bounded elements*: an element $`fK[[X]]`$ is bounded if
$$|f|:=sup(\{|m|\text{ }mSupp(f)\})<\mathrm{}.$$
Another way of putting this is the following.
###### Definition 3.1.
We define the *total-degree* filtration on $`K[[X]]`$ and its various subrings by
$$𝒯^dK[[X]]=\{fK[[X]]\text{ }|f|d\}$$
(4)
For $`K[[X]]f=_{m[X]}c_mm`$, we put
$$𝒯^df=\underset{\begin{array}{c}m[X]\\ |m|d\end{array}}{}c_mm$$
(5)
Then $`K[X]=_{d0}𝒯^dK[[X]]`$.
It is shown in that $`K[X]`$ is also the maximal subring of $`K[[X]]`$ with the following property: given any multiplicative total order $`>`$ on $`[X]`$ whose restriction to $`[X]_1`$ is order-isomorphic to $`\omega `$ (such a $`>`$ will be called a *term order* on $`[X]`$), the support of any non-constant element $`f`$ contains a maximal element $`\mathrm{in}_>(f)`$. Putting $`\mathrm{in}_>(1)=0`$, $`\mathrm{in}_>(0)=\mathrm{}`$, we can regard
$$\mathrm{in}_>:K[X][X]\{\mathrm{}\}$$
as a $`[X]`$-valuation, which induces a $`[X]`$-filtration of $`K[X]`$ by
$$\begin{array}{cc}\hfill 𝔉^{<m}K[X]& =\left\{fK[X]\mathrm{in}_>(f)<m\right\}\hfill \\ \hfill 𝔉^mK[X]& =\left\{fK[X]\mathrm{in}_>(f)m\right\}\hfill \end{array}$$
(6)
We then have a canonical map
$$\begin{array}{cc}\hfill K[X]& \mathrm{gr}(K[X])=\underset{m[X]}{}\frac{𝔉^mK[X]}{𝔉^{<m}K[X]}K[X]\hfill \\ \hfill f& \mathrm{in}_>(f)\hfill \end{array}$$
(7)
This map sends an ideal $`IK[X]`$ to its *initial ideal*
$$\mathrm{in}_>(I)=K[X]\left\{\mathrm{in}_>(f)\text{ }fI\right\}$$
(8)
which is a *monomial ideal*, that is, generated by monomials. We note that
1. Every monomial ideal is its own initial ideal,
2. Extension and contraction of ideals gives a bijection between monomial ideals in $`K[X]`$, $`K[X]`$ and $`K[[X]]`$,
3. Monomial ideals in $`K[X]`$, $`K[X]`$ or $`K[[X]]`$ correspond bijectively to monoid ideals in $`[X]`$.
Because of this identification, we shall say that a monoid ideal has a certain property whenever the corresponding monomial ideal has.
###### Theorem 3.2 (Snellman ).
For a $``$-graded ideal $`IK[X]`$, the following are equivalent:
1. $`I`$ is generated by a *locally finite* set, that is a set
$$F=\underset{d=1}{\overset{\mathrm{}}{}}F_i,i:F_iK[X]_i,i:\mathrm{\#}F_i<\mathrm{}$$
(9)
2. For each positive integer $`d`$,
$$dim_K\left(\frac{I_d}{_{i=1}^{d1}K[X]_iI_{di}}\right)<\mathrm{}$$
(10)
We call such ideals *locally finitely generated* (lfg). By our previous remark, we can talk about lfg monoid ideals, as well.
###### Theorem 3.3 (Snellman ).
Let $`>`$ be a term-order on $`[X]`$, and $`IK[X]`$ a homogeneous ideal. Then $`I`$ is lfg if and only if $`\mathrm{in}_>(I)`$ is.
### 3.1. Inverse limits
We shall need to relate elements in $`K[X]`$ with their *truncations* in $`K[X_n]`$. The necessary machinery is as follows.
For any positive integer $`n`$, we put $`X_n=\{x_1,\mathrm{},x_n\}`$, and let $`[X_n]`$ be the free abelian monoid on $`X_n`$. We define the polynomial ring $`K[X_n]`$ and the power series ring $`K[[X_n]]`$ as above. For $`i<j`$ there is a commutative diagram of $`K`$-multilinear maps
(11)
with the horisontal arrows given by inclusions, and the remaining ones given by (restrictions of) the truncation maps
$$\begin{array}{cc}\hfill \rho _n:[X]& [X_n]\{0\}\hfill \\ \hfill m& \{\begin{array}{cc}m\hfill & m[X_n]\hfill \\ 0\hfill & m[X_n]\hfill \end{array}\hfill \\ \hfill \rho _n:K[[X]]& K[[X_n]]\hfill \\ \hfill \underset{m[X]}{}c_mm& \underset{m[X]}{}c_m\rho _n\left(m\right),\hfill \end{array}$$
(12)
With respect to these inverse systems, we have that $`\underset{}{\mathrm{lim}}K[[X_n]]K[[X]]`$, whereas
$$K[X]K[X]\underset{}{\mathrm{lim}}K[X_n]K[[X]].$$
In fact,
$$\begin{array}{cc}\hfill \underset{}{\mathrm{lim}}K[X_n]& =\{fK[[X]]n:\rho _n\left(f\right)K[X_n]\}\hfill \\ \hfill K[X]& =\left\{f\underset{}{\mathrm{lim}}K[X_n]f\text{ is bounded }\right\}.\hfill \end{array}$$
(13)
Furthermore, the ring $`\underset{}{\mathrm{lim}}K[X_n]`$ is endowed with a natural topology, the inverse limit topology (where all $`K[X_n]`$ are discrete), and the ring $`K[X]`$ is a dense subring. The topology on $`K[X]`$ can be characterised by giving the closure of an arbitrary subset $`AK[X]`$:
$$\overline{A}=\{fK[X]\text{ }n:\rho _n\left(f\right)\rho _n\left(A\right)\}.$$
(14)
It was proved in that with respect to this topology, lfg ideals in $`K[X]`$ are closed. It was also proved that the closed monomial ideals are precisely the lfg monomial ideals.
### 3.2. Topologies on the set of ideals of $`K[X]`$, and a “continuity” result
###### Definition 3.4.
Let $`,𝔠,𝔥,𝔩,𝔪`$ denote the following sets of ideals in $`K[X]`$: all ideals, closed ideals, homogeneous ideals, lfg ideals, monomial ideals. We will also use combinations of letters to denote intersections, for instance
$$𝔩𝔪=𝔩𝔪$$
denotes the lfg monomial ideals.
###### Proposition 3.5.
1. The function
$$d(I,J)=2^n,n=\mathrm{max}\left\{n\text{ }\rho _n\left(I\right)=\rho _n\left(J\right)\right\}$$
(15)
gives a metric on $`𝔠`$.
2. The function
$$\widehat{d}(I,J)=2^d,d=\mathrm{max}\left\{d\text{ }𝒯^dI=𝒯^dJ\right\}$$
(16)
gives a metric on $`𝔥`$.
3. Define a convergence structure on $`𝔪`$ by dictating that $`I_n\stackrel{}{}I𝔪`$ if and only if,
$$m[X]:N(m)^+:n>N(m):mImI_n$$
(17)
Then the corresponding topology is weaker than both the previous topologies.
###### Proof.
1. It is clear the $`d(I,J)=d(J,I)0`$. Since $`I,J`$ are closed, $`d(I,J)=0`$ if and only if $`I=J`$. If $`A,B,C`$ are closed ideals, and if $`d(A,B)2^n`$, $`d(B,C)2^n`$, then $`\rho _n\left(A\right)=\rho _n\left(B\right)=\rho _n\left(C\right)`$, hence $`d(A,C)2^n`$. Thus the triangle inequality holds.
2. Obvious.
3. Let $`m[X_v]_d`$, let $`I,I_1,I_2,I_3,\mathrm{}`$ be monomial ideals, and suppose that either $`d(I_n,I)0`$ or $`\widehat{d}(I_n,I)0`$. In the first case, there is an $`N(v)`$ such that $`\rho _v\left(I_n\right)=\rho _v\left(I\right)`$ whenever $`nN(v)`$: since $`mIm\rho _v\left(I\right)`$, and similarly for $`I_n`$, it follows that $`mI`$ if and only if $`mI_n`$ for $`nN(v)`$. In the second case, there is an $`\widehat{N}(d)`$ such that $`𝒯^dI_n=𝒯^dI`$ whenever $`n\widehat{N}(d)`$: since $`mIm𝒯^dI`$, and similarly for $`I_n`$, it follows that $`mI`$ if and only if $`mI_n`$ for $`n\widehat{N}(d)`$.
###### Theorem 3.6.
Let $`>`$ be a term-order on $`[X]`$. Then the map
$$\begin{array}{cc}\hfill \mathrm{in}_>:𝔩& 𝔩𝔪\hfill \\ \hfill I& \mathrm{in}_>(I)\hfill \end{array}$$
(18)
is continuous with respect to the $`\widehat{d}`$-metric. If $`>`$ is the degree-reverse lexicographic term order, then
$$n:\rho _n\left(\mathrm{in}_>(I)\right)=\mathrm{in}_>(\rho _n\left(I\right))$$
(19)
from which it follows that (18) is continuous with respect to the $`d`$-metric.
###### Proof.
Using the results of , it is straight-forward to show that if $`I,J`$ are lfg ideals such that $`𝒯^dI=𝒯^dJ`$, then for any term-order $`>`$, $`𝒯^d\mathrm{in}_>(I)=𝒯^d\mathrm{in}_>(J)`$. Hence the first result follows.
It is immediate that the identity (19) implies continuity of (18). For all term orders, the LHS of (19) is included in the RHS, so we need to prove that the reverse inclusion holds for the degree-reverse lexicographic term order. Let $`fI`$ be homogeneous of degree $`d`$; then the monomials in $`Supp(f)`$ are ordered as follows: first the ones in $`[X_1]_dSupp(f)`$, if any, then the ones in $`([X_2]_d[X_1]_d)Supp(f)`$, and so on. Let $`m=\mathrm{in}(f)`$, then $`\mathrm{in}(\rho _n\left(f\right))=m`$ for $`n`$ sufficiently large, and $`0`$ otherwise. In the same way, $`\rho _n\left(m\right)=m`$ for sufficiently large $`n`$, and $`0`$ otherwise. So $`\mathrm{RHS}\mathrm{in}(\rho _n\left(f\right))=\rho _n\left(\mathrm{in}(f)\right)\mathrm{LHS}`$. ∎
The following result is a key one: it is what will allow us to define Hilbert numerators of lfg ideals by passing to their initial ideals.
###### Theorem 3.7 (Snellman ).
If $`I`$ is a lfg ideal, and $`>`$ is a term-order on $`[X]`$, then
$$\widehat{d}(\mathrm{in}_>(\rho _n\left(I\right))^e,\mathrm{in}_>(I))0\text{ as }n\mathrm{}.$$
(20)
## 4. Monoid ideals and arithmetic on $`[[X]]`$
### 4.1. Topologies on $`[[X]]`$
Unless otherwise stated, we henceforth assume that $`[[X]]=^{[X]}`$ is given the product topology. With this topology, $`f_nf`$ if for all $`m[X]`$, there is an $`N(m)`$ so that for $`nN(m)`$, $`[m]f=m[f_n]`$. An infinite sum $`_nf_n`$ is convergent if and only if each monomial $`m[X]`$ occurs in but finitely many of the sets $`Supp(f_n)`$.
We can also topologise $`[[X]]`$ by means of the total degree filtration: a sequence $`(f_n)_{n=1}^{\mathrm{}}`$, $`f_nf`$ if and only if
$$d:N(d)^+:v>N(d):m[X]_d:[m]f_v=[m]f.$$
This is a stronger topology than the previous one. We shall use it in particular for the study of the subring $`𝒮`$, to be defined later. For later use, we note the following:
###### Lemma 4.1.
The total degree filtration topology on $`[[X]]`$ gives a linear topology, and hence additiv translation with arbitrary elements, and multiplicative translation with invertible elements, are closed mappings.
###### Proof.
We put
$$J_d=\{0\}\left\{f[[X]]\text{ }Supp(f)_{v=d}[X]_v\right\}.$$
Then the $`J_d`$’s are clopen ideals which form a fundamental system of neighbourhoods of zero.
It follows that for any subset $`A[[X]]`$, the closure is given by
$$\overline{A}=_{d=1}^{\mathrm{}}(A+J_d).$$
Hence if $`h[[X]]`$ and $`A[[X]]`$, then
$$\overline{h+A}=_{d=1}^{\mathrm{}}(h+A+J_d)=h+_{d=1}^{\mathrm{}}(A+J_d)=h+\overline{A},$$
where the crucial inclusion $`_{d=1}^{\mathrm{}}(h+A+J_d)h+_{d=1}^{\mathrm{}}(A+J_d)`$ is proved as follows. If $`fh+A+J_d`$ for all $`d`$, then $`fhA+J_d`$ for all $`d`$, hence $`fh_{d=1}^{\mathrm{}}A+J_d`$, hence $`fh+_{d=1}^{\mathrm{}}A+J_d`$.
If $`h`$ has a multiplicative inverse $`h^1`$, then
$$\overline{hA}=_{d=1}^{\mathrm{}}(hA+J_d)=h_{d=1}^{\mathrm{}}(A+J_d)=h\overline{A};$$
the inclusion $`_{d=1}^{\mathrm{}}(hA+J_d)h_{d=1}^{\mathrm{}}(A+J_d)`$ is proved as follows. Suppose that $`fhA+J_d`$ for all $`d`$, then $`h^1fA+h^1J_dA+J_d`$ for all $`d`$, hence $`h^1f_{d=1}^{\mathrm{}}(A+J_d)`$, hence $`fh_{d=1}^{\mathrm{}}(A+J_d)`$. ∎
### 4.2. The ring of number-theorethic functions
Define $`\mathrm{\Gamma }`$ to be the set of all maps $`^+`$. With component-wise addition and multiplication by scalars, and with the *Dirichlet convolution*
$$fg(n)=\underset{k|n}{}f(k)g(n/k),$$
(21)
$`\mathrm{\Gamma }`$ becomes a commutative ring, often referred to as *the ring of number-theoretic functions* . The well-known isomorphism, given by unique factorisation of integers, between the multiplicative monoid $`(^+,)`$ of the positive integers and a denumerable sum of copies of $`(,+)`$, induces an isomorphism
$$\begin{array}{cc}\hfill \mathrm{\Gamma }& [[X]]\hfill \\ \hfill f& \underset{m=x_1^{\alpha _1}\mathrm{}x_n^{\alpha _n}[X]}{}f(p_1^{\alpha _1}\mathrm{}p_n^{\alpha _n})m,\hfill \end{array}$$
(22)
Define the elements $`\nu ,\mu [[X]]`$ by
$$\begin{array}{cc}\hfill \nu & =\underset{m[X]}{}m\hfill \\ \hfill \mu & =\underset{i=1}{\overset{\mathrm{}}{}}(1x_i)=1\underset{i=1}{\overset{\mathrm{}}{}}x_i+\underset{i<j}{}x_ix_j\underset{i<j<k}{}x_ix_jx_k+\mathrm{}.\hfill \end{array}$$
(23)
Then the image of $`\mu `$ in $`\mathrm{\Gamma }`$ is the well-known Möbius function, and Möbius inversion can be expressed by the formula
$$\nu \mu =1.$$
(24)
We note that we can write
$$\begin{array}{cc}\hfill \nu =1+\underset{i=1}{\overset{\mathrm{}}{}}\nu _i,& \mu =1+\underset{i=1}{\overset{\mathrm{}}{}}\mu _i\hfill \\ \hfill \nu _i=\underset{m[X]_i}{}m,& \mu _i=(1)^i\underset{\sigma \left(\genfrac{}{}{0pt}{}{[X]}{i}\right)}{}\sigma \hfill \end{array}$$
(25)
where $`\nu _i`$ is the $`i`$’th *complete symmetric function* and $`(1)^i\mu _i`$ is the $`i`$’th *elementary symmetric function*.
### 4.3. Characteristic/generating functions of monoid ideals
#### 4.3.1. Definitions
###### Definition 4.2.
If $`I`$ is a monoid ideal in $`[X]`$ then
$$W(I)=I\left([X]\{1\}\right)I$$
(26)
denote the canonical set of minimal generators of $`I`$. We define
$`\mathrm{char}(I)`$ $`={\displaystyle \underset{mI}{}}m`$ (27)
$`w(I)`$ $`={\displaystyle \underset{mW(I)}{}}m`$ (28)
$`q(I)`$ $`=\nu \mathrm{char}(I)`$ (29)
$`p(I)`$ $`=\mu q(I)`$ (30)
We call $`\mathrm{char}(I)`$ the *characteristic function* of $`I`$, $`q(I)`$ the $`[X]`$-graded *Hilbert series* of $`I`$, and $`p(I)`$ the $`[X]`$-graded *Hilbert numerator* of $`I`$.
For a monomial ideal $`J`$ in $`K[X]`$ or $`K[X]`$, we put
$$\begin{array}{cc}\hfill \mathrm{char}(J)& =\mathrm{char}(J[X])\hfill \\ \hfill w(J)& =w(J[X])\hfill \\ \hfill q(J)& =q(J[X])\hfill \\ \hfill p(J)& =p(J[X]).\hfill \end{array}$$
(31)
Similarly, if $`n`$ is a positive integer, and if $`I`$ is a monoid ideal in $`[X_n]`$, then we put
$`\mathrm{char}^n(I)`$ $`={\displaystyle \underset{mI}{}}m`$
$`q^n(I)`$ $`={\displaystyle \underset{m[X_n]I}{}}m`$
$`p^n(I)`$ $`=\rho _n\left(\mu \right)q^n(I).`$
###### Remark 4.3.
$`\mathrm{char}(I)`$ and $`q(I)`$ are the $`[X]`$-graded Hilbert series of $`I`$, regarded as a monomial ideal in $`K[X]`$, and $`\frac{K[X]}{I}`$, respectively. However, the ring $`K[X]`$ is not $`[X]`$-graded, so in order to attach a meaning to $`q(I)`$ for a monomial ideal we regard it as a limit of the Hilbert series of $`\rho _n\left(I\right)`$, that is, as a limit of $`q^n(I)`$.
We note that $`\mathrm{char}(I)`$, $`w(I)`$, $`q(I)`$, and $`p(I)`$ all lie in $`[[X]]`$.
#### 4.3.2. Distributiveness properties
###### Proposition 4.4.
Suppose that $`I,I_1,I_2,I_3,`$ are monomial ideals, and suppose that
$$\underset{i=1}{\overset{\mathrm{}}{}}I_n=I,$$
(32)
and that the sum is convergent with respect to the $`\stackrel{}{}`$ topology. Then
$$\mathrm{char}(I)=\underset{i}{}\mathrm{char}(I_i)\underset{i<j}{}\mathrm{char}(I_iI_j)+\underset{i<j<k}{}\mathrm{char}(I_iI_jI_k)\mathrm{},$$
(33)
and the sum is convergent (with respect to the product topology on $`[[X]]`$).
Putting $`\widehat{p}(I)=p(I)1`$, we also have that
$$\widehat{p}(I)=\underset{i}{}\widehat{p}(I_i)\underset{i<j}{}\widehat{p}(I_iI_j)+\underset{i<j<k}{}\widehat{p}(I_iI_jI_k)\mathrm{},$$
(34)
and this is a convergent sum.
###### Proof.
If we identify monomial ideals with their characteristic functions, and write $``$ for intersections of ideals, and $``$ for sum of ideals, then $``$ and $``$ correspond to component-wise minimum and maximum, and (33) to the identity
$$\underset{i=1}{\overset{\mathrm{}}{}}f_i=\underset{i}{}f_i\underset{i<j}{}f_if_j+\mathrm{},$$
(35)
where the sum is component-wise. The LHS of (35) is always defined; for the RHS to be defined, it is necessary and sufficient that
$$m[X]:N(M):n>N(M):f_i(m)=0.$$
If this holds, then denoting by $`S`$ the cardinality of the finite subset $`\left\{j^+\text{ }f_j(m)0\right\}`$, the formula (35) becomes
$$S\left(\genfrac{}{}{0pt}{}{S}{2}\right)+\left(\genfrac{}{}{0pt}{}{S}{3}\right)\mathrm{}=\{\begin{array}{cc}0\hfill & S=\mathrm{}\hfill \\ 1\hfill & S\mathrm{}\hfill \end{array}$$
a well-know binomial identity.
To prove (34), note that $`\widehat{p}(I_i)=\mu \mathrm{char}(I_i)`$, hence from (35) we get that
$`\widehat{p}(I)`$ $`=\mu \mathrm{char}(I)`$
$`=\mu \left({\displaystyle \underset{i}{}}\mathrm{char}(I_i){\displaystyle \underset{i<j}{}}\mathrm{char}(I_iI_j)+{\displaystyle \underset{i<j<k}{}}\mathrm{char}(I_iI_jI_k)\mathrm{}\right)`$
$`={\displaystyle \underset{i}{}}(\mu \mathrm{char}(I_i)){\displaystyle \underset{i<j}{}}(\mu \mathrm{char}(I_iI_j))+\mathrm{}`$
$`={\displaystyle \underset{i}{}}\widehat{p}(I_i){\displaystyle \underset{i<j}{}}\widehat{p}(I_iI_j)+{\displaystyle \underset{i<j<k}{}}\widehat{p}(I_iI_jI_k)\mathrm{}`$
#### 4.3.3. Inclusion-exclusion for Hilbert numerators
###### Theorem 4.5.
Let $`I[X]`$ be a monoid ideal. If $`\sigma W(I)`$ is finite, let $`\mathrm{𝐥𝐜𝐦}(\sigma )`$ be the least common multiple of the elements in $`\sigma `$, and let $`\mathrm{\#}\sigma `$ be the cardinality of $`\sigma `$. Then
$$p(I)=\underset{\sigma }{}(1)^{(\mathrm{\#}\sigma )}\mathrm{𝐥𝐜𝐦}(\sigma ),$$
(36)
where the sum is over all finite subsets of $`W(I)`$. Alternatively,
$$p(I)=1\underset{mW(I)}{}m+\underset{\sigma \left(\genfrac{}{}{0pt}{}{W(I)}{2}\right)}{}\mathrm{𝐥𝐜𝐦}(\sigma )\underset{\sigma \left(\genfrac{}{}{0pt}{}{W(I)}{3}\right)}{}\mathrm{𝐥𝐜𝐦}(\sigma )+\mathrm{}$$
(37)
###### Proof.
We have that $`I=_{mW(I)}(m)`$, and that $`(m_i)(m_j)=(\mathrm{𝐥𝐜𝐦}(m_i,m_j))`$. Hence the result follows from (34), once we have proved that that $`p((m))=1m`$ for all $`m[X]`$. But $`\mathrm{char}((m))=_{m|t}t`$, hence by Möbius inversion
$$p((m))=1\mu \mathrm{char}((m))=1\mu \left(\underset{m|t}{}t\right)=1\underset{m|t}{}\mu t=1m.$$
#### 4.3.4. Homology methods
###### Lemma 4.6.
Let $`I[X]`$ be a monoid ideal. Then
$$n^+:\rho _n\left(p(I)\right)=p^n(\rho _n\left(I\right))$$
(38)
###### Proof.
$`p^n(\rho _n\left(I\right))`$ $`=\rho _n\left(\mu \right)q^n(\rho _n\left(I\right))`$
$`=\rho _n\left(\mu \right){\displaystyle \underset{m[X_n]\rho _n\left(I\right)}{}}m`$
$`=\rho _n\left(\mu \right)\rho _n\left({\displaystyle \underset{m[X]I}{}}m\right)`$
$`=\rho _n\left(\mu \right)\rho _n\left(q(I)\right)=\rho _n\left(\mu q(I)\right)=\rho _n\left(p(I)\right).`$
Using this lemma, we can immediately extend the various homological methods for getting the multigraded Hilbert series of monoid ideals in $`[X_n]`$ (see ) to work for monoid ideals in $`[X]`$.
We get
###### Theorem 4.7.
Let $`I[X]`$ be a monoid ideal, and let $`m[X]`$. Let $`\mathrm{\Delta }_m=\mathrm{\Delta }_m(I)2^{(^+)}`$ be the following simplicial complex:
$$\sigma =\{\sigma _1,\mathrm{},\sigma _r\}\mathrm{\Delta }_m\frac{m}{_{i=1}^rx_{\sigma _i}}I.$$
(39)
Then $`\mathrm{\Delta }_m`$ is finite, and
$$[m]p(I)=\stackrel{~}{\chi }(\mathrm{\Delta }_m(I))=\underset{F\mathrm{\Delta }_m(I)}{}(1)^{|F|1}=\underset{i=1}{\overset{\mathrm{}}{}}(1)^idimH_i(\mathrm{\Delta }_m,K),$$
(40)
where $`\stackrel{~}{\chi }`$ denotes the the reduced Euler characteristic of an abstract simplicial complex, counting the empty set as a $`1`$-face.
###### Theorem 4.8.
Let $`I[X]`$ be a monoid ideal, let $`W=W(I)`$ be its minimal set of generators, and let $`L_I`$ be the lattice of all finite lcm’s of elements in $`W`$, ordered by divisibility. Let $`\widehat{0}`$ denote the minimal element in $`L_I`$, and let, for $`mL_I`$, $`\mu (\widehat{0},m)`$ denote the value of the Möbius function of the poset $`L_I`$, evaluated on the interval $`[\widehat{0},m]`$. Let $`\mathrm{\Delta }(\widehat{0},m)`$ denote the abstract simplicial complex of all chains in $`(\widehat{0},m)`$. Then we have:
$$m[X]:[m]p(I)=\{\begin{array}{cc}0\hfill & mL_I\hfill \\ \stackrel{~}{\chi }(\mathrm{\Delta }(\widehat{0},m))=\mu (\widehat{0},m)\hfill & mL_I\hfill \end{array}$$
(41)
###### Proof.
It follows from that $`c_m=0`$ for $`mL_I`$, and that $`c_m=\stackrel{~}{\chi }(\mathrm{\Delta }(0,m))`$ for $`mL_I`$. By , $`\stackrel{~}{\chi }(\mathrm{\Delta }(0,m))=\mu (\widehat{0},m)`$ whenever $`mL_I`$. ∎
#### 4.3.5. Classifications
###### Proposition 4.9.
Let $`f=_{m[X]}c_mm[[X]]`$. Then $`fp(𝔪)`$ if and only if the following conditions hold:
1. $`m[X]:_{s|m}c_s\{0,1\}`$,
2. If $`_{s|m}c_s=1`$ and $`t|m`$ then $`_{s|t}c_s=1`$.
###### Proof.
Suppose that $`I`$ is a monoid ideal in $`[X]`$, then $`q(I)=\nu \mathrm{char}(I)`$ is the characteristic function of $`I^c=[X]I`$. This is an *order ideal*, that is, if $`mI^c`$ and $`t|m`$, then $`tI^c`$. It follows that the set of $`q(I)`$’s is the set of $`g=_{m[X]}d_mm[[X]]`$ such that
1. $`m[X]:d_m\{0,1\}`$,
2. If $`d_m=1`$ and $`t|m`$ then $`d_t=1`$.
Since $`p(I)=\mu q(I)`$, the result follows by Möbius inversion. ∎
###### Proposition 4.10.
Let $`f=_{m[X]}c_mm[[X]]`$ be the $`[X]`$-graded Hilbert numerator of a monoid ideal. Let $`m=x_1^{\alpha _1}\mathrm{}x_n^{\alpha _n}`$. Then
$$\mathrm{abs}(c_m)\left(\genfrac{}{}{0pt}{}{n1}{\frac{n1}{2}}\right)$$
(42)
###### Proof.
From Theorem 4.7 we have that $`c_m`$ is the reduced Euler characteristic of some simplicial complex on $`n`$ vertices. Björner and Kalai showed that the absolute value of the reduced Euler characteristic of a simplicial complex on $`n`$ vertices is $`\left(\genfrac{}{}{0pt}{}{n1}{\frac{n1}{2}}\right)`$. ∎
###### Corollary 4.11.
Let $`f=_{m[X]}c_mm[[X]]`$ be the $`[X]`$-graded Hilbert numerator of a monoid ideal. Let $`m=x_1^{\alpha _1}\mathrm{}x_n^{\alpha _n}`$, and let $`r`$ be the number of $`1in`$ such that $`\alpha _i>0`$. Then
$$\mathrm{abs}(c_m)\left(\genfrac{}{}{0pt}{}{r1}{\frac{r1}{2}}\right)$$
(43)
###### Proof.
Let $`\sigma `$ be a permutation of $`X`$. Define $`\sigma (x_1^{a_1}\mathrm{}x_{\mathrm{}}^a_{\mathrm{}})=_{i=1}^{\mathrm{}}x_{\sigma (i)}^{a_i}`$, and $`\sigma (_{m[X]}c_mm)=_{m[X]}c_m\sigma (m)`$. We let $`\sigma `$ act on monoid ideals in $`[X]`$ in the obvious way. Then $`\mu `$ and $`\nu `$ are fix-points for the action of $`\sigma `$ on $`[[X]]`$, and $`\sigma (\mathrm{char}(I))=\mathrm{char}(\sigma (I))`$ for all monoid ideals $`I`$. Hence
$`p(\sigma (I))`$ $`=\mu (\nu \mathrm{char}(\sigma (I)))`$
$`=\mu (\sigma \nu \sigma (\mathrm{char}(I)))`$
$`=\sigma (\mu (\nu \mathrm{char}(I)))`$
$`=\sigma (p(I)).`$
Let $`i_1,\mathrm{},i_r`$ be the support of $`m`$, that is, $`\alpha _{i_1}>0,\mathrm{}\alpha _{i_r}>0`$, and let $`\sigma `$ be a permutation which takes $`i_1`$ to $`1`$, $`i_2`$ to $`2`$, and so on. Then $`\sigma (m)=x_1^{\alpha _{\sigma ^1(1)}}\mathrm{}x_r^{\alpha _{\sigma ^1(r)}}`$, and
$$c_m=[m]f=[\sigma (m)]\sigma (f),$$
hence the result follows by applying Proposition 4.10. ∎
## 5. The subring $`𝒮`$, locally finitely generated ideals, and their generalised Hilbert numerators
For this section, we fix a positive integer $`r`$ and set-partition $`Y`$ of the set of variables: $`X=_{\mathrm{}=1}^rY_{\mathrm{}}`$. There is an associated map $`y:^+\{1,\mathrm{},r\}`$ such that $`x_nY_{y(n)}`$. We denote by $`\mathrm{deg}`$ the associated $`r`$-multi-grading, that is, the monoid homomorphism
$$\begin{array}{cc}\hfill \mathrm{deg}:[X]& ^r\hfill \\ \hfill x_i& e_{y(i)}\hfill \\ \hfill x_1^{\alpha _1}\mathrm{}x_n^{\alpha _n}& \alpha _1\mathrm{deg}(x_1)+\mathrm{}\alpha _n\mathrm{deg}(x_n)\hfill \end{array}$$
(44)
where $`e_1,\mathrm{},e_r`$ are the natural basis elements of $`^r`$. In particular, if $`r=1`$, then $`\mathrm{deg}(m)=|m|`$. Note that, since $`r`$ is finite, $`K[X]`$ is indeed $`^r`$-graded by means of $`\mathrm{deg}`$, even if it is not $`[X]`$-graded. We say that an ideal is $`r`$-homogeneous if it is homogeneous with respect to this grading. Clearly, $`r`$-homogeneous ideals are homogeneous, and all monomial ideals are $`r`$-homogeneous. Furthermore we have:
###### Proposition 5.1.
Let $`I`$ be an ideal of $`K[X]`$. Then the following are equivalent:
1. $`I`$ is $`r`$-homogeneous and lfg,
2. $`I`$ can be generated by $`F=_{\alpha ^r}F_\alpha `$, where each $`F_\alpha `$ is a finite set of $`r`$-homogeneous elements of multi-degree $`\alpha `$.
3. For each $`\alpha ^r`$,
$$dim_K\left(\frac{I_\alpha }{_{\begin{array}{c}\beta +\gamma =\alpha \\ \beta ,\gamma 0\end{array}}K[X]_\beta I_\gamma }\right)<\mathrm{}$$
(45)
###### Proof.
For each total degree $`d`$, there are only finitely many multi-degrees in $`^r`$ of total degree $`d`$. Thus $`(i)`$ and $`(i)`$ are equivalent. The equivalence of $`(ii)`$ and $`(iii)`$ is parallel to Theorem 3.2 and is proved in the same way (see ). ∎
###### Definition 5.2.
Denote by $`𝒮[[X]]`$ the subring consisting of all $`f[[X]]`$ fulfilling the equivalent conditions below:
1. $`f=f_0+f_1+f_2+f_3+\mathrm{}`$ with $`f_i[X]_i`$,
2. $`f(t,t,t,\mathrm{})`$, the substitution of each $`x_i`$ with the new formal indeterminate $`t`$, is defined,
3. $`f=_{\alpha ^r}f_\alpha `$ with $`f_\alpha [X]_\alpha `$,
4. $`f(t_{y(1)},t_{y(2)},t_{y(3)},\mathrm{})`$, the substitution of each $`x_i`$ with the new formal indeterminate $`t_{y(i)}`$, is defined,
Denote the map $`𝒮ff(t_{y(1)},t_{y(2)},t_{y(3)},\mathrm{})[[t_1,\mathrm{},t_r]]`$ by $`=^y`$.
###### Theorem 5.3.
Let $`I[X]`$ be a monoid ideal. Then the following are equivalent:
1. $`p(I)𝒮`$,
2. $`I`$ is lfg,
3. $`w(I)𝒮`$.
###### Proof.
Write
$`w(I)`$ $`=0+w_1+w_2+w_3+\mathrm{},|w_i|=i`$
$`p(I)`$ $`=1+p_1+p_2+p_3+\mathrm{},|p_i|=i.`$
It is immediate that $`I`$ is lfg if and only if $`w(I)𝒮`$, which occurs precisely when every $`w_i`$ is a polynomial. We note that if $`\sigma W(I)`$ has cardinality $`u`$, and the minimal and maximal total degree of elements in $`\sigma `$ is $`c`$ and $`d`$, respectively, then $`c+1|\mathrm{𝐥𝐜𝐦}(\sigma )|ud`$. Clearly, the terms of $`w(I)`$ contributing to $`p_i`$ in (36) have total degree $`i`$.
Hence, if $`I`$ is lfg, so that each $`w_i`$ is a polynomial, then only the various lcm’s of elements in the support of $`w_1,\mathrm{},w_d`$ may contribute to $`p_d`$. The number of elements in the support of $`w_d`$ is thus $`2^{(\mathrm{\#}w_1+\mathrm{}\mathrm{\#}w_d)}<\mathrm{}`$.
Conversely, if $`I`$ is not lfg, suppose that $`w_1,\mathrm{},w_d`$ are polynomials, but that $`w_{d+1}`$ is not. Using (36) we see that $`p_{d+1}`$ receives contribution from a finite number of terms stemming from lcm’s of elements in the support of $`w_1,\mathrm{},w_d`$, and from the non-polynomial $`w_{d+1}`$. Thus $`p_{d+1}`$ is not a polynomial. ∎
###### Corollary 5.4.
Let $`f=_{m[X]}c_mm[[X]]`$. Then $`fp(𝔩𝔪)`$ if and only if the following conditions hold:
1. $`m[X]:_{s|m}c_s\{0,1\}`$ ,
2. If $`_{s|m}c_s=1`$ and $`t|m`$ then $`_{s|t}c_s=1`$,
3. $`f𝒮`$.
###### Proof.
This follows from Proposition 4.9 and Theorem 5.3. ∎
We henceforth regard $`𝒮`$ as a topological ring having the topology given by the total degree filtration. We have that this topology is the same as the one given by any $`r`$-multi degree filtration in the sense that if $`f_nf`$ if for each multi-degree $`\alpha `$, there is an $`N(\alpha )`$ so that $`f_n`$ and $`f`$ agrees in multi-degree $`\alpha `$ whenever $`nN(\alpha )`$: here $`\alpha `$ is with respect to some term-order on $`^r`$ which refines the total-degree partial order. Similarly, the $`\widehat{d}`$-metric on homogeneous ideals gives the same topology as an analogous $`r`$-multigraded metric.
###### Lemma 5.5.
$`𝒮`$ is a closed subset of $`[[X]]`$.
###### Proof.
Suppose that $`f_if`$, where $`f_i𝒮`$. Fix a total degree $`d`$. By the definition of the total degree filtration topology, there exists an $`N`$ such that $`𝒯^df_i=𝒯^df`$ for all $`i>N`$. Since for all $`f_i`$, $`𝒯^df_i`$ is a polynomial, this is true for $`𝒯^f`$, as well. ∎
###### Theorem 5.6.
$`:𝒮[[t_1,\mathrm{},t_r]]`$ is continuous and clopen, when $`Z[[t_1,\mathrm{},t_r]]`$ is given the $`(t_1,\mathrm{},t_r)`$-adic topology.
###### Proof.
We assume for simplicity that $`r=1`$. Suppose that $`f_if`$ in $`𝒮`$. Fix an integer $`d`$, and choose an $`N(d)`$ such that $`f_if\overline{𝔪^d}`$ for $`i>N(d)`$. Thus for $`i>N(d)`$ we have that the $`t^d`$ coefficient of $`(f_i)`$ and $`(f)`$ coincides. This shows that $`(f_i)(f)`$.
To show that this map is clopen, we pick a basic clopen subset $`O_{f,d}=\left\{g𝒮\text{ }𝒯^df=𝒯^dg\right\}`$, where $`d`$ is a positive integer, and $`f𝒮`$. Then $`(O_{f,d})=\left\{h[[t]]\text{ }𝒯^dh=𝒯^d(f)\right\}`$, and this is a basic clopen set of $`[[t]]`$. ∎
###### Lemma 5.7.
The characteristic function is a continuous mapping from the set of lfg monomial ideals, with the $`\widehat{d}`$ metric, to $`𝒮`$. In fact, it is a homeomorphism onto its image.
###### Proof.
Obvious. ∎
###### Lemma 5.8.
The set of characteristic functions of lfg monoid ideals is a closed subset of $`𝒮`$ (and of $`[[X]]`$).
###### Proof.
This follows from the previous Lemma and from Lemma 5.5. ∎
###### Lemma 5.9.
The set $`p(𝔩𝔪)𝒮`$ is closed.
###### Proof.
By the previous Lemma, the set of characteristic functions of lfg monoid ideals is a closed subset of $`[[X]]`$. By Lemma 4.1, the mapping $`f\mu (\nu f)`$ is a closed mapping, hence $`p(𝔩𝔪)`$ is a closed subset of $`[[X]]`$. From Theorem 5.3 we have that $`p(𝔩𝔪)𝒮`$, hence it is closed in there. ∎
We henceforth assume that $`𝔩𝔪`$ have the $`\widehat{d}`$-topology.
###### Proposition 5.10.
Let $`I,I_1,I_2,I_3,`$ be lfg monomial ideals in $`K[X]`$. The following are equivalent:
1. $`\widehat{d}(I_n,I)0`$,
2. $`\mathrm{char}(I_n)\mathrm{char}(I)`$,
3. $`q(I_n)q(I)`$,
4. $`p(I_n)p(I)`$,
5. $`w(I_n)w(I)`$.
Furthermore, if the conditions above are satisfied, then
$$(p(I_n))(p(I)).$$
###### Proof.
By the previous lemma, $`I_nI`$ if and only if $`\mathrm{char}(I_n)\mathrm{char}(I)`$. Since the endomorphism given by multiplication with a fixed element in a topological ring is continuous,
$$\mathrm{char}(I_n)\mathrm{char}(I)q(I_n)q(I)p(I_n)p(I).$$
If $`w(I_n)w(I)`$, then fixing a total degree $`d`$, we get that there exists an $`N(d)`$ such that $`w(I_n)w(I)\widehat{𝔪^d}`$ for $`nN(d)`$. It then follows that $`\mathrm{char}(I_n)\mathrm{char}(I)\widehat{𝔪^d}`$ for $`nN(d)`$. The converse also holds.
The last assertion follows immediately from the fact that $``$ is continuous. ∎
We now recall a theorem by Macaulay, which says that if $`IK[X_n]`$ is a homogeneous ideal, and $`>`$ is a term-order on $`[X_n]`$, then $`\frac{K[X_n]}{I}`$ and $`\frac{K[X_n]}{\mathrm{in}_>(I)}`$ have the same $``$-graded Hilbert series (see for instance ). It is also true that if $`I`$ is $`r`$-multigraded, when $`K[X_n]`$ is $`^r`$-graded using the partition $`YX_n`$, then the above algebras have in fact the same $`^r`$-graded Hilbert series. Using this, and our previous results, we get:
###### Theorem 5.11.
Suppose that $`>`$ is a term-order on $`[X]`$. Let $`IK[X]`$ be a $`r`$-homogeneous lfg ideal, and define $`g_n[t_1,\mathrm{},t_r]`$ by requiring that
$$\frac{g_n}{_{i=1}^n(1t_{y(i)})}$$
is the $`^r`$-graded Hilbert series of
$$\frac{K[X_n]}{\rho _n\left(I\right)}.$$
Then, $`g_n(p(I))`$ as $`n\mathrm{}`$, and $`(p(I))[[t_1,\mathrm{},t_r]]`$.
###### Proof.
From Theorem 3.7 we know that
$$\widehat{d}(\mathrm{in}_>(\rho _n\left(I\right))^e,\mathrm{in}_>(I))0.$$
Then Proposition 5.10 gives that
$$p(\mathrm{in}_>(\rho _n\left(I\right))^e)p(\mathrm{in}_>(I)),$$
hence, using Theorem 5.6, we have that
$$(p(\mathrm{in}_>(\rho _n\left(I\right))^e))(p(\mathrm{in}_>(I))).$$
It is clear that
$$p^n(\mathrm{in}_>(\rho _n\left(I\right)))=p(\mathrm{in}_>(\rho _n\left(I\right))^e).$$
As we remarked above, a $`r`$-homogeneous ideal have the same $`^r`$-graded Hilbert series as its initial ideal, so
$$g_n=(p^n(\mathrm{in}_>(\rho _n\left(I\right)))),$$
hence
$$g_n(p(\mathrm{in}_>(I))).$$
In , the result above (for $`r=1`$) was proved through a different route, and the power series $`(p(I))`$ was called the *generalised Hilbert numerator* of $`I`$.
We note two simple corollaries:
###### Corollary 5.12.
If $`I,J`$ are $`r`$-homogeneous lfg ideals of $`K[X]`$, and if $`\rho _n\left(I\right)`$ and $`\rho _n\left(J\right)`$ have the same $`^r`$-graded Hilbert series for all $`n`$, then $`I`$ and $`J`$ have the same $`r`$-graded Hilbert numerator.
###### Corollary 5.13.
If $`I`$ is an $`r`$-homogeneous lfg ideal of $`K[X]`$, and if $`>`$ is a term-order on $`[X]`$, then $`I`$ and $`\mathrm{in}_>(I)`$ have the same $`r`$-graded Hilbert numerator.
In particular, $`(p(𝔩))=(p(𝔩𝔪))`$, that is, all $`r`$-graded Hilbert numerators of lfg ideals can be obtained from lfg monomial ideals.
### 5.1. Polynomial $`r`$-graded Hilbert numerators
###### Definition 5.14.
We put $`𝒯^Y=^1([t_1,\mathrm{},t_r])`$. If $`p(I)𝒯^Y`$ we say that $`I`$ has *polynomial $`r`$-graded Hilbert numerator*.
###### Lemma 5.15.
Suppose that $`Y^{}`$ refines $`Y`$. Denote by $`0`$ the partition $`X=X`$. Then $`[X]𝒯^Y^{}𝒯^Y𝒯^0𝒮.`$
###### Proof.
The inclusions are obvious. To see that the strict ones are indeed strict, consider the following examples: $`_{i=1}^{\mathrm{}}(x_1^ix_2^i)𝒯^0[X],`$ $`_{i=1}^{\mathrm{}}x_1^i𝒮𝒯^Y.`$
###### Example 5.16.
There are lfg monomial ideals which have Hilbert numerators in $`𝒮𝒯^0`$. Let $`I`$ be generated by $`a_i=x_1x_2x_3\mathrm{}x_{i1}x_i^2`$, for $`i1`$, and $`b_j=x_1x_2x_3\mathrm{}x_{j2}x_j^6`$, for $`j2`$. Put $`p_n=p^n(\rho _n\left(I\right))`$ for $`n>0`$, and define $`p_0=1`$. We claim that
$$p_n=p_{n1}+(1)^nv_n,v_n=(x_{n1}+x_n^4)x_1x_2\mathrm{}x_{n2}x_n^2\underset{i=1}{\overset{n1}{}}(x_i1)$$
(46)
To see this, we first note that $`(x_{n1}+x_n^4)x_1x_2\mathrm{}x_{n2}x_n^2=a_n+b_n.`$ By induction, we have that $`(x_{n1}+x_n^4)x_1x_2\mathrm{}x_{n2}x_n^2_{i=1}^{n1}(x_i1)`$ consists of those monomials which can be formed as a lcm of $`\{a_n\}S`$ or of $`\{b_n\}S`$ or of $`\{a_n,b_n\}S`$, where $`S\{a_1,a_2,\mathrm{},a_{n1},b_2,b_3,\mathrm{},b_{n1}\}`$. Note that monomials $`x_1^{\alpha _1}\mathrm{}x_n^{\alpha _n}`$ with $`\alpha _i=\alpha _j=6`$ for $`i<j`$ does not occur. This can be readily explained: every such monomial can be expressed as a lcm in two different ways, by either including or omitting the superfluous generator $`a_i`$.
Hence, it follows that
$$\begin{array}{cc}\hfill p_n& =1+\underset{i=1}{\overset{n}{}}(1)^iv_i\hfill \\ \hfill limp_n=p& =1+\underset{i=1}{\overset{\mathrm{}}{}}(1)^iv_i\hfill \end{array}$$
(47)
We have that $`p𝒮[X]`$. As we shall see, $`p𝒯`$.
Setting each $`x_i=t`$ in (47) we have that
$$\begin{array}{cc}\hfill p(t)& =1t^2+\underset{n=2}{\overset{\mathrm{}}{}}(1)^n(t1)^{n1}(t+t^4)t^n\hfill \\ & =1t^2+\frac{t+t^4}{t1}\underset{n=2}{\overset{\mathrm{}}{}}(1)^n(t1)^nt^n\hfill \\ & =1t^2+\frac{t+t^4}{t1}\left(\frac{1}{1+t(t1)}1+(t1)t\right)\hfill \\ & =1t^2t^3+t^5\hfill \end{array}$$
(48)
###### Lemma 5.17.
Let $`r=1`$. Let $`d,a_1,\mathrm{},a_d`$ be integers, with $`d>0`$. Then the set of all $`|p(I)|`$, where $`I`$ is finitely generated and generated in degrees $`d`$, and has
$$(p(I))=1+a_1+\mathrm{}a_dt^d+O(t^{d+1}),$$
(49)
is either empty, or has a maximum, which we denote by $`Q_d(a_1,\mathrm{},a_d)`$.
###### Proof.
We claim that there are positive integers $`u_1,\mathrm{},u_d`$ such that if $`I`$ is a finitely generated monomial ideal generated in degrees $`d`$ satisfying (49), then $`(w(I))=w_1t+\mathrm{}w_dt^d`$ with $`w_iu_i`$ for $`1id`$. Assuming the claim, it is clear that the total degree of $`p(I)`$ is $`_{i=1}^diu_i`$, since this is a bound of the lcm of all the generators.
To establish the claim, we note that $`a_1=w_1`$, and assume by induction that we have shown that $`u_1,\mathrm{},u_i`$ exist. We note that the minimal generators which affect $`a_{i+1}`$ are those of degree $`i+1`$, which each contribute with $`1`$, and also $`s`$-tuples $`m_1,\mathrm{},m_s`$ with $`m_{\mathrm{}}W(I)`$, $`|m_{\mathrm{}}|<i+1`$, and with $`|\mathrm{𝐥𝐜𝐦}(m_1,\mathrm{},m_s)|=i+1`$, which each contribute $`(1)^s`$. If we pick $`\lambda _1`$ elements of $`W(I)_1`$, $`\lambda _2`$ elements of $`W(I)_2`$, et cetera, then for the resulting lcm to be of total degree $`i+1`$ it is necessary that $`\lambda _1+\lambda _2+\mathrm{}+\lambda _ii+1`$ and that $`\lambda _{\mathrm{}}<i+1`$ for all $`\mathrm{}`$; thus only finitely many $`\lambda =(\lambda _1,\mathrm{},\lambda _i)`$ are relevant.
We thus have that
$$a_{i+1}=w_{i+1}+\underset{\lambda }{}(1)^{c(\lambda )}R_\lambda ,$$
(50)
where $`\lambda =(\lambda _1,\mathrm{},\lambda _i)`$, $`|\lambda |i+1`$, $`0\lambda _{\mathrm{}}<i+1`$ for $`1\mathrm{}ji`$, $`c(\lambda )`$ is the number of non-zero entries in $`\lambda `$. The symbol $`R_\lambda `$ denotes a finite interval $`[0,L]`$ of integers, where $`L`$ is the maximal numbers of lcm’s of $`\lambda _1`$ elements of degree 1, drawn from a set of cardinality $`u_1`$, $`\lambda _2`$ elements of degree 2, drawn from a set of cardinality $`u_2`$, and so on, which have total degree $`i+1`$. For instance, if $`i=1`$ and $`\lambda =(2)`$ then $`L=\left(\genfrac{}{}{0pt}{}{u_1}{2}\right)`$, if $`i=2`$ and $`\lambda =(1,1)`$ then $`L=u_1u_2`$.
These finite intervals are added using interval arithmetic, so that $`[a,b]+[c,d]=[a+c,b+d]`$. We can the deduce that
$$w_{i+1}=a_{i+1}+\underset{\lambda }{}(1)^{c(\lambda )}R_\lambda =[C,D]$$
(51)
for some integers $`C,D`$. Putting $`u_{i+1}=D`$ we have the desired bound. ∎
###### Lemma 5.18.
Let $`r=1`$ and suppose that $`f(t)(p(𝔩))`$, in other words, that $`f(t)`$ is the $``$-graded Hilbert numerator of some lfg ideal. Suppose that $`f(t)=1+a_1t+\mathrm{}+a_dt^d+t^{d+r+1}g(t)`$, with $`r>Q_d(a_1,\mathrm{},a_d)`$. Then $`1+a_1t+\mathrm{}+a_dt^d(p(𝔩))`$.
###### Proof.
Let $`f=(p(I))`$, where $`I`$ is a lfg monomial ideal. Let $`I_d`$ denote the ideal generated by everything in $`I`$ of total degree $`d`$. Since the maximal degree of a lcm of the generators of degree $`d`$ is $`Q_d(a_1,\mathrm{},a_r)`$, it follows from (36) and Lemma 5.17 that $`(p(I_d)=1+a_1t+\mathrm{}+a_dt^d`$. ∎
###### Theorem 5.19.
Let $`r=1`$. If $`I[X]`$ is a lfg monoid ideal with $`p(I)𝒯`$, then there exists a positive integer $`N`$ and a monoid ideal $`J[X_N]`$ so that $`(p(J^e))=(p(I))`$.
###### Proof.
Let $`f=(p(I))=1+a_1t+a_2t^2+\mathrm{}a_dt^d`$, and let $`f_n=(p^n(\rho _n\left(I\right)))`$. Then $`f_nf`$ in $`[[t]]`$, with respect to the $`(t)`$-adic topology. Let $`r>Q_d(a_1,\mathrm{},a_d)`$, and choose $`N`$ such that for $`nN`$, $`f_nf(t^r)`$. Then Lemma 5.18 shows that there is a monoid ideal $`J`$ in $`[X_n]`$ with $`f`$ as its $``$-graded Hilbert numerator. ∎
###### Corollary 5.20.
The set of polynomial $``$-graded Hilbert numerators of lfg ideals in $`K[X]`$ is equal to the set of $``$-graded Hilbert numerators of homogeneous ideals in finitely many variables. This set is dense in the set of all possible $``$-graded Hilbert numerators of lfg ideals.
###### Proof.
From Theorem 5.19 we get that all polynomial $``$-graded Hilbert numerators can be obtained from ideals generated in finitely many variables. To prove the second assertion, we note that if $`I`$ is lfg, $`d>0`$, and $`I_d`$ is the ideal generated by everything in $`I`$ of degree $`d`$, then $`(p(I))(p(I_d))mod(t^{d+1})`$, and since $`I_d`$ is finitely generated, $`p(I)[X]`$ hence $`(p(I))[t]`$. ∎
###### Theorem 5.21.
Let, for every pair of integers $`0<ab`$, $`G_{a,b}`$ denote the set
$$\{(1t)^b(1+a_1t+a_2t^2+a_3t^3+\mathrm{})\text{ }a_1=a,i:0a_{i+1}a_i^{<i>}\},$$
where
$$u^{<d>}=\left(\genfrac{}{}{0pt}{}{k(d)+1}{d+1}\right)+\left(\genfrac{}{}{0pt}{}{k(d1)+1}{d}\right)+\mathrm{}+\left(\genfrac{}{}{0pt}{}{k(1)+1}{2}\right)$$
when $`u`$ has $`d`$-th Macaulay expansion
$$u=\left(\genfrac{}{}{0pt}{}{k(d)}{d}\right)+\left(\genfrac{}{}{0pt}{}{k(d1)}{d1}\right)+\mathrm{}+\left(\genfrac{}{}{0pt}{}{k(1)}{1}\right)$$
(see ). Then the set of polynomial $``$-graded Hilbert numerators is $`_{0<ab}G_{a,b}`$, and the closure of this set in $`[[t]]`$ is exactly the set of $``$-graded Hilbert numerators of lfg ideals in $`K[X]`$.
###### Proof.
It follows from a well-know classification by Macaulay (see \[4, Theorem 4.2.10\]) that the set of (generating functions of) Hilbert functions of homogeneous quotients of polynomial rings with finitely many indeterminates is
$$\{\mathrm{\hspace{0.17em}1}+a_1t+a_2t^2+a_3t^3+\mathrm{}\text{ }i:0a_{i+1}a_i^{<i>}\}.$$
The function $`1+a_1t+a_2t^2+a_3t^3+\mathrm{}`$ can be realised as the Hilbert function of a quotient of $`K[x_1,\mathrm{},x_{a_1}]`$ with a monomial ideal; we are of course free to use more variables, if we so desire. The first part of the theorem is therefore demonstrated.
To prove the second part, we proceeds as follows. We know by Lemma 5.9 that
$$p(𝔩𝔪)=\left\{\mu (\nu \mathrm{char}(I))\text{ }I𝔩𝔪\right\}$$
is a closed subset of $`𝒮`$. We have that $`:𝒮[[t]]`$ is a closed map (Theorem 5.6), hence $`(p(𝔩𝔪))`$ is closed in $`[[t]]`$.
Now, Corollary 5.20 shows that the set of polynomial $``$-graded Hilbert numerators is dense in the set of all $``$-graded Hilbert numerators; since this latter set is closed, the second part of the theorem follows. ∎
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# Fourier Tauberian Theorems and Applications
## 1. Tauberian theorems I: basic estimates
Let $`F`$ be a non-decreasing function on $``$. For the sake of definiteness, we shall always be assuming that
(1.1)
$$F(\tau )=\frac{1}{2}\left[F(\tau +0)+F(\tau 0)\right],\tau .$$
### 1.1. Auxiliary functions
We shall deal with continuous functions $`\rho `$ on $``$ satisfying the following conditions:
1. $`|\rho (\tau )|\mathrm{const}\tau ^{2m2}`$, where $`m>\frac{1}{2}`$;
2. $`c_{\rho ,0}:=\rho (\tau )d\tau =1`$;
3. $`\rho `$ is even;
4. $`\rho 0`$;
5. $`\mathrm{supp}\widehat{\rho }[1,1]`$.
For every $`m`$ the functions $`\rho `$ satisfying (1<sub>m</sub>)–(5) do exist (see, for example, \[H2\], Section 17.5, or Example 1.1 below).
###### Example 1.1.
Let $`l`$ be a positive integer and
(1.2)
$$\begin{array}{c}\gamma (\tau ):=_{\frac{\pi }{2}}^{\frac{\pi }{2}}(\frac{\tau }{2l}+s)^{2l}\mathrm{sin}^{2l}(\frac{\tau }{2l}+s)ds.\end{array}$$
The function $`\gamma `$ satisfies (3), (5), and
(1.3)
$$c_\gamma ^{}\tau ^{2l}\gamma (\tau )c_\gamma ^+\tau ^{2l}$$
with some positive constants $`c_\gamma ^\pm `$. Indeed, (3) and (1.3) are obvious, and (5) follows from the fact that $`2(2\pi )^{1/2}\tau ^1\mathrm{sin}\tau `$ is the Fourier transform of the characteristic function of the interval $`[1,1]`$. If $`\rho (\tau ):=c_{\gamma ,0}^1\gamma (\tau )`$ then the conditions (2)–(5) are fulfilled and (1<sub>m</sub>) holds with $`m=l1`$.
We shall always be assuming (1<sub>m</sub>). Let
$$\rho _{1,1}(\tau ):=\{\begin{array}{cc}_\tau ^{\mathrm{}}\rho (\mu )d\mu ,\hfill & \tau >0,\hfill \\ 0,\hfill & \tau =0,\hfill \\ _{\mathrm{}}^\tau \rho (\mu )d\mu ,\hfill & \tau <0,\hfill \end{array}$$
and, if (1<sub>m</sub>) holds with $`m>0`$,
$$\rho _{1,0}(\tau ):=_\tau ^{\mathrm{}}\mu \rho (\mu )d\mu ,\rho _{1,2}(\tau ):=\{\begin{array}{cc}_\tau ^{\mathrm{}}_\mu ^{\mathrm{}}\rho (\lambda )d\lambda d\mu ,\hfill & \tau 0,\hfill \\ _{\mathrm{}}^\tau _{\mathrm{}}^\mu \rho (\lambda )d\lambda d\mu ,\hfill & \tau 0.\hfill \end{array}$$
One can easily see that
$$\rho _{1,0}(\tau )\mathrm{const}\tau ^{2m},\rho _{1,1}(\tau )\mathrm{const}\tau ^{2m1},\rho _{1,2}(\tau )\mathrm{const}\tau ^{2m}$$
for all $`\tau 0`$. Integrating by parts, we obtain
(1.4)
$$\rho _{1,0}(\tau )=_\tau ^{\mathrm{}}\mu \rho _{1,1}^{}(\mu )d\mu =\rho _{1,2}(\tau )+\tau \rho _{1,1}(\tau ),\tau 0.$$
Denote
$$c_{\rho ,\kappa }:=|\mu |^\kappa \rho (\mu )d\mu ,\kappa (1,2m+1).$$
Under condition (2), by Jensen’s inequality, we have
(1.5)
$$c_{\rho ,r}^\kappa c_{\rho ,\kappa }^r,\kappa r0.$$
If the condition (3) is fulfilled then $`\rho _{1,0}`$ and $`\rho _{1,2}`$ are even continuous functions, $`\rho _{1,1}`$ is an odd function continuous outside the origin and
(1.6)
$$\begin{array}{c}\rho _{1,1}(\pm 0)=\pm \frac{1}{2}c_{\rho ,0},\rho _{1,0}(0)=\rho _{1,2}(0)=\frac{1}{2}c_{\rho ,1}.\end{array}$$
Indeed, the first two equalities in (1.6) are obvious, and the last follows from (1.4).
The condition (4) and (1.4) imply that
(1.7)
$$\begin{array}{cc}0\rho _{1,2}(\tau )\rho _{1,0}(\tau ),\hfill & \hfill \tau 0,\\ 0\rho _{1,k}(\mu )\rho _{1,k}(\tau ),\hfill & \hfill k=0,1,2,\mu \tau 0.\end{array}$$
Let
(1.8)
$$\rho _\delta (\tau ):=\delta \rho (\delta \tau ),\rho _{\delta ,k}(\tau ):=\delta ^{1k}\rho _{1,k}(\delta \tau ),k=0,1,2,$$
where $`\delta `$ is an arbitrary positive number. If (5) is fulfilled then
(1.9)
$$\mathrm{supp}\widehat{\rho }_{\delta ,0}\mathrm{supp}\widehat{\rho }_\delta [\delta ,\delta ].$$
Indeed, these inclusions follow from (1.8) and the fact that $`\rho _{1,0}`$ is the convolution of the functions $`\mu \rho (\mu )`$ and $`\chi _{}(\mu )`$.
### 1.2. Main estimates
If $`f`$ is a piecewise continuous function on $`^1`$, we denote
$`fF(\tau )`$ $`:=`$ $`\underset{R\mathrm{}}{lim}{\displaystyle _R^R}f(\tau \mu )F(\mu )d\mu ,`$
$`fF^{}(\tau )`$ $`:=`$ $`\underset{R\mathrm{}}{lim}{\displaystyle _{(R,R)}}f(\tau \mu )dF(\mu ),`$
whenever the limits exist. We shall deduce the estimates for $`F(\tau )`$ from the following simple lemma.
###### Lemma 1.2.
Let $`\rho `$ satisfy the conditions (1<sub>m</sub>)–(3) and $`\rho _{T,1}(\tau s)F(s)0`$ as $`s\pm \mathrm{}`$ for some $`T>0`$ and $`\tau `$. Then $`\rho _{T,1}F^{}(\tau )`$ is well defined if and only if $`\rho _TF(\tau )`$ is well defined, and
(1.10)
$$F(\tau )\rho _TF(\tau )=\rho _{T,1}F^{}(\tau ).$$
###### Proof.
Integrating by parts, we obtain
$$\begin{array}{c}_{(R,R)}\rho _{T,1}(\tau \mu )dF(\mu )=_{(R,\tau )}\rho _{T,1}(\tau \mu )dF(\mu )+_{(\tau ,R)}\rho _{T,1}(\tau \mu )dF(\mu )\hfill \\ \hfill =_R^R\rho _T(\tau \mu )F(\mu )d\mu +\rho _{T,1}(+0)F(\tau 0)\rho _{T,1}(0)F(\tau +0)\\ \hfill \rho _{T,1}(\tau +R)F(R+0)+\rho _{T,1}(\tau R)F(R0).\end{array}$$
In view of (1.1), (1.6) and (2), we have
$$\rho _{T,1}(+0)F(\tau 0)\rho _{T,1}(0)F(\tau +0)=F(\tau ).$$
Now the lemma is proved by passing to the limit as $`R\mathrm{}`$. ∎
###### Theorem 1.3.
Let $`\rho `$ satisfy the conditions (1<sub>m</sub>)–(4) with $`m>0`$. Assume that $`\rho _{\delta ,0}(\tau s)F(s)0`$ as $`s\pm \mathrm{}`$ and $`\rho _{\delta ,0}F^{}(\tau )<\mathrm{}`$ for some $`\delta >0`$ and $`\tau `$. Then $`\rho _TF(\tau )<\mathrm{}`$ and
(1.11)
$$|F(\tau )\rho _TF(\tau )|c_{\rho ,1}^1\delta ^1\rho _{\delta ,0}F^{}(\tau )$$
for all $`T\delta `$.
###### Proof.
The identity (1.4) and (4) imply that
$$\frac{\mathrm{d}}{\mathrm{d}\tau }\left(\frac{\rho _{1,1}(\tau )}{\rho _{1,0}(\tau )}\right)=\frac{\rho (\tau )\left(\tau \rho _{1,1}(\tau )\rho _{1,0}(\tau )\right)}{\left(\rho _{1,0}(\tau )\right)^2}0,\tau >0.$$
Therefore, in view of (2) and (1.6),
$$\frac{|\rho _{1,1}(\tau )|}{\rho _{1,0}(\tau )}\frac{|\rho _{1,1}(+0)|}{\rho _{1,0}(0)}=c_{\rho ,1}^1,\tau >0.$$
Taking into account (3), (1.8) and the second inequality (1.7), we obtain
(1.12)
$$|\rho _{T,1}(\tau )|\frac{\rho _{T,0}(\tau )}{c_{\rho ,1}T}\frac{\rho _{\delta ,0}(\tau )}{c_{\rho ,1}\delta },T\delta >0,\tau .$$
The inequality (1.12) implies that $`\rho `$ and $`F`$ satisfy the conditions of Lemma 1.2 and that $`\rho _TF(\tau )<\mathrm{}`$. Obviously, (1.11) follows from (1.10) and (1.12). ∎
###### Remark 1.4.
If $`T=\delta `$ then the estimate (1.11) can be rewritten in the form
(1.13)
$$\rho _\delta ^+F(\tau )F(\tau )\rho _\delta ^{}F(\tau ),$$
where $`\rho _\delta ^\pm (\tau ):=\rho _\delta (\tau )\pm c_{\rho ,1}^1\delta \tau \rho _\delta (\tau )`$.
###### Remark 1.5.
The inequality (1.11) is not precise in the sense that, apart from some degenerate situations, it never turns into an equality. The crucial point in our proof is the estimate $`|\rho _{T,1}|c_{\rho ,1}^1\delta ^1\rho _{\delta ,0}`$ which implies that $`|\rho _{T,1}F^{}|c_{\rho ,1}^1\delta ^1\rho _{\delta ,0}F^{}(\tau )`$. However, the function $`\rho _{T,1}`$ is negative on one half-line and positive on another, so $`|\rho _{T,1}F^{}|`$ may well admit much a better estimate. Using this observation, one can try to improve our results under additional conditions on the function $`F`$.
###### Theorem 1.6.
Let $`[a,b]`$ be a bounded interval. Assume that the conditions of Theorem 1.3 are fulfilled for every $`\tau [a,b]`$ and that $`\rho _{\delta ,0}F^{}(\tau )`$ is uniformly bounded on $`[a,b]`$. Then
$`T^1\delta ^1f(b)\rho _{\delta ,0}F^{}(b)`$
(1.14) $`{\displaystyle _a^b}f(\tau )\left[F(\tau )\rho _TF(\tau )\right]d\tau `$
$`T^1\delta ^1f(a)\rho _{\delta ,0}F^{}(a)+T^1\delta ^1{\displaystyle _a^b}f^{}(\tau )\rho _{\delta ,0}F^{}(\tau )d\tau `$
for every non-negative non-decreasing function $`fC^1[a,b]`$ and all $`T\delta `$.
###### Proof.
In view of (1.7) and (1.8) we have
(1.15)
$$T\rho _{T,2}(\tau )T^1\rho _{T,0}(\tau )\delta ^1\rho _{\delta ,0}(\tau ),T\delta >0,\tau .$$
This estimates, (1.12) and Lemma 1.2 imply that the functions $`\rho _{T,2}F^{}(\tau )`$, $`|\rho _{T,1}|F^{}(\tau )`$ and $`\rho _TF(\tau )`$ are uniformly bounded on $`[a,b]`$. Since $`\rho _{T,2}^{}(s)=\rho _{T,1}(s)`$ whenever $`s0`$, integrating by parts with respect to $`\tau `$ we obtain
$$\begin{array}{c}_a^bf(\tau )\rho _{T,1}(\tau \mu )dF(\mu )d\tau =f(a)\rho _{T,2}(a\mu )dF(\mu )\hfill \\ \hfill f(b)\rho _{T,2}(b\mu )dF(\mu )+_a^bf^{}(\tau )\left(\rho _{T,2}(\tau \mu )dF(\mu )\right)d\tau .\end{array}$$
Now (1.14) follows from Lemma 1.2 and (1.15). ∎
If $`f1`$ then (1.14) turns into
(1.16)
$$\rho _{\delta ,0}F^{}(b)T\delta _a^b\left[F(\mu )\rho _TF(\mu )\right]d\mu \rho _{\delta ,0}F^{}(a).$$
This estimate and the obvious inequalities
(1.17)
$$\epsilon ^1_{\tau \epsilon }^\tau F(\mu )d\mu F(\tau )\epsilon ^1_\tau ^{\tau +\epsilon }F(\mu )d\mu ,\epsilon >0,$$
imply the following corollary.
###### Corollary 1.7.
Under conditions of Theorem 1.6
(1.18) $`F(b)`$ $``$ $`\epsilon ^1{\displaystyle _{b\epsilon }^b}\rho _TF(\mu )d\mu \epsilon ^1T^1\delta ^1\rho _{\delta ,0}F^{}(b),`$
(1.19) $`F(a)`$ $``$ $`\epsilon ^1{\displaystyle _a^{a+\epsilon }}\rho _TF(\mu )d\mu +\epsilon ^1T^1\delta ^1\rho _{\delta ,0}F^{}(a)`$
for all $`\epsilon (0,ba]`$ and $`T\delta `$.
If (4) is fulfilled then $`\rho _TF`$ is a non-decreasing function. Therefore (1.18) and (1.19) imply that
(1.20) $`F(b)`$ $``$ $`\rho _TF(b\epsilon )\epsilon ^1T^1\delta ^1\rho _{\delta ,0}F^{}(b),`$
(1.21) $`F(a)`$ $``$ $`\rho _TF(a+\epsilon )+\epsilon ^1T^1\delta ^1\rho _{\delta ,0}F^{}(a).`$
###### Remark 1.8.
It is clear from the proof that Theorems 1.3 and 1.6 remain valid (with some other constants independent of $`\delta `$ and $`T`$) if we drop the condition (4) and replace $`\rho _{\delta ,0}(\tau )`$ with an arbitrary non-negative function $`\gamma _\delta `$ such that $`|\rho _{T,1}(\tau )|\mathrm{const}\delta ^1\gamma _\delta (\tau )`$ and $`|\rho _{T,2}(\tau )|\mathrm{const}T^1\delta ^1\gamma _\delta (\tau )`$. In particular, one can take $`\gamma _\delta (\tau )=\delta \gamma (\delta \tau )`$, where $`\gamma `$ is the function defined by (1.3) with $`l=m`$.
## 2. Tauberian theorems II: applications
### 2.1. General remarks
From now on we shall be assuming that the function $`F`$ is polynomially bounded. Then the conditions of Theorems 1.3 and 1.6 are fulfilled for all $`\tau ,a,b^1`$ and $`T\delta >0`$ whenever $`\rho `$ satisfies (1<sub>m</sub>) with a sufficiently large $`m`$.
So far we have not used the condition (5), which is not needed to prove the estimates. However, this condition often appears in applications. It implies that the convolutions $`\rho _TF`$ and $`\rho _{T,0}F^{}`$ are determined by the restrictions of $`\widehat{F}`$ to the interval $`(T,T)`$. If
(2.1)
$$\widehat{F}_0(t)|_{(T,T)}=\widehat{F}(t)|_{(T,T)}$$
then, under condition (5), $`\rho _TF=\rho _TF_0`$ and $`\rho _{\delta ,0}F^{}=\rho _{\delta ,0}F_0^{}`$ for all $`\delta T`$. If $`F_0(\tau )`$ behaves like a linear combination of homogeneous functions for large $`\tau `$ then $`\rho _{\delta ,0}F_0^{}`$ is of lower order than $`\rho _TF_0`$, so it plays the role of an error term in asymptotic formulae.
It is not always possible to find a model function $`F_0`$ satisfying (2.1). However, one can often construct $`\stackrel{~}{F}_0`$ in such a way that the convolutions $`\rho _T(F\stackrel{~}{F}_0)(\tau )`$ and $`\rho _{\delta ,0}(F^{}\stackrel{~}{F}_0^{})(\tau )`$ admit good estimates for large $`\tau `$ (roughly speaking, it happens if the Fourier transforms of $`F`$ and $`\stackrel{~}{F}_0`$ have similar singularities on the corresponding interval). Then the Tauberian theorems imply estimates with the error term
$$\pm \left(|\rho _T(F\stackrel{~}{F}_0)(\tau )|+|\rho _{\delta ,0}(F^{}\stackrel{~}{F}_0^{})(\tau )|\right).$$
In particular, if $`F`$ is the spectral or counting function of an elliptic partial differential operator with smooth coefficients then (1.11) gives a precise reminder estimate in the Weyl asymptotic formula, and the refined estimates (1.20), (1.21) allow one to obtain the second asymptotic term by letting $`T\mathrm{}`$ (see \[SV\] for details).
In applications to the second order differential operators it is usually more convenient to deal with the cosine Fourier transform of $`F^{}`$. The following elementary observation enables one to apply our results in the case where information on the sine Fourier transform of $`F^{}`$ is not available.
###### Proposition 2.1.
If the cosine Fourier transforms of the derivatives $`F^{}`$ and $`F_0^{}`$ coincide on an interval $`(\delta ,\delta )`$ then the Fourier transforms of the functions $`F(\tau )F(\tau )`$ and $`F_0(\tau )F_0(\tau )`$ coincide on the same interval.
### 2.2. Test functions $`\rho `$
In this subsection we consider a class of functions $`\rho `$ satisfying (1<sub>m</sub>)–(5) and estimate the constants $`c_{\rho ,\kappa }`$.
###### Lemma 2.2.
Let $`\zeta C^{m+1}[\frac{1}{2},\frac{1}{2}]`$ be a real-valued even function such that $`\zeta _{L_2}=1`$ and $`\zeta ^{(k)}(\pm \frac{1}{2})=0`$ for $`k=0,1,\mathrm{}m1`$, where $`\zeta ^{(k)}`$ denotes the $`k`$th derivative. If we extend $`\zeta `$ to $``$ by zero then $`\rho :=(\widehat{\zeta })^2`$ satisfies (1<sub>m</sub>)–(5) and
(2.2)
$$c_{\rho ,2k}=\zeta ^{(k)}_{L_2}^2,k=0,1,\mathrm{},m.$$
###### Proof.
The conditions (3) and (4) are obviously fulfilled; (2), (5) and (2.2) follow from the fact that $`\widehat{\rho }=(2\pi )^{1/2}\zeta \zeta `$. Finally, (1<sub>m</sub>) holds true because the $`(m+1)`$th derivative of the extended function $`\zeta `$ coincides with a linear combination of an $`L_1`$-function and two $`\delta `$-functions. ∎
The following lemma is a consequence of the uncertainty principle.
###### Lemma 2.3.
If $`\rho `$ is defined as in Lemma 2.2 then
(2.3)
$$\begin{array}{c}c_{\rho ,1}\frac{\pi }{2}.\end{array}$$
###### Proof.
Let $`\mathrm{\Pi }_a`$ be the multiplication operator and $`\widehat{\mathrm{\Pi }}_a`$ be the Fourier multiplier generated by the characteristic function of the interval $`[a,a]`$. Then the Hilbert-Schmidt norm of the operator $`\widehat{\mathrm{\Pi }}_{a_1}\mathrm{\Pi }_{a_2}`$ acting in $`L_2()`$ is equal to $`\sqrt{2\pi ^1a_1a_2}`$. Therefore
$$2_0^\mu \widehat{\zeta }^2(\tau )d\tau =\widehat{\mathrm{\Pi }}_\mu \mathrm{\Pi }_{1/2}\zeta _{L_2}^2\pi ^1\mu \zeta _{L_2}^2=\pi ^1\mu ,$$
which implies that
$$\begin{array}{c}c_{\rho ,1}=2_0^{\mathrm{}}\mu \widehat{\zeta }^2(\mu )d\mu =2_0^{\mathrm{}}_\mu ^{\mathrm{}}\widehat{\zeta }^2(\tau )d\tau d\mu \hfill \\ \hfill 2_0^\pi _\mu ^{\mathrm{}}\widehat{\zeta }^2(\tau )d\tau d\mu _0^\pi (1\pi ^1\mu )d\mu =\frac{\pi }{2}.\end{array}$$
###### Remark 2.4.
As follows from Nazarov’s theorem (see \[Na\] or \[HJ\]),
$$\begin{array}{c}_\mu ^{\mathrm{}}\widehat{\varphi }^2(\tau )d\tau b_1e^{b_2\mu },\varphi C_0^{\mathrm{}}(\frac{1}{2},\frac{1}{2}),\mu 0,\end{array}$$
where $`b_1,b_2>0`$ are some absolute constants. Using the estimates for $`b_1,b_2`$ obtained in \[Na\], one can slightly improve the estimate (2.3).
###### Example 2.5.
Let $`\stackrel{~}{\nu }_m`$ be the first eigenvalue of the operator $`{\displaystyle \frac{\mathrm{d}^{2m}}{\mathrm{d}t^{2m}}}`$ on the interval $`(\frac{1}{2},\frac{1}{2})`$ subject to Dirichlet boundary condition, and let $`\zeta _m`$ be the corresponding real even normalized eigenfunction. Denote $`\nu _m:=\left(\stackrel{~}{\nu }_m\right)^{\frac{1}{2m}}`$. If we define $`\rho `$ as in Lemma 2.2 then, in view of (2.2) and (1.5),
(2.4)
$$c_{\rho ,2m}=\nu _m^{2m},c_{\rho ,\kappa }\nu _m^\kappa ,\kappa <2m.$$
The eigenvalues $`\stackrel{~}{\nu }_m=\nu _m^{2m}`$ grow very fast as $`m\mathrm{}`$. The following lemma gives a rough estimate for $`\nu _m`$.
###### Lemma 2.6.
We have $`\nu _m2m\sqrt[2m]{3}`$ for all $`m2`$.
###### Proof.
If $`\varphi (t)=\left(\frac{1}{4}t^2\right)^m`$ and $`_{L_2}`$ is the norm in $`L_2(\frac{1}{2},\frac{1}{2})`$ then
(2.5)
$$\stackrel{~}{\nu }_m\frac{\varphi ^{(m)}_{L_2}^2}{\varphi _{L_2}^2}=\frac{(4m+1)!(m!)^2}{(2m+1)!(2m)!}2^{2m+1}(2m)!.$$
One can easily see that
$$\frac{2^{2m}(2m)!}{(2m)^{2m}}=\frac{2(m^21)\mathrm{}(m^2(m1)^2)}{m^{2m2}}\frac{2(m^2(m1)^2)}{m^2}\frac{3}{2}.$$
Therefore (2.5) implies the required estimate. ∎
### 2.3. Power like singularities
Assume that $`|F(\tau )|\mathrm{const}(|\tau |+1)^n`$ with a non-negative integer $`n`$ and define
$$\sigma _n:=\{\begin{array}{cc}0,\hfill & \text{if }n\text{ is odd,}\hfill \\ 1,\hfill & \text{if }n\text{ is even,}\hfill \end{array}m_n:=\{\begin{array}{cc}\frac{n+1}{2},\hfill & \text{if }n\text{ is odd,}\hfill \\ \frac{n+2}{2},\hfill & \text{if }n\text{ is even,}\hfill \end{array}$$
$$P_n^+(\tau ,\mu ):=\frac{(\tau +\mu )^n+(\tau \mu )^n}{2},P_n^{}(\tau ,\mu ):=\frac{\mu (\tau +\mu )^n\mu (\tau \mu )^n}{2}.$$
Clearly, $`P_n^\pm `$ are homogeneous polynomials in $`(\tau ,\mu )`$ with positive coefficients, which contain only even powers of $`\mu `$.
###### Lemma 2.7.
Let $`\rho `$ be a function satisfying (3), (5) and (1<sub>m</sub>) with $`m>\frac{n}{2}`$. If $`\mathrm{supp}F(0,+\mathrm{})`$ and the cosine Fourier transform of $`F^{}(\tau )`$ coincides on the interval $`(\delta ,\delta )`$ with the cosine Fourier transform of the function $`n\tau _+^{n1}`$ then
(2.6) $`\rho _\delta F(\tau )`$ $``$ $`{\displaystyle \left[P_n^+(\tau ,\delta ^1\mu )\sigma _n\delta ^n|\mu |^n\right]\rho (\mu )d\mu },`$
(2.7) $`\rho _\delta F(\tau )`$ $``$ $`{\displaystyle P_n^+(\tau ,\delta ^1\mu )\rho (\mu )d\mu },`$
(2.8) $`\rho _{\delta ,0}F^{}(\tau )`$ $``$ $`\delta ^2{\displaystyle \left[P_n^{}(\tau ,\delta ^1\mu )+\sigma _n\delta ^{n1}|\mu |^{n+1}\right]\rho (\mu )d\mu }`$
for all $`\tau >0`$.
###### Proof.
According to Proposition 2.1, the Fourier transform of $`F(\tau )F(\tau )`$ coincides on the interval $`(\delta ,\delta )`$ with the Fourier transform of
$$\mathrm{sign}\tau |\tau |^n=\left(12\sigma _n\chi _{}(\tau )\right)\tau ^n.$$
Since $`\rho `$ is even, this implies that
$$\begin{array}{c}\rho _\delta F(\tau )=\delta \left(12\sigma _n\chi _{}(\tau \mu )\right)(\tau \mu )^n\rho (\delta \mu )d\mu \hfill \\ \hfill =P_n^+(\tau ,\delta ^1\mu )\rho (\mu )d\mu \mathrm{\hspace{0.33em}2}\sigma _n_{\delta \tau }^{\mathrm{}}(\delta ^1\mu \tau )^n\rho (\mu )d\mu ,\end{array}$$
$$\begin{array}{c}\rho _{\delta ,0}F^{}(\tau )=\rho _{\delta ,0}^{}F(\tau )=\delta ^3\left(12\sigma _n\chi _{}(\tau \mu )\right)(\tau \mu )^n\mu \rho (\delta \mu )d\mu \hfill \\ \hfill =\delta ^2P_n^{}(\tau ,\delta ^1\mu )\rho (\mu )d\mu +\mathrm{\hspace{0.33em}2}\sigma _n\delta _{\delta \tau }^{\mathrm{}}(\delta ^1\mu \tau )^n\mu \rho (\mu )d\mu \end{array}$$
for all $`\tau >0`$. Estimating $`\mathrm{\hspace{0.17em}0}(\delta ^1\mu \tau )\delta ^1\mu `$ in the integrals on the right hand sides, we arrive at (2.6)–(2.8). ∎
The obvious inequalities
$`\tau ^n+\sigma _n|\nu |^nP_n^+(\tau ,\nu )\tau ^n+n|\nu |(\tau +|\nu |)^{n1},`$
$`P_n^{}(\tau ,\nu )+\sigma _n|\nu |^{n+1}n\nu ^2(\tau +|\nu |)^{n1}`$
and (2.6)–(2.8) imply that, for all $`\tau >0`$,
(2.9) $`0\rho _\delta F(\tau )\tau ^n`$ $``$ $`n\delta ^1{\displaystyle |\mu |(\tau +\delta ^1|\mu |)^{n1}\rho (\mu )d\mu },`$
(2.10) $`\rho _{\delta ,0}F^{}(\tau )`$ $``$ $`n{\displaystyle \mu ^2(\tau +\delta ^1|\mu |)^{n1}\rho (\mu )d\mu }.`$
Note that $`m_n`$ is the minimal positive integer which is greater than $`\frac{n}{2}`$. If $`\rho `$ is defined as in Lemma 2.2 with $`m=m_n`$ then, by (2.2),
(2.11)
$$P_n^\pm (\tau ,\delta ^1\mu )\rho (\mu )d\mu =(P_n^\pm (\tau ,\delta ^1D_t)\zeta ,\zeta )_{L_2}.$$
Applying (2.6)–(2.11) and (1.11) or (1.18), (1.19), one can obtain various estimates for $`F(\tau )`$.
###### Example 2.8.
Let $`n=3`$ and $`\zeta `$ be an arbitrary function satisfying conditions of Lemma 2.2 with $`m=m_n=2`$. If the conditions of Lemma 2.7 are fulfilled then (2.6)–(2.8), (2.11) and (1.19), (1.20) with $`T=\delta `$ imply that
$`F(\tau )`$ $``$ $`\tau ^3{\displaystyle \frac{3\epsilon \tau ^2}{2}}+\epsilon ^2\tau {\displaystyle \frac{\epsilon ^3}{4}}+{\displaystyle \frac{3}{2\delta ^2}}\left(\tau {\displaystyle \frac{\tau ^2}{\epsilon }}{\displaystyle \frac{\epsilon }{2}}\right)\zeta ^{}_{L_2}^2{\displaystyle \frac{1}{\epsilon \delta ^4}}\zeta ^{\prime \prime }_{L_2}^2,`$
$`F(\tau )`$ $``$ $`\tau ^3+{\displaystyle \frac{3\epsilon \tau ^2}{2}}+\epsilon ^2\tau +{\displaystyle \frac{\epsilon ^3}{4}}+{\displaystyle \frac{3}{2\delta ^2}}\left(\tau +{\displaystyle \frac{\tau ^2}{\epsilon }}+{\displaystyle \frac{\epsilon }{2}}\right)\zeta ^{}_{L_2}^2+{\displaystyle \frac{1}{\epsilon \delta ^4}}\zeta ^{\prime \prime }_{L_2}^2`$
for all $`\epsilon >0`$ and $`\tau >0`$. Thus, $`F(\tau )`$ lies between the first Dirichlet eigenvalues of ordinary differential operators generated by the quadratic forms on the right hand sides of the above inequalities.
###### Corollary 2.9.
Under conditions of Lemma 2.7
(2.12) $`F(\tau )`$ $``$ $`\tau ^n2\pi ^1\nu _{m_n}^2n\delta ^1(\tau +\delta ^1\nu _{m_n})^{n1},`$
(2.13) $`F(\tau )`$ $``$ $`\tau ^n+(2\pi ^1\nu _{m_n}^2+\nu _{m_n})n\delta ^1(\tau +\delta ^1\nu _{m_n})^{n1}`$
for all $`\tau >0`$.
###### Proof.
If we define $`\rho `$ as in Lemma 2.2 with $`\zeta =\zeta _m`$ (see Example 2.5) then (2.12), (2.13) follow from (1.11) with $`T=\delta `$, (2.9), (2.10), (2.3) and (2.4). ∎
###### Corollary 2.10.
Under conditions of Lemma 2.7
(2.14) $`{\displaystyle _0^{\lambda ^2}}F(\sqrt{\mu })d\mu `$ $``$ $`{\displaystyle \frac{2\lambda ^{n+2}}{n+2}}2n\nu _{m_n}^2\delta ^2\lambda (\lambda +\delta ^1\nu _{m_n})^{n1},`$
(2.15) $`{\displaystyle _0^{\lambda ^2}}F(\sqrt{\mu })d\mu `$ $``$ $`{\displaystyle \frac{2\lambda ^{n+2}}{n+2}}+(n+1)\nu _{m_n}^2\delta ^2(\lambda +\delta ^1\nu _{m_n})^n`$
for all $`\lambda >0`$.
###### Proof.
Since $`_0^{\lambda ^2}F(\sqrt{\mu })d\mu =2_0^\lambda F(\tau )\tau d\tau `$, Theorem 1.6 with $`T=\delta `$, $`a=0`$, $`b=\lambda `$ and $`f(\tau )=\tau `$ implies
(2.16) $`{\displaystyle _0^{\lambda ^2}}F(\sqrt{\mu })d\mu `$ $``$ $`2{\displaystyle _0^\lambda }\tau \rho _\delta F(\tau )d\tau 2\delta ^2\lambda \rho _{\delta ,0}F^{}(\lambda ),`$
(2.17) $`{\displaystyle _0^{\lambda ^2}}F(\sqrt{\mu })d\mu `$ $``$ $`2{\displaystyle _0^\lambda }\left(\tau \rho _\delta F(\tau )+\delta ^2\rho _{\delta ,0}F^{}(\tau )\right)d\tau .`$
Let $`\rho `$ be defined as in Lemma 2.2 with $`\zeta =\zeta _m`$. Then (2.14) follows from (2.16), (2.9), (2.10) and (2.4). Since $`\tau P_n^+(\tau ,\nu )+P_n^{}(\tau ,\nu )=P_{n+1}^+(\tau ,\nu )`$, the inequality (2.17) and (2.7), (2.8) imply that
$$_0^{\lambda ^2}F(\sqrt{\mu })d\mu 2_0^\lambda \left(P_{n+1}^+(\tau ,\delta ^1\mu )+\sigma _n|\delta ^1\mu |^{n+1}\right)\rho (\mu )d\mu d\tau .$$
Estimating
$$_0^\lambda \left[P_{n+1}^+(\tau ,\nu )+\sigma _n|\nu |^{n+1}\right]d\tau \frac{\lambda ^{n+2}}{n+2}+\frac{n+1}{2}\nu ^2(\lambda +|\nu |)^n$$
with $`\nu =\delta ^1\mu `$ and applying (2.4), we obtain (2.15). ∎
## 3. Applications to the Laplace operator
Let $`\mathrm{\Omega }^n`$ be an open domain and $`d(x)`$ be the distance from $`x\mathrm{\Omega }`$ to the boundary $`\mathrm{\Omega }`$.
### 3.1. Estimates of the spectral function
Consider the Laplacian $`\mathrm{\Delta }_B`$ in $`\mathrm{\Omega }`$ subject to a self-adjoint boundary condition $`B(x,D_x)u|_\mathrm{\Omega }=0`$, where $`B`$ is a differential operator. Assume that the operator $`\mathrm{\Delta }_B`$ is non-negative and denote by $`\mathrm{\Pi }(\lambda )`$ its spectral projection corresponding to the interval $`[0,\lambda )`$. Let $`e(x,y;\lambda )`$ be the integral kernel of the operator $`\frac{\mathrm{\Pi }(\lambda 0)+\mathrm{\Pi }(\lambda +0)}{2}`$ (the so-called spectral function). The Sobolev embedding theorem implies that $`e(x,y;\lambda )`$ is a smooth function on $`\mathrm{\Omega }\times \mathrm{\Omega }`$ for each fixed $`\lambda `$ and that $`e(x,x;\lambda )`$ is a non-decreasing polynomially bounded function of $`\lambda `$ for each fixed $`x\mathrm{\Omega }`$.
Let $`\mathrm{\Delta }_0`$ be the Laplacian on $`^n`$, and $`e_0(x,y;\lambda )`$, $`\stackrel{~}{e}_0(x,y;\lambda )`$, $`\stackrel{~}{e}(x,y;\lambda )`$ be the spectral functions of the operators $`\mathrm{\Delta }_0`$, $`\sqrt{\mathrm{\Delta }}_0`$, $`\sqrt{\mathrm{\Delta }_B}`$ respectively. Then
$`\chi _+(\tau )e(x,x;\tau ^2)`$ $`=`$ $`\stackrel{~}{e}(x,x;\tau ),`$
$`\chi _+(\tau )e_0(x,x;\tau ^2)`$ $`=`$ $`\stackrel{~}{e}_0(x,x;\tau )=C_n\tau _+^n,`$
where
(3.1)
$$C_n:=(2\pi )^n\mathrm{meas}\{\xi ^n:|\xi |<1\}.$$
By the spectral theorem, the cosine Fourier transform of $`\frac{\mathrm{d}}{\mathrm{d}\tau }\stackrel{~}{e}(x,y;\tau )`$ coincides with the fundamental solution $`u(x,y;t)`$ of the wave equation in $`\mathrm{\Omega }`$,
$$u_{tt}=\mathrm{\Delta }u,Bu|_\mathrm{\Omega }=0,u|_{t=0}=\delta (xy),u_t|_{t=0}=0.$$
Due to the finite speed of propagation, $`u(x,x;t)`$ is equal to $`u_0(x,x;t)`$ whenever $`t(d(x),d(x))`$, where $`u_0(x,y;t)`$ is the fundamental solution of the wave equation in $`^n`$. Thus, the cosine Fourier transforms of the derivatives $`\frac{\mathrm{d}}{\mathrm{d}\tau }\stackrel{~}{e}_0(x,x;\tau )`$ and $`\frac{\mathrm{d}}{\mathrm{d}\tau }\stackrel{~}{e}(x,x;\tau )`$ coincide on the time interval $`(d(x),d(x))`$. Applying (2.12)–(2.15) to $`F(\tau )=C_n^1\stackrel{~}{e}(x,x;\tau )`$ we obtain the following corollary.
###### Corollary 3.1.
For every $`x\mathrm{\Omega }`$ and all $`\lambda >0`$ we have
(3.2) $`e(x,x;\lambda )C_n\lambda ^{n/2}{\displaystyle \frac{nC_n\mathrm{\hspace{0.17em}2}\pi ^1\nu _{m_n}^2}{d(x)}}\left(\lambda ^{1/2}+{\displaystyle \frac{\nu _{m_n}}{d(x)}}\right)^{n1},`$
(3.3) $`e(x,x;\lambda )C_n\lambda ^{n/2}+{\displaystyle \frac{nC_n(2\pi ^1\nu _{m_n}^2+\nu _{m_n})}{d(x)}}\left(\lambda ^{1/2}+{\displaystyle \frac{\nu _{m_n}}{d(x)}}\right)^{n1},`$
(3.4) $`{\displaystyle _0^\lambda }e(x,x;\mu )d\mu {\displaystyle \frac{2C_n\lambda ^{n/2+1}}{n+2}}{\displaystyle \frac{2nC_n\nu _{m_n}^2\lambda ^{1/2}}{(d(x))^2}}\left(\lambda ^{1/2}+{\displaystyle \frac{\nu _{m_n}}{d(x)}}\right)^{n1}`$
(3.5) $`{\displaystyle _0^\lambda }e(x,x;\mu )d\mu {\displaystyle \frac{2C_n\lambda ^{n/2+1}}{n+2}}+{\displaystyle \frac{(n+1)C_n\nu _{m_n}^2}{(d(x))^2}}\left(\lambda ^{1/2}+{\displaystyle \frac{\nu _{m_n}}{d(x)}}\right)^n.`$
### 3.2. Estimates of the counting function of the Dirichlet Laplacian
In this subsection we shall be assuming that $`|\mathrm{\Omega }|<\mathrm{}`$, where $`||`$ denotes the $`n`$-dimensional Lebesgue measure.
Consider the positive operator $`\mathrm{\Delta }_D`$, where $`\mathrm{\Delta }_D`$ is the Dirichlet Laplacian in $`\mathrm{\Omega }`$. Let $`N(\lambda )`$ be the number of its eigenvalues lying below $`\lambda `$. The following theorem is due to F. Berezin \[B\].
###### Theorem 3.2.
For all $`\lambda 0`$ we have
(3.6)
$$_0^\lambda N(\mu )d\mu \frac{2}{n+2}C_n|\mathrm{\Omega }|\lambda ^{n/2+1}.$$
This results was reproduced in \[La\]. A. Laptev also noticed that the famous Li–Yau estimate
(3.7)
$$N(\lambda )(1+2/n)^{n/2}C_n|\mathrm{\Omega }|\lambda ^{n/2},\lambda 0,$$
(see \[LY\]) is a one line consequence of (3.6). Indeed, (3.7) can be proved by estimating
(3.8)
$$N(\lambda )(\theta \lambda )^1_0^{\lambda +\theta \lambda }N(\mu )d\mu \frac{2(1+\theta )^{n/2+1}}{(n+2)\theta }C_n|\mathrm{\Omega }|\lambda ^{n/2}$$
and optimizing the choice of $`\theta >0`$.
###### Remark 3.3.
In \[B\] F. Berezin proved an analogue of (3.6) for general operators with constant coefficients subject to Dirichlet boundary condition. In the same way as above, applying the first inequality (3.8) and Berezin’s estimates, one can easily obtain upper bounds for the corresponding counting functions (see \[La\]).
According to the Weyl asymptotic formula
(3.9)
$$N(\lambda )=C_n|\mathrm{\Omega }|\lambda ^{n/2}+o(\lambda ^{n/2}),\lambda +\mathrm{},$$
(in the general case (3.9) was proved in \[BS\]). The coefficient in the right hand side of (3.7) contains an extra factor $`(1+2/n)^{n/2}`$. G. Pólya conjectured \[P\] that (3.7) holds without this factor. However, this remains an open problem.
Given a positive $`\epsilon `$, denote
$$\mathrm{\Omega }_\epsilon ^\mathrm{b}:=\{x\mathrm{\Omega }:d(x)\epsilon \},\mathrm{\Omega }_\epsilon ^\mathrm{i}:=\{x\mathrm{\Omega }:d(x)>\epsilon \}.$$
If
(3.10)
$$|\mathrm{\Omega }_\epsilon ^\mathrm{b}|\mathrm{const}\epsilon ^r,r(0,1],$$
then, using the variational method \[CH\], one can prove that
(3.11)
$$|N(\lambda )C_n|\mathrm{\Omega }|\lambda ^{n/2}|\{\begin{array}{cc}\mathrm{const}\lambda ^{(n1)/2}\mathrm{ln}\lambda ,\hfill & r=1,\hfill \\ \mathrm{const}\lambda ^{(nr)/2},\hfill & r<1.\hfill \end{array}$$
It is well known that in the smooth case
$$N(\lambda )C_n|\mathrm{\Omega }|\lambda ^{n/2}=O(\lambda ^{(n1)/2})$$
(see, for example, \[I1\] or \[SV\]), but it is not clear whether this estimate remains valid for an arbitrary domain satisfying (3.10) with $`r=1`$.
There is a number of papers devoted to estimates of the remainder term in the Weyl formula. In \[BL\] the authors, applying the variational technique, obtain explicit estimates for the constants in (3.11). In order to prove the estimate of $`N(\lambda )`$ from above, they imposed an additional condition on the outer neighbourhood of the boundary $`\mathrm{\Omega }`$, but this condition can probably be removed \[Ne\]. In \[Kr\] the author estimated the remainder term with the use of a different technique (similar to that in \[LY\]); his results seem to be less precise than those obtained in \[BL\].
Let
$$N_\epsilon ^\mathrm{b}(\lambda ):=_{\mathrm{\Omega }_\epsilon ^\mathrm{b}}e(x,x;\lambda )dx,N_\epsilon ^\mathrm{i}(\lambda ):=_{\mathrm{\Omega }_\epsilon ^\mathrm{i}}e(x,x;\lambda )dx.$$
Then $`N(\lambda )=N_\epsilon ^\mathrm{b}(\lambda )+N_\epsilon ^\mathrm{i}(\lambda )`$ for every $`\epsilon >0`$.
###### Corollary 3.4.
For all $`\lambda >0`$ and $`\epsilon >0`$ we have
(3.12) $`N_\epsilon ^\mathrm{i}(\lambda )C_n|\mathrm{\Omega }_\epsilon ^\mathrm{i}|\lambda ^{n/2}C_{n,1}(\lambda ^{1/2}+\epsilon ^1\nu _{m_n})^{n1}{\displaystyle _{\mathrm{\Omega }_\epsilon ^\mathrm{i}}}{\displaystyle \frac{\mathrm{d}x}{d(x)}},`$
(3.13) $`N_\epsilon ^\mathrm{i}(\lambda )C_n|\mathrm{\Omega }_\epsilon ^\mathrm{i}|\lambda ^{n/2}+C_{n,2}(\lambda ^{1/2}+\epsilon ^1\nu _{m_n})^{n1}{\displaystyle _{\mathrm{\Omega }_\epsilon ^\mathrm{i}}}{\displaystyle \frac{\mathrm{d}x}{d(x)}},`$
(3.14) $`N_\epsilon ^\mathrm{b}(\lambda )C_{n,3}|\mathrm{\Omega }_\epsilon ^\mathrm{b}|\lambda ^{n/2}+C_{n,4}\lambda ^{1/2}(\lambda ^{1/2}+\epsilon ^1\nu _{m_n})^{n1}{\displaystyle _{\mathrm{\Omega }_\epsilon ^\mathrm{i}}}{\displaystyle \frac{\mathrm{d}x}{(d(x))^2}},`$
where
$$\begin{array}{ccc}& C_{n,1}=nC_n\mathrm{\hspace{0.17em}2}\pi ^1\nu _{m_n}^2,\hfill & C_{n,2}=nC_n(2\pi ^1\nu _{m_n}^2+\nu _{m_n}),\hfill \\ & & \\ & C_{n,3}=(1+2/n)^{n/2}C_n,\hfill & C_{n,4}=(1+2/n)^{n/2}n^2C_n\nu _{m_n}^2.\hfill \end{array}$$
###### Proof.
The inequalities (3.12), (3.13) are proved by straightforward integration of (3.2), (3.3). Theorem 3.2 and (3.4) imply that
(3.15)
$$\begin{array}{c}_0^\lambda N_\epsilon ^\mathrm{b}(\mu )d\mu =_0^\lambda N(\lambda )d\mu _0^\lambda N_\epsilon ^\mathrm{i}(\mu )d\mu \hfill \\ \hfill \frac{2}{n+2}C_n|\mathrm{\Omega }_\epsilon ^\mathrm{b}|\lambda ^{n/2+1}+2nC_n\nu _{m_n}^2\lambda ^{1/2}\left(\lambda ^{1/2}+\epsilon ^1\nu _{m_n}\right)^{n1}_{\mathrm{\Omega }_\epsilon ^\mathrm{i}}\frac{\mathrm{d}x}{(d(x))^2}.\end{array}$$
Now, applying the first inequality (3.8) with $`\theta =2/n`$, we arrive at (3.14). ∎
Adding up the inequalities (3.13) and (3.14) we obtain
(3.16)
$$\begin{array}{c}N(\lambda )C_n|\mathrm{\Omega }_\epsilon ^\mathrm{i}|\lambda ^{n/2}+C_{n,3}|\mathrm{\Omega }_\epsilon ^\mathrm{b}|\lambda ^{n/2}\hfill \\ \hfill +(\lambda ^{1/2}+\epsilon ^1\nu _{m_n})^{n1}_{\mathrm{\Omega }_\epsilon ^\mathrm{i}}\frac{C_{n,2}d(x)+C_{n,4}\lambda ^{1/2}}{(d(x))^2}dx,\epsilon >0.\end{array}$$
Since
(3.17)
$$_{\mathrm{\Omega }_\epsilon ^\mathrm{i}}\frac{\mathrm{d}x}{(d(x))^j}=_\epsilon ^{\mathrm{}}s^j\mathrm{d}(|\mathrm{\Omega }_s^\mathrm{b}|)=j_\epsilon ^{\mathrm{}}s^{j1}|\mathrm{\Omega }_s^\mathrm{b}|ds\epsilon ^j|\mathrm{\Omega }_\epsilon ^\mathrm{b}|,$$
(3.10) and the inequalities (3.12), (3.16) with $`\epsilon =\lambda ^{1/2}`$ imply (3.11).
By (3.12) and (3.16) we have
$`C_n|\mathrm{\Omega }_\epsilon ^\mathrm{b}||\mathrm{\Omega }_\epsilon ^\mathrm{i}|{\displaystyle \frac{C_{n,1}}{\epsilon \lambda ^{1/2}}}\left(1+{\displaystyle \frac{\nu _{m_n}}{\epsilon \lambda ^{1/2}}}\right)^{n1}`$
(3.18) $`\lambda ^{n/2}N(\lambda )C_n|\mathrm{\Omega }|`$
$`(C_{n,3}C_n)|\mathrm{\Omega }_\epsilon ^\mathrm{b}|+|\mathrm{\Omega }_\epsilon ^\mathrm{i}|\left({\displaystyle \frac{C_{n,2}}{\epsilon \lambda ^{1/2}}}+{\displaystyle \frac{C_{n,4}}{\epsilon ^2\lambda }}\right)\left(1+{\displaystyle \frac{\nu _{m_n}}{\epsilon \lambda ^{1/2}}}\right)^{n1}`$
for all $`\epsilon >0`$. If $`\epsilon \mathrm{}`$ then the second inequality (3.18) turns into (3.7). Since $`|\mathrm{\Omega }_\epsilon ^\mathrm{b}|0`$ as $`\epsilon 0`$, (3.18) implies (3.9). Moreover, taking $`\epsilon =\lambda ^\kappa `$ with an arbitrary $`\kappa (0,\frac{1}{2})`$, we obtain the Weyl formula with a remainder estimate
$$\lambda ^{n/2}N(\lambda )C_n|\mathrm{\Omega }|=O(|\mathrm{\Omega }_{\lambda ^\kappa }^\mathrm{b}|+\lambda ^{\kappa 1/2}),\lambda +\mathrm{}.$$
###### Remark 3.5.
If the condition (3.10) is fulfilled then integrating (3.4) over $`\mathrm{\Omega }_{\lambda ^{1/2}}^\mathrm{i}`$, applying (3.17) and taking into account (3.6), we see that
(3.19)
$$\lambda ^1_0^\lambda N(\mu )d\mu =\frac{2}{n+2}C_n|\mathrm{\Omega }|\lambda ^{n/2}+O(\lambda ^{(nr)/2})$$
for all $`r(0,1]`$.
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# Thermodynamics, strange quark matter, and strange stars
## I Introduction
During the 10-plus years which have elapsed since Witten’s conjecture that strange quark matter (SQM), rather than the normal nuclear matter, might be the true ground state of quantum chromodynamics (QCD), many theoretical and observational efforts have been made on the investigation of its properties and potential astrophysical significance . Because of the well-known difficulty of QCD in the nonperturbative domain, phenomenological models reflecting the characteristic of strong interactions are widely used in the study of hadron, and many of them have been successfully applied to investigating the stability and properties of SQM. One of the most famous models is the MIT bag model with which Jaffe et. al find that SQM is absolutely stable around the normal nuclear density for a wide range of parameters. A vast number of further investigations are performed with fruitful results. A recent important result is that young millisecond pulsars are most likely to be strange stars rather than neutron stars . Another alternative model is the mass-density-dependent model with which Chakrabarty et al. obtained significantly different results . However, Benvenuto and Lugones pointed out that it is caused by the wrong thermodynamic treatment. They added an extra term to the expression of both pressure and energy, and got similar results as that in the bag model. A recent investigation indicates a link of SQM to the study of quark condensates while a more recent work has carefully studied the relation between the charge and critical density of SQM .
Latterly, we have demonstrated that the previous treatments have unreasonable vacuum limits . In addition to this problem, there exist another serious problem, i.e., the zero pressure does not appear in the lowest energy state. In fact, there are two important problems in the quark mass-density-dependent model. One is how to determine the quark mass scaling. The other is how to deal with the thermodynamics with density-dependent particle masses self consistently. We have mainly concentrated on the first problem in Ref. . The present paper will concentrate more on the second problem. We find that the extra term provided in Ref. should indeed be appended to the expression of pressure. However, it should not appear in that of energy according to both the general ensemble theory and basic thermodynamic principle. After our modification, the zero pressure point appears exactly at the lowest energy state, and thus the thermodynamics with density-dependent particle masses becomes self-consistent, which leads to completely different density behaviour of the sound velocity in SQM and different structure of strange stars.
We organize this paper as such. In the subsequent section, we give detail arguments on why the additional term in the pressure should not appear in the energy. The thermodynamic expression needed later are all derived carefully in this section. Then in Sec. III, we apply the new thermodynamic formulas to investigating strange quark matter, and find that the density behaviour of the sound velocity is opposite to the previous treatment but consistent with our recent publication. On application of our obtained equation of state, we integrate the equations of steller structure for strange stars in Sec. IV, and find that strange stars are dimensionally smaller and less massive than the previous calculation if SQM is absolutely stable. Sec. V is a short summary.
## II Thermodynamics with density-dependent particle masses
Let us explore directly from the general ensemble theory what the expression of pressure and energy should look like if the particle masses are dependent on density.
We express the density matrix as
$$\rho =\frac{1}{\mathrm{\Xi }}e^{\beta \left(E_{N_i,\alpha }_i\mu _iN_i\right)},$$
(1)
where $`\mathrm{\Xi }`$ is the partition function, $`\beta `$ is the reverse temperature, $`N_i`$ are the particle numbers, and $`\mu _i`$ are the corresponding chemical potentials. The microscopic energy $`E_{N_i,\alpha }`$ is a function of the system volume $`V`$, the particle masses $`m_i`$, the particle numbers $`N_i`$, and the other quantum numbers $`\alpha `$. The pressure of the system is
$`P`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Xi }}}{\displaystyle \underset{\{N_i\},\alpha }{}}\left({\displaystyle \frac{E_{N_i,\alpha }}{V}}\right)e^{\beta \left(E_{N_i,\alpha }_i\mu _iN_i\right)}`$ (2)
$`=`$ $`{\displaystyle \frac{1}{\mathrm{\Xi }}}{\displaystyle \underset{\{N_i\},\alpha }{}}\left[{\displaystyle \frac{1}{\beta }}{\displaystyle \frac{}{V}}e^{\beta \left(E_{N_i,\alpha }_i\mu _iN_i\right)}\right]`$ (3)
$`=`$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \frac{\mathrm{ln}\mathrm{\Xi }}{V}}|_{T,\{\mu _i\}}={\displaystyle \frac{(V\mathrm{\Omega })}{V}}|_{T,\{\mu _i\}},`$ (4)
where
$$\mathrm{\Omega }\frac{1}{V\beta }\mathrm{ln}\mathrm{\Xi }$$
(5)
is the thermodynamic potential density which is generally a function of the temperature $`T`$, the chemical potentials $`\mu _i`$, and the particle masses $`m_i`$. If the particle masses have nothing to do with the baryon number density $`n_b=N/(3V)`$ (N is the total particle number), we simply get
$$P=\mathrm{\Omega }.$$
(6)
If the masses depend on density or volume, one should have
$$P=\mathrm{\Omega }+n_b\frac{\mathrm{\Omega }}{n_b}.$$
(7)
This is just the right thing the authors have done in Ref. . The authors derive it as such:
$$P=\frac{(\mathrm{\Omega }/n_b)}{(1/n_b)}|_{T,\mu _i}=n_b\frac{\mathrm{\Omega }}{n_b}\mathrm{\Omega }.$$
(8)
For canonical ensemble, the particle numbers $`N_i`$ keep fixed. This derivation is thus obvious. However, it is not so obvious for grand canonical ensemble because the particle number is not necessarily constant when the temperature T and chemical potentials $`\mu _i`$ unchanged. We will give a more convincing derivation a little later.
The additional term is of crucial importance for pressure balance. Not as done in Ref. , however, the extra term does not appear in the expression of energy. Now, let’s calculate the statistic average for the energy:
$`\overline{E}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Xi }}}{\displaystyle \underset{\{N_i\},\alpha }{}}E_{N_i,\alpha }e^{\beta \left(E_{N_i,\alpha }_i\mu _iN_i\right)}`$ (9)
$`=`$ $`{\displaystyle \frac{1}{\mathrm{\Xi }}}{\displaystyle \underset{\{N_i\},\alpha }{}}\left({\displaystyle \frac{}{\beta }}+{\displaystyle \underset{i}{}}\mu _iN_i\right)e^{\beta \left(E_{N_i,\alpha }_i\mu _iN_i\right)}`$ (10)
$`=`$ $`{\displaystyle \frac{}{\beta }}\mathrm{ln}\mathrm{\Xi }+{\displaystyle \underset{i}{}}\mu _i\overline{N_i},`$ (11)
where
$`\overline{N_i}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Xi }}}{\displaystyle \underset{\{N_i\},\alpha }{}}N_ie^{\beta \left(E_{N_i,\alpha }_i\mu _iN_i\right)}`$ (12)
$`=`$ $`{\displaystyle \frac{1}{\mathrm{\Xi }}}{\displaystyle \underset{\{N_i\},\alpha }{}}\left[{\displaystyle \frac{1}{\beta }}{\displaystyle \frac{}{\mu _i}}e^{\beta \left(E_{N_i,\alpha }_i\mu _iN_i\right)}\right]_{V,T,\{m_k\}}`$ (13)
$`=`$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \frac{}{\mu _i}}\mathrm{ln}\mathrm{\Xi }=V{\displaystyle \frac{\mathrm{\Omega }}{\mu _i}}|_{T,\{m_k\}}`$ (14)
is the average number for particle type $`i`$. Therefore, the energy density of the system is
$`E`$ $`=`$ $`{\displaystyle \frac{\overline{E}}{V}}={\displaystyle \frac{(\beta \mathrm{\Omega })}{\beta }}+{\displaystyle \underset{i}{}}\mu _in_i`$ (15)
$`=`$ $`\mathrm{\Omega }+\beta {\displaystyle \frac{\mathrm{\Omega }}{\beta }}+{\displaystyle \underset{i}{}}\mu _in_i`$ (16)
$`=`$ $`\mathrm{\Omega }+{\displaystyle \underset{i}{}}\mu _in_iT{\displaystyle \frac{\mathrm{\Omega }}{T}},`$ (17)
where $`n_i`$ is the number density of particle type $`i`$:
$$n_i\frac{\overline{N_i}}{V}=\frac{\mathrm{\Omega }}{\mu _i}|_{T,\{m_k\}}.$$
(18)
It is clear from Eq. (17) that only when Eq. (6) holds can one get the Eq. (8) in Ref. . Therefore, we should not, as done in Ref. , use that expression to calculate the energy density. Instead, we will calculate $`E`$ directly from Eq. (17) in this paper.
For more evident arguments, let’s see the following derivation starting from the basic derivative relation for an open system:
$$d(VE)=Td(VS)PdV+\underset{i}{}\mu _id\overline{N_i},$$
(19)
where S is the entropy density of the system.
Choosing T, V, and {$`\overline{N_i}`$} as the independent macroscopic state variables, the combined statement of the first and second laws of thermodynamics, Eq. (19), can be expressed as
$$d(VA)=VSdTPdV+\underset{i}{}\mu _id\overline{N_i},$$
(20)
where $`AETS`$ is the Helmholtz free energy density by which we have
$`P`$ $`=`$ $`{\displaystyle \frac{d(VA)}{dV}}|_{T,\{\overline{N_i}\}}`$ (21)
$`=`$ $`AV{\displaystyle \frac{dA}{dV}}|_{T,\{\overline{N_i}\}}`$ (22)
$`=`$ $`A+{\displaystyle \underset{j}{}}n_j{\displaystyle \frac{dA(T,\{n_i\})}{dn_j}}|_T.`$ (23)
This is a general expression for pressure. In obtaining the third equality, we have used the chain relation
$$V\frac{d}{dV}f\left(\{n_i=\frac{\overline{N_i}}{V}\}\right)|_{\{\overline{N_i}\}}=\underset{j}{}n_j\frac{d}{dn_j}f(\{n_i\}),$$
(24)
where $`f`$ is an arbitrary function.
According to the basic relation between thermodynamics and statistics, we have
$$A=\mathrm{\Omega }+\underset{i}{}\mu _in_i$$
(25)
where $`\mathrm{\Omega }`$ is the thermodynamic potential density. For a free Fermi system, it is
$`\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{g_iT}{2\pi ^2}}{\displaystyle _0^{\mathrm{}}}\mathrm{ln}\left[1+e^{\beta \left(\sqrt{p^2+m_i^2}\mu _i\right)}\right]p^2𝑑p`$ (26)
$``$ $`{\displaystyle \underset{i}{}}\mathrm{\Omega }_i(T,\mu _i,m_i),`$ (27)
where $`g_i`$ is the degeneracy factor which is 6 for quarks and 2 for electrons. In order to include the interaction between particles, we regard the particle masses as density-dependent, namely
$$m_i=m_i\left(n_b\underset{j}{}n_j/3\right).$$
(28)
Because we have chosen $`T`$, $`V`$, and $`\{\overline{N_i}\}`$ as independent state variables, the chemical potential $`\mu _i`$ should also be regarded as a function of T and $`\{n_k\}`$, namely
$$\mu _i=\mu _i(T,\{n_k\}).$$
(29)
So, the total derivative of $`\mathrm{\Omega }(T,\{\mu _k\},\{m_k\})`$ with respect to $`n_j`$ should be taken like this:
$`{\displaystyle \frac{d\mathrm{\Omega }}{dn_j}}|_T`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{\mathrm{\Omega }}{\mu _i}}|_{T,\{m_k\}}{\displaystyle \frac{d\mu _i}{dn_j}}|_T+{\displaystyle \frac{\mathrm{\Omega }}{n_b}}|_{T,\{\mu _k\}}{\displaystyle \frac{n_b}{n_j}}`$ (30)
$`=`$ $`{\displaystyle \underset{i}{}}n_i{\displaystyle \frac{d\mu _i}{dn_j}}|_T+{\displaystyle \frac{1}{3}}{\displaystyle \frac{\mathrm{\Omega }}{n_b}}|_{T,\{\mu _k\}}.`$ (31)
Here we have used the Eq. (18) and the fact that $`n_b/n_j=1/3`$.
Substituting Eq. (25) into Eq. (23) gives
$`P`$ $`=`$ $`A+{\displaystyle \underset{j}{}}n_j{\displaystyle \frac{d}{dn_j}}\left[\mathrm{\Omega }+{\displaystyle \underset{i}{}}\mu _in_i\right]_T`$ (32)
$`=`$ $`A+{\displaystyle \underset{j}{}}n_j\left[{\displaystyle \frac{d\mathrm{\Omega }}{dn_j}}|_T+{\displaystyle \underset{i}{}}\left(n_i{\displaystyle \frac{d\mu _i}{dn_j}}|_T+\mu _i{\displaystyle \frac{dn_i}{dn_j}}\right)\right]`$ (33)
$`=`$ $`A+{\displaystyle \underset{i}{}}\mu _in_i+{\displaystyle \underset{j}{}}{\displaystyle \frac{n_j}{3}}{\displaystyle \frac{\mathrm{\Omega }}{n_b}}|_{T,\{\mu _k\}}`$ (34)
$`=`$ $`\mathrm{\Omega }+n_b{\displaystyle \frac{\mathrm{\Omega }}{n_b}}|_{T,\{\mu _k\}}`$ (35)
$`=`$ $`{\displaystyle \underset{i}{}}\left(\mathrm{\Omega }_i+n_b{\displaystyle \frac{m_i}{n_b}}{\displaystyle \frac{\mathrm{\Omega }_i}{m_i}}\right).`$ (36)
At zero temperature, the corresponding thermodynamic potential density can be obtained from Eq. (26) by first taking the limit $`T0`$ and then carrying out the resulting integration:
$`\mathrm{\Omega }={\displaystyle \underset{i}{}}{\displaystyle \frac{g_i}{48\pi ^2}}[\mu _i(\mu _i^2m_i^2)^{1/2}(2\mu _i^25m_i^2)`$ (37)
$`+3m_i^4\mathrm{ln}{\displaystyle \frac{\mu _i+\sqrt{\mu _i^2m_i^2}}{m_i}}].`$ (38)
We thus have from Eqs. (18), (17), and (36):
$`n_i`$ $`=`$ $`{\displaystyle \frac{g_i}{6\pi ^2}}(\mu _i^2m_i^2)^{3/2},`$ (39)
$`E`$ $`=`$ $`{\displaystyle \underset{i}{}}m_in_iF(x_i),`$ (40)
$`P`$ $`=`$ $`{\displaystyle \underset{i}{}}m_in_ix_i^2G(x_i){\displaystyle \underset{i}{}}m_in_if(x_i)`$ (41)
where
$$x_i\frac{p_{f,i}}{m_i}=\frac{\left(\frac{6\pi ^2}{g_i}n_i\right)^{1/3}}{m_i}=\frac{\sqrt{\mu _i^2m_i^2}}{m_i}$$
(42)
is the ratio of the Fermi momentum to the mass that related to particle type $`i`$. With the hyperbolic sine function sh$`{}_{}{}^{1}(x)\mathrm{ln}(x+\sqrt{x^2+1})`$, the functions $`F(x_i)`$, $`G(x_i)`$, and $`f(x_i)`$ are defined as
$`F(x_i){\displaystyle \frac{3}{8}}\left[x_i\sqrt{x_i^2+1}(2x_i^2+1)\text{sh}^1(x_i)\right]/x_i^3,`$ (43)
$`G(x_i){\displaystyle \frac{1}{8}}\left[x_i\sqrt{x_i^2+1}(2x_i^23)+3\text{sh}^1(x_i)\right]/x_i^5,`$ (44)
$`f(x_i){\displaystyle \frac{3}{2}}{\displaystyle \frac{n_b}{m_i}}{\displaystyle \frac{dm_i}{dn_b}}\left[x_i\sqrt{x_i^2+1}\text{sh}^1(x_i)\right]/x_i^3.`$ (45)
One can see, from Eqs. (40) and (41), that an additional term appears in the pressure expression, but not in the energy expression. We can specially confirm this result further as such.
From Eq. (19), one has an alternative general expression for pressure:
$`P`$ $`=`$ $`{\displaystyle \frac{d(VE)}{dV}}|_{S,\{\overline{N_k}\}}`$ (46)
$`=`$ $`E+{\displaystyle \underset{j}{}}n_j{\displaystyle \frac{dE}{dn_j}}|_S.`$ (47)
According to the Pauli principle and the relativistic energy-momentum relation $`\epsilon _i=\sqrt{p^2+m_i^2}`$, the energy density of the system at zero temperature should be
$$E(\{n_i\},\{m_j(n_b)\})=\underset{i}{}\frac{g_i}{2\pi ^2}_0^{p_{f,i}}\epsilon _ip^2𝑑p,$$
(48)
which, when the integration is carried out, is just the same as Eq. (40).
Because the entropy is also zero (or constant) at zero temperature, we can substitute Eq. (48) into Eq. (47), and accordingly get
$`P`$ $`=`$ $`E+{\displaystyle \underset{j}{}}n_j\left({\displaystyle \frac{E}{n_j}}+{\displaystyle \underset{i}{}}{\displaystyle \frac{E}{m_i}}{\displaystyle \frac{m_i}{n_b}}{\displaystyle \frac{n_b}{n_j}}\right)`$ (49)
$`=`$ $`E+{\displaystyle \underset{j}{}}n_j{\displaystyle \frac{E}{n_j}}+{\displaystyle \underset{i}{}}{\displaystyle \underset{j}{}}{\displaystyle \frac{n_j}{3}}{\displaystyle \frac{E}{m_i}}{\displaystyle \frac{m_i}{n_b}}`$ (50)
$`=`$ $`\mathrm{\Omega }+{\displaystyle \underset{i}{}}n_b{\displaystyle \frac{m_i}{n_b}}{\displaystyle \frac{E}{m_i}},`$ (51)
which leads to Eq. (41) exactly.
## III Properties of strange quark matter in the new thermodynamic treatment
Having derived in detail the thermodynamics with variable particle masses in the previous section, we now apply it to the investigation of strange quark matter.
As usually done in the previous literature , We assume the SQM to be a Fermi gas mixture of $`u`$, $`d`$, $`s`$ quarks and electrons with chemical equilibrium maintained by the weak interactions:
$$d,su+e+\overline{\nu }_e,s+uu+d,\mathrm{}$$
Because of these reactions, the chemical potentials $`\mu _i(i=u,d,s,e)`$ should satisfy
$`\mu _d=\mu _s\mu ,`$ (52)
$`\mu _u+\mu _e=\mu .`$ (53)
For the bulk SQM in weak equilibrium, the previous investigations got a slightly positive charge . Our recent study demonstrates that negative charges could lower the critical density. However, too much negative charge can make it impossible to maintain flavor equilibrium. Therefore, the charge of SQM is not allowed to shift too far away from zero at both positive and negative directions. Therefore, one also has another two equations for a given baryon number density $`n_b`$:
$`{\displaystyle \frac{1}{3}}(n_u+n_d+n_s)=n_b,`$ (54)
$`{\displaystyle \frac{2}{3}}n_u{\displaystyle \frac{1}{3}}n_d{\displaystyle \frac{1}{3}}n_sn_e=0.`$ (55)
The first is the definition of baryon number density; the second is from the charge neutrality requirement. $`n_i(i=u,d,s,e)`$ is related to $`\mu _i`$ and $`m_i`$ by Eq. (39).
Because the results from lattice calculations show that quark matter does not become asymptotically free soon after the phase transition (instead, it approaches the free gas equation of state very slowly), one should consider the strong interaction between quarks in a proper way. We do this by including the interaction effect within the variable quark masses. Because of the characteristic of the quark confinement and asymptotic freedom of the strong interaction, one can write down the simplest and most symmetric parametrization for the quark masses $`m_q(q=u,d,s)`$ :
$$m_q=m_{q0}+\frac{D}{n_b^z},$$
(56)
where $`m_{q0}`$ is the corresponding quark current mass, $`z`$ is a fixed exponential. Previously, $`z`$ is regarded as 1. Our recent study indicates that it is more reasonable to take $`z=1/3`$. The parameter $`D`$ is usually determined by stability arguments, i.e., at zero pressure $`(P=0)`$, the energy per baryon, $`E/n_b`$, is great than 930 MeV for two flavor quark matter in order not to contradict standard nuclear physics, but less than 930 MeV for three flavor symmetric quark matter so that SQM can have the possibility of absolute stability. Obviously, the rang of $`D`$ determined by this method depends on different thermodynamic treatments. Within the thermodynamics derived in the preceding section, $`D`$ is in the range (155—171 MeV)<sup>2</sup> when taking $`z=1/3`$.
Because the light quark current masses are very small, their value uncertainties are not important. So we take the fixed central values $`m_{u0}=5`$ MeV and $`m_{d0}=10`$ MeV in our calculation. The electron mass is very small (0.511 MeV). As for $`s`$ quarks, we take 80 and 90 MeV, corresponding respectively to $`D^{1/2}`$ = 156 and 160 MeV.
For a given $`n_b`$, we solve for $`\mu _i(i=u,d,s,e)`$ from Eqs. (52)—(55), and calculate the energy density and pressure of SQM respectively from Eqs. (40) and (41) with the quark messes replace by Eq. (56).
Firstly, we draw the configuration of the SQM for the parameter set $`m_{s0}=80`$ MeV and $`D^{1/2}=156`$ MeV in Fig. 1. At high densities, all of the $`u,d`$, and $`s`$ quarks tend to become a triplicate. When the density becomes lower, $`d`$ fraction will increases while $`s`$ fraction decreases, and will become zero at a definite density which is called critical density in Ref. because SQM does not exist below that density. The $`u`$ fraction is nearly unchanged. It in fact increases very slowly. To keep charge neutrality, the electron fraction will also increase. However, because of its very small mass, the electron fraction is so little that we multiply by one thousand to draw it in the figure.
In Fig. 2, we show the density dependence of the energy per baryon, $`VE/N=E/n_b`$, vs baryon number density $`n_b`$ for the parameter set I: $`m_{s0}=80`$ MeV, $`D^{1/2}=156`$ MeV, and II: $`m_{s0}=90`$ MeV, $`D^{1/2}=160`$ MeV. For the first parameter set, SQM is absolutely stable while for the second set it is nearly meta-stable. The point marked with a circle ‘$``$’ is the zero pressure point where the system pressure becomes zero. It can be clearly seen that the zero pressure points are exactly located at the lowest energy state. In fact, this is a basic requirement of thermodynamics because one can obtain from Eq. (46)
$$P=\frac{d(VE)}{dn_b}\frac{dn_b}{dV}|_{\{\overline{N_k}\}}=n_b^2\frac{d(E/n_b)}{dn_b}.$$
(57)
However, this is not the case for most of the previous thermodynamic treatments of strange quark matter in the mass-density-dependent model , which is their another serious flaw in addition to the unreasonable vacuum limits mentioned before.
In Fig. 3, we give the relation between the pressure $`P`$ and energy density $`E`$, i.e., the equation of state. It approaches the free gas equation of state at high densities. However, its shape is a little sunken at lower densities, contrary to previous calculation which is protuberant. This will leads to completely different lower density behaviour of the sound velocity in strange quark matter.
The velocity of sound is plotted in the lower part of Fig. 4. The upper part is calculated by the same method in Ref. with parameter set B there. Simultaneously given with a full horizontal line is the ultra-relativistic case ($`1/\sqrt{3}`$) for purpose of comparison. Obviously, they become nearly identical at high densities while the lower density behaviour is opposite. The sound velocity in the previous treatment is higher than the ultra-relativistic case and will eventually exceed the speed of light at lower densities, which is unreasonable from the point of view of the theory of relativity.
## IV Strange stars
It has long been proposed that the currently called neutron stars might be composed of strange quark matter and thus be in fact strange stars. A recent investigation shows that young millisecond pulsars are most likely to be strange stars rather than neutron stars . Previous authors have investigated the properties of strange stars by applying their obtained equation of state with interesting results . We have now modified the thermodynamic treatment and updated the quark mass scaling. Therefore, it is meaningful to study the structure of strange stars in the new context from the astrophysical view point.
As generally done, we assume the strange star to be a spherically symmetric object. Its stability is determined by the general relativistic equation of hydrostatic equilibrium known as Tolman-Oppenheimer-Volkov equation
$$\frac{dP}{dr}=\frac{GmE}{r^2}\frac{(1+4\pi r^3P/m)(1+P/E)}{12Gm/r},$$
(58)
with the subsidiary condition
$$dm/dr=4\pi r^2E,$$
(59)
where $`G=6.70710^{45}`$ MeV<sup>-2</sup> is the gravitational constant, $`r`$ is the distance from the core of the star, $`E=E(r)`$ is the energy density or mass density, $`P=P(r)`$ is the pressure, and $`m=m(r)`$ is the mass within the radius $`r`$.
For an initial baryon number density $`n_0`$ (accordingly $`P_0`$ and $`E_0`$), we can numerically solve Eqs. (58) and (59) with the aid of the equation of state, and obtain the corresponding $`P=P(r,n_0)`$ and $`m=m(r,n_0)`$, and consequently $`n=n(r,n_0)`$, the baryon number density at the radius $`r`$ for the central density $`n_0`$. The radius $`R`$ of the strange star is determined by the condition
$$P(R,n_0)=0,$$
(60)
namely,
$$R=R(n_0).$$
(61)
Accordingly, the mass of the strange star is
$$M=m(R(n_0),n_0)M(n_0).$$
(62)
To make strange stars stable, we must require $`dM/dn_0>0`$. For the above obtained equation of state, M firstly increases with $`n_0`$ up to a definite value $`M_{max}`$ corresponding to the highest acceptable central density $`n_{0max}`$. After that, $`M`$ decreases with $`n_0`$, and the star becomes unstable.
For the parameter set I, i.e. $`m_{s0}=80`$ MeV and $`D=(156`$ MeV$`)^2`$, we give the density profiles $`n(r,n_0)`$ in Fig. 5 as an example. The upmost line is for the largest acceptable central density $`n_{0max}`$ $`(1.35`$ fm<sup>-3</sup>). The lowest horizontal line corresponds to the surface density $`n_s`$ $`(0.25`$ fm <sup>-3</sup>) of strange stars which is independent of the central density, but a functional of the equation of state. Each line will intersect with it. The cross points correspond to the radius $`R`$ of the star. The maximum radius of the star appears in $`n_0`$ $`0.65`$ fm<sup>-3</sup>.
In Fig. 6, we show the mass-radius relation of strange stars with a solid line. The point marked with a full dot ‘$``$’ represents the largest acceptable mass $`M_{max}`$ ($`1.58`$ times the solar mass). For comparison, we have also plotted the result from the bag model calculation with the bag constant $`B^{1/4}=144`$ MeV, and that in Re. with the parameter set B there. We can see that the strange stars in our case is dimensionally smaller and less massive than the previous calculation if SQM is absolutely stable. Naturally, this observation depends on the parameters employed. If we choose a bigger $`m_{s0}`$ and larger $`D`$, the case would be different. However, SQM would have no possibility of absolute stability in that case.
## V Summary
We have derived the thermodynamics with density-dependent particle masses self-consistently, which overcome the serious flaws of the previous treatment of SQM in the quark mass-density-dependent model. We find that an additional term should be appended to the expression of pressure, but it should not appear in that of energy. When applying the new formulas to the investigation of SQM, we find that the density behaviour of the sound velocity is opposite to the previous calculation, but consistent with our recent publication, which leads to different structure of strange stars.
## ACKNOWLEDGMENTS
The authors would like to thank the National Natural Science Foundation of China for financial support under Grant No. 19905011.
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# 1 Introduction
## 1 Introduction
The objective of the present lecture is to illustrate the property of the Schrödinger equation which is called here the dyon–oscillator duality. The property is in the following. The Schödinger equation for an oscillator possesses two parameters – the energy $`E`$ and the cyclic frequency $`\omega `$. The quantization leads to the constraint $`E=\mathrm{}\omega (N+D/2)`$ where $`N=0,1,2,\mathrm{}`$, and $`D`$ is the dimension of the configuration space of the oscillator. If $`\omega `$ is fixed, then $`E`$ is quantized and that is the standard situation. Imagine for a moment that now $`E`$ is fixed. Whence, necessarily $`\omega `$ is quantized, and we are in a nonstandard situation. The question is whether the nonstandard situation corresponds to any physics, i.e., whether it is possible to find such a transformation that converts the oscillator into a physical system with a coupling constant $`\alpha `$, being a function of $`E`$, and energy $`\epsilon `$, depending on $`\omega `$. If there exists such a transformation, we can confirm that the ”nonstandard oscillator” is identical to that physical system. Below will be shown the validity of the described picture for dimensions $`D=1,2,4,8`$ and that the final system is a bound system of charge–dyon (remind, that dyon is the hypothetical particle introduced by Schwinger which is unlike the Dirac monopole, endowed with not just magnetic but electric charge as well). As the ”standard” and ”nonstandard” regimes are mutually exclusive, the initial oscillator and the final ”charge–dyon” system are dual to each other, and that explains the relevancy of the term ”dyon–oscillator duality”. Note also that in the initial system the spectrum is discrete only, i.e. the particle has just a finite motion (for such cases it is usually said that we have a model with confinement). Generally speaking, the spectrum of the final system includes the discrete spectrum as well as the continuous one, i.e. in that model there is no confinement. However, unlike the first model, in the second model we have monopoles. There is some analogy between the dyon–oscillator and the Seiberg–Witten duality, according to which the gauge theories with strong interactions are equivalent to the theories having weak interaction on the one hand and topological nontrivial objects, such as monopoles and dyons are, on the other hand.
## 2 Radial Equations
Let consider the equation
$$\frac{d^2R}{du^2}+\frac{D1}{u}\frac{dR}{du}\frac{L(L+D2)}{u^2}R+\frac{2\mu }{\mathrm{}^2}\left(E\frac{\mu \omega ^2u^2}{2}\right)R=0.$$
(1)
Here $`R`$ is the radial part of the wave function for the $`D`$-dimensional oscillator ($`D>2`$) and $`L=0,1,2,\mathrm{}`$ are the eigenvalues of the global angular momentum.
Introduce $`r=u^2`$ and take into account that
$$\frac{1}{u}\frac{d}{du}=2\frac{d}{dr},\frac{d^2}{du^2}=2\frac{d}{dr}+4r\frac{d^2}{dr^2}.$$
Then, equation (1) transforms into
$$\frac{d^2R}{dr^2}+\frac{d1}{r}\frac{dR}{dr}\frac{l(l+d2)}{r^2}R+\frac{2\mu }{\mathrm{}^2}\left(\epsilon +\frac{\alpha }{r}\right)R=0,$$
(2)
where
$$d=D/2+1,l=L/2,$$
(3)
$$\epsilon =\mu \omega ^2/8,\alpha =E/4.$$
(4)
This is quite an unexpected result. If $`D=4,6,8,10,\mathrm{}`$, then $`d=3,4,5,6,\mathrm{}`$, and equation (2) is formally identical to the radial equation for a $`d`$-dimensional hydrogen atom (for odd $`D>2`$ the value of $`d`$ is half-integer and so cannot have the meaning of the dimension of the space in a usual sense). Then, $`l`$ takes not just integer but half-integer values as well, and a question arises about the origin of the fermion degree of freedom. The answer to the question will be given later. Finally, as has been mentioned in the first section, equations (1) and (2) are dual to each other and the duality transformation is $`r=u^2`$.
Earlier, just the radial part of the wave function of the oscillator was considered. For the Schrödinger equation we must take into account the angular part as well. Thus, the duality transformation must also include the transformation of angular variables. If we interpret the change of variables $`r=u^2`$ as a mechanism of generation of electric charge, then (as will be shown later) the transformation of some angular variables is responsible for the generation of magnetic charges.
In the next sections, we study dimensions $`D=1`$ and $`D=2`$ not considered in equation (1). Then, we analyze the dimensions $`D=4`$ and $`D=8`$. The dyon–oscillator duality is limited to these four dimensions. We postpone for a while the discussion of the problem of selection of the dimensions $`D=1,2,4,8`$.
## 3 1D Coulomb Anyon
Consider the one-dimensional Schrödinger equation
$$\frac{d^2\mathrm{\Psi }}{du^2}+\frac{2\mu }{\mathrm{}^2}\left(E\frac{\mu \omega ^2u^2}{2}\right)\mathrm{\Psi }=0,$$
(5)
where $`\mathrm{}<u<\mathrm{}`$. We define a new variable
$$x=u^2,$$
and using the identity
$$\frac{d^2}{du^2}=4|x|\left(\frac{d^2}{dx^2}+\frac{1}{2x}\frac{d}{dx}\right)$$
and setting
$$\mathrm{\Psi }=Cx^{1/4}\mathrm{\Phi },$$
(6)
arrive at the equation
$$\frac{d^2\mathrm{\Phi }}{dx^2}+\frac{2\mu }{\mathrm{}^2}\left(\epsilon _+\frac{\alpha }{|x|}+\frac{\mathrm{}^2}{2\mu }\frac{3}{16x^2}\right)\mathrm{\Phi }=0,$$
(7)
where $`\epsilon `$ and $`\alpha `$ are the same as in (4).
Let us introduce the quantity $`\nu `$ which takes two values: $`\nu =1/4`$ and $`\nu =3/4`$, and rewrite the last equation in the form
$$\frac{d^2\mathrm{\Phi }^{(\nu )}}{dx^2}+\frac{2\mu }{\mathrm{}^2}\left(\epsilon V_cV_{cs}\right)\mathrm{\Phi }^{(\nu )}=0,$$
(8)
where $`V_c=\alpha /|x|`$ and $`V_{cs}`$ is the Calogero–Sutherland potential $`V_{cs}=\mathrm{}^2\nu (1\nu )/2\mu x^2`$.
In one spatial dimension, a particle moving in the Calogero–Sutherland potential has a very unusual property. Unlike the potential $`V_{cs}`$, the wave function is not invariant under the replacement $`\nu (1\nu )`$. It describes a boson for even $`\nu `$ and a fermion for odd $`\nu `$. Statistics corresponding to the other values of $`\nu `$ is called the fractional statistics, and the system influenced along with $`V_{cs}`$ by a potential binding the particle to the center is called the 1D anyon. So, we have started from the 1D quantum oscillator and arrived at the 1D Coulomb anyon.
Comparing Eq. (5) with Eq. (8), we summarize that there are two alternative possibilities connected with Eq. (5) – explicit and hidden. In the first case, the parameter $`\omega `$ is fixed ($`\omega =fix.>0`$) and plays a role of the coupling constant, the parameter $`E`$ is quantized and has the meaning of energy, and the system is the 1D quantum oscillator. For a hidden possibility, the parameter $`E`$ is fixed ($`E=fix.>0`$), the coupling constant is equal to $`E/4`$, $`\omega `$ is quantized, the quantity $`\epsilon =\mu \omega ^2/8`$ takes the meaning of energy, and the system is the 1D Coulomb anyon. Since the 1D Coulomb anyon includes the $`1/x^2`$ interaction, it pretends to be a magnetic monopole in one spatial dimension. So, the anyon–oscillator duality is a prototype of the dyon–oscillator duality in 1D Quantum Mechanics.
Now we can calculate the energy levels $`\epsilon _n`$ and the wave functions $`\mathrm{\Phi }_n^{(\nu )}`$ in the following way. For energy levels we have
$$\epsilon =\frac{\mu \omega ^2}{8}=\frac{\mu }{8}\left[\frac{E}{\mathrm{}(2n+2\nu )}\right]^2=\frac{\mu }{8}\left[\frac{4\alpha }{\mathrm{}(2n+2\nu )}\right]^2=\frac{\mu \alpha ^2}{2\mathrm{}^2(n+\nu )^2},$$
where $`N=2n+2\nu 1/2`$ with $`N`$ numerating the energy levels $`E=\mathrm{}\omega (N+1/2)`$ and $`n`$ being integer and nonnegative.
Consider the wave functions. It follows from (6) that
$$\mathrm{\Phi }_n^{(\nu )}=\frac{1}{C}x^{1/4}\mathrm{\Psi }_n^{(\nu )},$$
where $`\mathrm{\Psi }_n^{(\nu )}\mathrm{\Psi }`$, and therefore
$$\underset{\mathrm{}}{\overset{\mathrm{}}{}}|\mathrm{\Phi }_n^{(\nu )}|^2𝑑x=\frac{1}{|C|^2}\underset{\mathrm{}}{\overset{\mathrm{}}{}}x^{1/2}|\mathrm{\Psi }_n^{(\nu )}|^2𝑑x.$$
The integral in the left-hand side is equal to 1, from which it follows that
$$|C|^2=\underset{\mathrm{}}{\overset{\mathrm{}}{}}u^2|\mathrm{\Psi }_N(u)|^2𝑑u=\overline{u^2}=\frac{2(n+\nu )\mathrm{}}{\mu \omega }.$$
Thus,
$$\mathrm{\Phi }_n^{(\nu )}=\frac{(1)^n}{\sqrt{2}}\sqrt{\frac{\mu \omega }{\mathrm{}(n+\nu )}}x^{1/4}\mathrm{\Psi }_n^{(\nu )}$$
if we choose the phase factor as $`(1)^n`$.
Remind that according to the theory of quantum oscillator,
$$\mathrm{\Psi }_N^{(\nu )}=\left(\frac{\mu \omega }{\pi \mathrm{}}\right)^{1/4}\frac{1}{2^NN!}e^{\mu \omega u^2/2}H_N\left(u\sqrt{\frac{\mu \omega }{\mathrm{}}}\right),$$
where $`H_N(\xi )`$ is the Hermite polynomial
$$H_N(\xi )=(1)^Ne^{\xi ^2}\frac{d^N}{d\xi ^N}e^{\xi ^2}.$$
Further, it is known that Hermite polynomials could be expressed in terms of confluent hypergeometric functions. For our case ($`s=0,1/2`$)
$$H_{2n+2s}(z)=(1)^n\frac{(2n+2s)!}{n!}(2z)^{2s}F(n,2s+1/2,z^2).$$
Using the identification $`y=x\mu \omega /\mathrm{}`$ and the relations $`2s+1/2=2\nu `$ and $`\mu \omega /\mathrm{}=2\mu \alpha /\mathrm{}^2(n+\nu )`$, we get
$$\mathrm{\Phi }_n^{(\nu )}=\sqrt{\frac{\mu \alpha }{\mathrm{}^2}}\frac{1}{2^{n\nu +1/4}}\frac{\sqrt{\mathrm{\Gamma }(2n+2\nu +1/2)}}{\pi ^{1/4}n!(n+\nu )}y^\nu e^{|y|/2}F(n,2\nu ,y),$$
and after taking into account the duplication formula for Euler’s gamma-function
$$\mathrm{\Gamma }(2z)=2^{2z1}\pi ^{1/2}\mathrm{\Gamma }(z)\mathrm{\Gamma }(z+1/2)$$
we arrive at the formula
$$\mathrm{\Phi }_n^{(\nu )}=\frac{\sqrt{\mu \alpha }}{\mathrm{}}\frac{1}{n+\nu }\frac{1}{\mathrm{\Gamma }(2\nu )}\sqrt{\frac{\mathrm{\Gamma }(n+2\nu )}{n!}}y^\nu e^{|y|/2}F(n,2\nu ,y).$$
So, we have two types of 1D Coulomb anyons with $`\nu =1/4`$ and $`\nu =3/4`$, respectively.
## 4 Magnetic Vortex
Now turn to the cyclic oscillator. Here is the first example where along with the radial variable there appears an angular one. In the polar coordinates $`(u,\phi )`$, where $`0u<\mathrm{}`$, $`0\phi <2\pi `$, the Schrödinger equation takes the form
$$\frac{^2\mathrm{\Psi }}{u^2}+\frac{1}{u}\frac{\mathrm{\Psi }}{u}+\frac{1}{u^2}\frac{^2\mathrm{\Psi }}{\phi ^2}+\frac{2\mu }{\mathrm{}^2}\left(E\frac{\mu \omega ^2u^2}{2}\right)\mathrm{\Psi }=0.$$
(9)
Input new variables
$$r=u^2,\varphi =2\phi $$
(10)
and rewrite equation (9) as
$$\frac{^2\mathrm{\Psi }}{r^2}+\frac{1}{r}\frac{\mathrm{\Psi }}{r}+\frac{1}{r^2}\frac{^2\mathrm{\Psi }}{\varphi ^2}+\frac{2\mu }{\mathrm{}^2}\left(\epsilon +\frac{\alpha }{r}\right)\mathrm{\Psi }=0.$$
(11)
where $`\epsilon `$ and $`\alpha `$ are given by expressions (4). Equation (10) is identical to the Schrödinger equation for a two-dimensional hydrogen atom; however, $`\varphi [0,4\pi )`$. Thus, instead of a plane, we have two-sheeted Riemann surface. As a consequence, the single-valuedness condition $`\mathrm{\Psi }(r,\varphi +4\pi )=\mathrm{\Psi }(r,\varphi )`$ leads to the integer as well as half-integer eigenvalues for the angular momentum. The solution of the first type for $`\varphi (\varphi +2\pi )`$ does not change the sign, while the second type solutions under the same transformation change the sign. Whence, without loss of information we can think of $`\mathrm{\Psi }(r,\varphi )`$ defined in the region $`0\varphi <2\pi `$ and having two modifications that differ from each other by the quantum number $`s=0`$ or $`1/2`$. In addition $`\mathrm{\Psi }^{(0)}(r,\varphi +2\pi )=\mathrm{\Psi }^{(0)}(r,\varphi )`$ and $`\mathrm{\Psi }^{(1/2)}(r,\varphi +2\pi )=\mathrm{\Psi }^{(1/2)}(r,\varphi )`$. We say that these wave functions describe the system with full inner momentum $`s=0`$ and $`s=1/2`$, respectively.
Introduce now the important substitution
$$\mathrm{\Psi }^{(s)}(r,\varphi )=e^{is\varphi }\overline{\mathrm{\Psi }}^{(s)}(r,\varphi ),$$
(12)
where $`\overline{\mathrm{\Psi }}^{(s)}(r,\varphi +2\pi )=\overline{\mathrm{\Psi }}^{(s)}(r,\varphi )`$ for $`s=0`$ as well as for $`s=1/2`$. From (11) and (12) it follows that the function $`\overline{\mathrm{\Psi }}^{(s)}`$ satisfies the equation
$$\frac{^2\overline{\mathrm{\Psi }}^{(s)}}{r^2}+\frac{1}{r}\frac{\overline{\mathrm{\Psi }}^{(s)}}{r}+\frac{1}{r^2}\left(\frac{}{\varphi }+is\right)^2\overline{\mathrm{\Psi }}^{(s)}+\frac{2\mu }{\mathrm{}^2}\left(\epsilon +\frac{\alpha }{r}\right)\overline{\mathrm{\Psi }}^{(s)}=0.$$
(13)
Now let us clear up to what system there corresponds equation (13). Input the Cartesian coordinates
$$x_1=r\mathrm{cos}\varphi ,x_2=r\mathrm{sin}\varphi .$$
As $`/\varphi =x_1/x_2x_2/x_1`$, then instead of (13) we have
$$\left(\frac{}{x_1}\frac{isx_2}{r^2}\right)^2\overline{\mathrm{\Psi }}^{(s)}+\left(\frac{}{x_2}+\frac{isx_1}{r^2}\right)^2\overline{\mathrm{\Psi }}^{(s)}+\frac{2\mu }{\mathrm{}^2}\left(\epsilon +\frac{\alpha }{r}\right)\overline{\mathrm{\Psi }}^{(s)}=0.$$
(14)
To this equation there corresponds the Hamiltonian
$$\widehat{H}=\frac{1}{2\mu }\left[\left(\widehat{p}_1\frac{\mathrm{}sx_2}{r^2}\right)^2+\left(\widehat{p}_2+\frac{\mathrm{}sx_1}{r^2}\right)^2\right]\frac{\alpha }{r}.$$
(15)
Input a vector
$$\stackrel{}{A}=\frac{g}{r^2}(x_2,x_1)$$
where $`g=\mathrm{}cs/e`$ and $`e=\sqrt{\alpha }`$. As
$$rot\stackrel{}{A}=\frac{A_2}{x_1}\frac{A_1}{x_2}=g\left[\frac{}{x_1}\left(\frac{}{x_1}\frac{1}{r}\right)+\frac{}{x_2}\left(\frac{}{x_2}\frac{1}{r}\right)\right]=g\left(\frac{^2}{x_1^2}+\frac{^2}{x_2^2}\right)\frac{1}{r}=2\pi g\delta (\stackrel{}{x}),$$
then $`\stackrel{}{A}`$ is the vector potential created by the magnetic vortex of the magnetic charge $`g`$ and placed in the origin of coordinates.
Now instead of (15) we get the Hamiltonian
$$\widehat{H}=\frac{1}{2\mu }\left(\widehat{p}_\mu \frac{e}{c}A_\mu \right)^2\frac{e^2}{r},$$
corresponding to the two-dimensional charge–dyon system. So it is proved that the cyclic oscillator is dual to the charge–dyon system, being a generalization of a usual two-dimensional hydrogen atom.
Let us discuss the correspondence between the cyclic oscillator and the charge–dyon system in detail. It is well-known that in the polar coordinates $`(u,\phi )`$ the energy and wave function of the cyclic oscillator are given by the formulas $`E=\mathrm{}\omega (2n+|M|+1)`$ and
$$\mathrm{\Psi }_{n,M}(u,\phi )=A_{n,M}u^{|M|}e^{\mu \omega u^2/\mathrm{}}F(n,|M|+1,\frac{\mu \omega }{\mathrm{}}u^2)e^{iM\phi },$$
where $`n=0,1,2,\mathrm{}`$, $`M=0,\pm 1,\pm 2,\mathrm{}`$. To even and odd wave functions there correspond even and odd values of $`M`$. Formally, the allowance for parity can be realized by introducing the quantum numbers $`s=0,1/2`$ and $`m=0,\pm 1,\pm 2,\mathrm{}`$, so that $`M=2(m+s)`$.
The energy $`\epsilon `$ is calculated similarly to the one in the previous section
$$\epsilon =\frac{\mu e^4}{2\mathrm{}^2(n+|m+s|+1/2)^2}.$$
(16)
Next, going over to $`r=u^2`$ and $`\varphi =2\phi `$ we have
$$\mathrm{\Psi }_{n,m}^{(s)}=A_{n,m}^{(s)}r^{|m+s|}e^{\mu \omega r/\mathrm{}}F(n,2|m+s|+1,\mu \omega r/\mathrm{})e^{i(m+s)\varphi }.$$
It remains to pass from the two-sheeted Riemann surface $`(0\varphi <4\pi )`$ to the plane $`(0\varphi <2\pi )`$, then take into account (12) and the last formula with the expression $`\mu \omega /\mathrm{}=2\mu e^2/\mathrm{}^2(n+|m+s|+1/2)`$ and, introducing a new variable $`\rho =2\mu e^2r/\mathrm{}^2(n+|m+s|+1/2)`$, write
$$\overline{\mathrm{\Psi }}_{n,m}^{(s)}(\rho ,\varphi )=C_{n,m}^{(s)}\rho ^{|n+m|}e^\rho F(n,2|n+m|+1,\rho )e^{im\varphi },$$
(17)
where the normalization constant $`C_{n,m}^{(s)}`$ is determined by the condition
$$2\pi \underset{0}{\overset{\mathrm{}}{}}\left|\overline{\mathrm{\Psi }}_{n,m}^{(s)}(\rho ,\phi )\right|^2r𝑑r=1.$$
Go back to the transformation (10) and pass there from the polar coordinates $`(r,\varphi )`$ to the Cartesian ones $`(x_1,x_2)`$. Note that $`\varphi [0,4\pi )`$. We have
$`x_1`$ $`=`$ $`r\mathrm{cos}\varphi =u^2[\mathrm{cos}^2(\varphi /2)\mathrm{sin}^2(\varphi /2)]=u_1^2u_2^2,`$
$`x_2`$ $`=`$ $`r\mathrm{sin}\varphi =2u^2\mathrm{sin}^2(\varphi /2)\mathrm{cos}(\varphi /2)=2u_1u_2.`$ (18)
This transformation is known from celestial mechanics as the Levi–Civita transformation. In terms of the complex coordinates $`z=x_1+ix_2`$, $`v=u_1+iu_2`$ it takes the form $`z=v^2`$, i.e., corresponds to the square of the complex variable $`v`$. The Levi–Civita transformation together with the transformations (10) and the $`𝐙_\mathrm{𝟐}`$-reduction compose the duality transformation. Note that $`x=\sqrt{x_1^2+x_2^2}=u^2u_1^2+u_2^2`$. The last expression is known as the Euler’s identity. Thus, the Levi–Civita transformation is bilinear coordinate transformation obeying Euler’s identity. This fact is quite noteworthy from the mathematical point of view, and we will have an opportunity to discuss it.
## 5 Charge–Dyon System
Unlike the two-dimensional space, in four-dimensions there are several types of ”spherical coordinates”. We take the ones used in the theory of symmetrical top
$$u_1+iu_2=u\mathrm{cos}(\beta /2)e^{i(\alpha +\gamma )/2},u_3+iu_4=u\mathrm{sin}(\beta /2)e^{i(\alpha \gamma )/2}.$$
(19)
For $`u=const`$, the position on a sphere is parametrized by the coordinates $`(\alpha ,\beta ,\gamma )`$ that cover the sphere completely when
$$\alpha [0,2\pi ),\beta [0,\pi ),\gamma [0,4\pi ).$$
In the coordinates (19), the length-element and the Laplacian are given by
$`dl^2=du^2+{\displaystyle \frac{u^2}{4}}\left(d\alpha ^2+d\beta ^2+d\gamma ^2+2\mathrm{cos}\beta d\alpha d\gamma \right),`$
$`{\displaystyle \frac{^2}{u_\mu ^2}}={\displaystyle \frac{1}{u^3}}{\displaystyle \frac{}{u}}\left(u^3{\displaystyle \frac{}{u}}\right){\displaystyle \frac{4}{u^2}}\widehat{J}^2,`$
where
$`\widehat{J}^2={\displaystyle \frac{1}{\mathrm{sin}\beta }}{\displaystyle \frac{}{\beta }}\left(\mathrm{sin}\beta {\displaystyle \frac{}{\beta }}\right){\displaystyle \frac{1}{\mathrm{sin}^2\beta }}\left({\displaystyle \frac{^2}{\alpha ^2}}2\mathrm{cos}\beta {\displaystyle \frac{^2}{\alpha \gamma }}+{\displaystyle \frac{^2}{\gamma ^2}}\right).`$
Thus, in terms of the coordinates (19) the isotropic oscillator is described by the equation
$$\frac{^2\mathrm{\Psi }}{u^2}+\frac{3}{u}\frac{\mathrm{\Psi }}{u}\frac{4}{u^2}\widehat{J}^2\mathrm{\Psi }+\frac{2\mu }{\mathrm{}^2}\left(E\frac{\mu \omega ^2u^2}{2}\right)\mathrm{\Psi }=0.$$
The operators $`\widehat{J}^2,\widehat{J}_3=i/\gamma `$, $`\widehat{J}_3^{}=i/\alpha `$ are mutually commuting, and their eigenfunction is represented by the matrix of finite rotations
$$D_{ms}^j(\alpha ,\beta ,\gamma )=e^{im\alpha }d_{ms}^j(\beta )e^{is\gamma }.$$
The explicit form of the function $`d_{ms}^j(\beta )`$ is rather complicated, it can be found in manuals on Quantum Mechanics. It is important that the quantities $`j,m`$ and $`s`$ run the values $`j=0,1/2,1,\mathrm{}`$ and $`m,s=0,\pm 1/2,\pm 1,\mathrm{},\pm j`$.
Now it is clear that the function $`\mathrm{\Psi }`$ should be of the form
$$\mathrm{\Psi }=R(u)D_{ms}^j(\alpha ,\beta ,\gamma ).$$
As the eigenvalues of the operator $`\widehat{J}^2`$ are equal to $`j(j+1)`$, the radial function $`R(u)`$ satisfies the equation
$$\frac{d^2R}{du^2}+\frac{3}{u}\frac{dR}{du}\frac{4j(j+1)}{u^2}R+\frac{2\mu }{\mathrm{}^2}\left(E\frac{\mu \omega ^2u^2}{2}\right)R=0.$$
This equation is solved as follows. First, use the dimensionless variable $`v=au`$ with $`a=(\mu \omega /\mathrm{})^{1/2}`$, and rewrite the last equation as
$$\frac{d^2R}{dv^2}+\frac{3}{v}\frac{dR}{dv}\frac{4j(j+1)}{v^2}R+\lambda Rv^2R=0,$$
where $`\lambda =2\mu E/\mathrm{}^2a^2=2E/\mathrm{}\omega `$. The next step is in the passage to the variable $`\rho =v^2`$ and the equation
$$\frac{d^2R}{d\rho ^2}+\frac{2}{\rho }\frac{dR}{dv}\frac{j(j+1)}{\rho ^2}R+\left(\frac{\lambda }{4\rho }\frac{1}{4}\right)R=0.$$
The analysis of this equation as $`\rho 0`$ and $`\rho \mathrm{}`$ verifies the appropriateness of the substitution
$$R(u)=\rho ^je^{\rho /2}W(\rho )$$
leading to the equation for a confluent hypergeometric function
$$\rho \frac{d^2W}{d\rho ^2}+(2j+2\rho )\frac{dW}{d\rho }(j+1\lambda /4)W=0.$$
A further scenario is usual to any student who masters the course of Quantum Mechanics. The result is
$$W=F(j+1\lambda /4,2j+2;\rho ),$$
$$j+1\lambda /4=n,n=0,1,2\mathrm{}.$$
Concluding, we receive
$`E_N=\mathrm{}\omega (N+2),N=2n+2j=0,1,2,\mathrm{}`$ (20)
$`\mathrm{\Psi }=Const(au)^{2j}e^{a^2u^2/2}F(n,2j+2,a^2u^2)d_{ms}^j(\beta )e^{im\alpha }e^{is\gamma }.`$ (21)
For fixed $`j`$ to the energy level $`E_N`$ there correspond $`(2j+1)^2`$ states (degeneracy by $`m`$ and $`s`$). As $`j=\frac{N}{2},\frac{N}{2}1,\mathrm{}`$, the total degeneracy for the $`N`$-th energy level is
$$g_N=\frac{1}{6}(N+1)(N+2)(N+3).$$
Observe now how the charge–dyon system could be obtained from the four-dimensional oscillator.
Using the variable $`r=u^2`$, we obtain
$`{\displaystyle \frac{^2}{u_\mu ^2}}`$ $`=`$ $`4r\{{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{}{r}}\left(r^2{\displaystyle \frac{}{r}}\right)+{\displaystyle \frac{1}{r^2}}[{\displaystyle \frac{1}{\mathrm{sin}\beta }}{\displaystyle \frac{}{\beta }}\left(\mathrm{sin}\beta {\displaystyle \frac{}{\beta }}\right)+{\displaystyle \frac{1}{\mathrm{sin}^2\beta }}{\displaystyle \frac{^2}{\alpha ^2}}]`$
$`+`$ $`{\displaystyle \frac{1}{r^2\mathrm{sin}^2\beta }}[{\displaystyle \frac{^2}{\gamma ^2}}2\mathrm{cos}\beta {\displaystyle \frac{^2}{\alpha \gamma }}]\}.`$
Thus, the Schrödinger equation
$$\frac{d^2\mathrm{\Psi }}{du_\mu ^2}+\frac{2\mu }{\mathrm{}^2}\left(E\frac{\mu \omega ^2u^2}{2}\right)\mathrm{\Psi }=0$$
gains (in terms of the coordinates (19)) the form
$`{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{}{r}}\left(r^2{\displaystyle \frac{\mathrm{\Psi }}{r}}\right)`$ $`+`$ $`{\displaystyle \frac{1}{r^2}}\left[{\displaystyle \frac{1}{\mathrm{sin}\beta }}{\displaystyle \frac{}{\beta }}\left(\mathrm{sin}\beta {\displaystyle \frac{\mathrm{\Psi }}{\beta }}\right)+{\displaystyle \frac{1}{\mathrm{sin}^2\beta }}{\displaystyle \frac{^2\mathrm{\Psi }}{\alpha ^2}}\right]`$
$`+`$ $`{\displaystyle \frac{1}{r^2\mathrm{sin}^2\beta }}\left[{\displaystyle \frac{^2\mathrm{\Psi }}{\gamma ^2}}2\mathrm{cos}\beta {\displaystyle \frac{^2\mathrm{\Psi }}{\alpha \gamma }}\right]+{\displaystyle \frac{2\mu }{\mathrm{}^2}}\left(\epsilon +{\displaystyle \frac{e^2}{r}}\right)\mathrm{\Psi }=0,`$
with $`e^2=E/4`$, $`\epsilon =\mu \omega ^2/8`$.
Perform a substitution
$$\mathrm{\Psi }(r,\alpha ,\beta ,\gamma )=\overline{\mathrm{\Psi }}^{(s)}(r,\alpha ,\beta )e^{is(\alpha +\gamma )},$$
(22)
where $`s`$ is any real parameter. It is easy to show that the function $`\overline{\mathrm{\Psi }}^{(s)}`$ satisfies the equation
$`{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{}{r}}\left(r^2{\displaystyle \frac{\overline{\mathrm{\Psi }}^{(s)}}{r}}\right)+{\displaystyle \frac{1}{r^2}}\left[{\displaystyle \frac{1}{\mathrm{sin}\beta }}{\displaystyle \frac{}{\beta }}\left(\mathrm{sin}\beta {\displaystyle \frac{\overline{\mathrm{\Psi }}^{(s)}}{\beta }}\right)+{\displaystyle \frac{1}{\mathrm{sin}^2\beta }}{\displaystyle \frac{^2\overline{\mathrm{\Psi }}^{(s)}}{\alpha ^2}}\right]`$
$`+{\displaystyle \frac{2is}{r^2(1+\mathrm{cos}\beta )}}{\displaystyle \frac{\overline{\mathrm{\Psi }}^{(s)}}{\alpha }}{\displaystyle \frac{2s^2}{r^2(1+\mathrm{cos}\beta )}}\overline{\mathrm{\Psi }}^{(s)}+{\displaystyle \frac{2\mu }{\mathrm{}^2}}\left(\epsilon +{\displaystyle \frac{e^2}{r}}\right)\overline{\mathrm{\Psi }}^{(s)}=0.`$ (23)
As in the previous sections, the expression for $`\epsilon `$ is easily found to be
$$\epsilon =\frac{\mu e^4}{2\mathrm{}^2(n+j+1)^2}.$$
The wave function $`\overline{\mathrm{\Psi }}^{(s)}`$ is obtained by comparing formulae (22) and (21), i.e.,
$$\overline{\mathrm{\Psi }}^{(s)}(r,\alpha ,\beta )=C\rho ^je^{\rho /2}F(n,2j+2,\rho )d_{ms}^j(\beta )e^{i(ms)\alpha },$$
where $`\rho =a^2u^2=2\mu e^2r/\mathrm{}^2(n+j+1)`$.
The first line of this equation is nothing but $`^2\overline{\mathrm{\Psi }}^{(s)}/x_j^2`$ where
$$x_1+ix_2=r\mathrm{sin}\beta e^{i\alpha },x_3=r\mathrm{cos}\beta .$$
Then,
$$\frac{2is}{r^2(1+\mathrm{cos}\beta )}\frac{}{\alpha }=\frac{is}{r^2(1+\mathrm{cos}\beta )}\frac{}{\alpha }+\frac{}{\alpha }\frac{is}{r^2(1+\mathrm{cos}\beta )}.$$
Substituting here the formula
$$\frac{}{\alpha }=x_1\frac{}{x_2}x_2\frac{}{x_1},$$
we obtain
$`{\displaystyle \frac{2is}{r^2(1+\mathrm{cos}\beta )}}{\displaystyle \frac{}{\alpha }}`$ $`=`$ $`{\displaystyle \frac{isx_1}{r^2(1+\mathrm{cos}\beta )}}{\displaystyle \frac{}{x_2}}+{\displaystyle \frac{}{x_2}}{\displaystyle \frac{isx_1}{r^2(1+\mathrm{cos}\beta )}}`$
$``$ $`{\displaystyle \frac{isx_2}{r^2(1+\mathrm{cos}\beta )}}{\displaystyle \frac{}{x_1}}{\displaystyle \frac{}{x_1}}{\displaystyle \frac{isx_2}{r^2(1+\mathrm{cos}\beta )}}.`$
Also note that
$$\frac{s^2(x_1^2+x_2^2)}{r^4(1+\mathrm{cos}^2\beta )}\frac{2s^2}{r^2(1+\mathrm{cos}\beta )}=\frac{s}{r^2}.$$
Then, equation (23) can be rewritten as
$`\left({\displaystyle \frac{}{x_1}}+{\displaystyle \frac{isx_2}{r^2(1+\mathrm{cos}\beta )}}\right)^2\overline{\mathrm{\Psi }}^{(s)}`$ $`+`$ $`\left({\displaystyle \frac{}{x_2}}{\displaystyle \frac{isx_1}{r^2(1+\mathrm{cos}\beta )}}\right)^2\overline{\mathrm{\Psi }}^{(s)}+{\displaystyle \frac{^2\overline{\mathrm{\Psi }}^{(s)}}{x_3^2}}`$
$`+`$ $`{\displaystyle \frac{2\mu }{\mathrm{}^2}}\left(\epsilon +{\displaystyle \frac{e^2}{r}}{\displaystyle \frac{\mathrm{}^2}{2\mu }}{\displaystyle \frac{s^2}{r^2}}\right)\overline{\mathrm{\Psi }}^{(s)}=0.`$
This equation is identical to the Pauli equation
$`\left({\displaystyle \frac{}{x_j}}i{\displaystyle \frac{e}{c}}A_j\right)^2\overline{\mathrm{\Psi }}^{(s)}+{\displaystyle \frac{2\mu }{\mathrm{}^2}}\left(\epsilon +{\displaystyle \frac{e^2}{r}}{\displaystyle \frac{\mathrm{}^2}{2\mu }}{\displaystyle \frac{s^2}{r^2}}\right)\overline{\mathrm{\Psi }}^{(s)}=0,`$ (24)
where the vector potential $`A_j`$ is expressed as follows:
$$\stackrel{}{A}=\frac{g\mathrm{sin}\beta }{r(1+\mathrm{cos}\beta )}(\mathrm{sin}\alpha ,\mathrm{cos}\alpha ,0)$$
(25)
with $`g=\mathrm{}cs/e`$.
The vector potential (25) corresponds to the Dirac monopole with the magnetic charge $`g`$. So, equation (24) describes the motion for a charged particle $`e`$ in the field of a dyon with charges $`(e,g)`$. The presence of the charge $`(e)`$ is indicated by the term $`e^2/r`$. The part $`(\mathrm{}^2s^2/2\mu r^2)`$ presents a potential introduced by Goldhaber with the argument of conservation of the angular momentum in scattering of a charged particle from a magnetic monopole. As has been proved by Zwanziger, the addition of such a term makes a problem, corresponding to equation (24), superintegrable.
Thus, we are lucky to ”synthesize” from the isotropic oscillator the bound charge–dyon system. It remains to clear up one important detail. As was shown by Dirac, the introduction of magnetic monopole in Quantum Mechanics leads to the quantization of an electric charge
$$e=\frac{\mathrm{}c}{g}s,s=0,\pm 1/2,\pm 1,\pm 3/2,\mathrm{}.$$
In our approach, the Dirac quantization condition is deduced from formula (22). The transformation $`\gamma (\gamma +4\pi )`$ is identical, as it is seen from the coordinate definition (19). Requiring the single-valuedness for the wave function $`\mathrm{\Psi }(r,\alpha ,\beta ,\gamma )`$, we come to the condition $`s=0,\pm 1/2,\pm 1,\mathrm{}`$ which, together with the formula $`g=\mathrm{}cs/e`$, leads to the quantization of an electric charge.
So, we have shown that the dyon–oscillator duality is valid for the four-dimensional oscillator.
Now focus on the duality transformation. So we have
$`x_1+ix_2`$ $`=`$ $`r\mathrm{sin}\beta e^{i\alpha }=2r\mathrm{sin}(\beta /2)\mathrm{cos}(\beta /2)e^{i\alpha }=2r{\displaystyle \frac{u_1+iu_2}{u}}e^{i(\alpha +\gamma )/2}`$
$``$ $`{\displaystyle \frac{u_3+iu_4}{u}}e^{i(\alpha \gamma )/2}e^{i\alpha }=2(u_1+iu_2)(u_3+iu_4),`$
$`x_3`$ $`=`$ $`r\mathrm{cos}\beta =r[\mathrm{cos}^2(\beta /2)\mathrm{sin}^2(\beta /2)]=r{\displaystyle \frac{u_1^2+u_2^2}{u^2}}r{\displaystyle \frac{u_3^2+u_4^2}{u^2}}`$
$`=`$ $`u_1^2+u_2^2u_3^2u_4^2,`$
or, otherwise,
$`x_1`$ $`=`$ $`2(u_1u_3u_2u_4),`$
$`x_2`$ $`=`$ $`2(u_1u_4+u_2u_3),`$
$`x_3`$ $`=`$ $`u_1^2+u_2^2u_3^2u_4^2.`$
This bilinear transformation satisfies Euler’s condition $`r=u^2`$ and is called the Kustaanheimo–Stiefel transformation. It corresponds to the mapping $`\mathrm{I}\mathrm{R}^4(\stackrel{}{\mathrm{u}})\mathrm{I}\mathrm{R}^3(\stackrel{}{\mathrm{x}})`$ that, along with the formula
$$\gamma =\frac{i}{2}ln\left\{\frac{u_1+iu_2}{u_1iu_2}\frac{u_3iu_4}{u_3+iu_4}\right\}$$
and the ansatz $`\mathrm{\Psi }\overline{\mathrm{\Psi }}^{(s)}`$, composes the duality transformation.
## 6 Magic Numbers
Let us answer the question why the dyon–oscillator duality is valid just for the oscillators with the configuration spaces of dimensions $`D=1,2,4,8`$. We have already mentioned that the duality transformation must satisfy the Euler’s identity
$$(u_1^2+u_2^2+\mathrm{}+u_D^2)^2=x_1^2+x_2^2+\mathrm{}+x_d^2,$$
(26)
where $`d=1`$ for $`D=1`$ and $`d=D/2+1`$ for $`D>1`$. It was proved by Hurwitz that in the cases of $`x_i`$ being a bilinear combination of $`u_i`$, the identity
$$(u_1^2+u_2^2+\mathrm{}+u_D^2)^2=x_1^2+x_2^2+\mathrm{}+x_D^2$$
(27)
is true for $`D=1,2,4,8`$. These magic numbers are directly related to the existence of the four fundamental algebraic structures: real numbers, complex numbers, quaternions and octonions. Putting in (27) $`x_{d+1}=x_{d+2}=\mathrm{}=x_D=0`$, we come to (26).
## 7 Hurwitz Transformation
The question arises, of how to find a transformation converting $`\mathrm{I}\mathrm{R}^8(\mathrm{u})`$ into $`\mathrm{I}\mathrm{R}^5(\mathrm{x})`$, i.e. the transformation with the last of the magic numbers presented above. Begin to write down the transformation in the form
$$x=H(u;D)u.$$
Here $`D`$ is the dimension of the space, $`H`$ is the matrix $`D\times D`$ with the elements $`u_\mu `$, and $`x,u`$ are the $`D`$-dimensional columns composed from $`x_j,u_\mu `$ and, possibly, zeroes. So for the Levi–Civita and Kustaanheimo–Stiefel transformations, we have
$$\left|\begin{array}{c}x_1\\ x_2\end{array}\right|=\left|\begin{array}{cc}u_1& u_2\\ u_2& u_1\end{array}\right|\left|\begin{array}{c}u_1\\ u_2\end{array}\right|,$$
$$\left|\begin{array}{c}x_1\\ x_2\\ x_3\\ 0\end{array}\right|=\left|\begin{array}{cccc}u_3& u_4& u_1& u_2\\ u_4& u_3& u_2& u_1\\ u_1& u_2& u_3& u_4\\ u_2& u_1& u_4& u_3\end{array}\right|\left|\begin{array}{c}u_1\\ u_2\\ u_3\\ u_4\end{array}\right|.$$
The matrices $`H(u;2)`$ and $`H(u;4)`$ have the property
$$H(u;2)H^T(u;2)=u^2E(2),H(u;4)H^T(u;2)=u^2E(4),$$
where $`\mathrm{"}T\mathrm{"}`$ means the sign of transposition, $`E(2)`$ and $`E(4)`$ are the unit matrices. Due to these properties the Euler’s identities are fulfilled. Now, one can easily deduce that the transformation $`\mathrm{I}\mathrm{R}^8(\stackrel{}{\mathrm{u}})\mathrm{I}\mathrm{R}^5(\stackrel{}{\mathrm{x}})`$ must take the form
$$\left|\begin{array}{c}x_0\\ x_1\\ x_2\\ x_3\\ x_4\\ 0\\ 0\\ 0\end{array}\right|=\left|\begin{array}{cccccccc}u_0& u_1& u_2& u_3& u_4& u_5& u_6& u_7\\ u_4& u_5& u_6& u_7& u_0& u_1& u_2& u_3\\ u_5& u_4& u_7& u_6& u_1& u_0& u_3& u_2\\ u_6& u_7& u_4& u_5& u_2& u_3& u_0& u_1\\ u_7& u_6& u_5& u_4& u_3& u_2& u_1& u_0\\ u_1& u_0& u_3& u_2& u_5& u_4& u_7& u_6\\ u_2& u_3& u_0& u_1& u_6& u_7& u_4& u_5\\ u_3& u_2& u_1& u_0& u_7& u_6& u_5& u_4\end{array}\right|\left|\begin{array}{c}u_0\\ u_1\\ u_2\\ u_3\\ u_4\\ u_5\\ u_6\\ u_7\end{array}\right|.$$
Whence it follows that
$`x_0`$ $`=`$ $`u_0^2+u_1^2+u_2^2+u_3^2u_4^2u_5^2u_6^2u_7^2,`$
$`x_1`$ $`=`$ $`2(u_0u_4+u_1u_5u_2u_6u_3u_7),`$
$`x_2`$ $`=`$ $`2(u_0u_5u_1u_4+u_2u_7u_3u_6),`$ (28)
$`x_3`$ $`=`$ $`2(u_0u_6+u_1u_7+u_2u_4+u_3u_5),`$
$`x_4`$ $`=`$ $`2(u_0u_7u_1u_6u_2u_5+u_3u_4).`$
It is easy to prove that for the matrix $`H(u;8)`$ there is a condition
$$H(u;8)H^T(u;8)=u^2E(8)$$
that guarantees the validity of Euler’s identity.
Adding to (28) the transformations
$`\alpha _T`$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{ln}{\displaystyle \frac{(u_0+iu_1)(u_2iu_3)}{(u_0iu_1)(u_2+iu_3)}},`$
$`\beta _T`$ $`=`$ $`2\mathrm{arctan}\left({\displaystyle \frac{u_0^2+u_1^2}{u_2^2+u_3^2}}\right)^{1/2},`$ (29)
$`\gamma _T`$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{ln}{\displaystyle \frac{(u_0iu_1)(u_2iu_3)}{(u_0+iu_1)(u_2+iu_3)}},`$
we obtain a transformation converting $`\mathrm{I}\mathrm{R}^8`$ to the direct product $`\mathrm{I}\mathrm{R}^5𝐒^\mathrm{𝟑}`$ of the space $`\mathrm{I}\mathrm{R}^5(\stackrel{}{\mathrm{x}})`$ and a three-dimensional sphere $`𝐒^\mathrm{𝟑}(\alpha _T,\beta _T,\gamma _T)`$.
## 8 Yang–Coulomb Monopole
In the coordinates (28)-(29)the eight-dimensional isotropic oscillator is described by the equation
$$\frac{1}{2\mu }\left(i\mathrm{}\frac{}{x_j}\mathrm{}A_j^2\widehat{T}_a\right)^2\mathrm{\Psi }+\frac{\mathrm{}^2}{2\mu r^2}\widehat{T}^2\mathrm{\Psi }\frac{e^2}{r}\mathrm{\Psi }=\epsilon \mathrm{\Psi }$$
(30)
where $`\epsilon `$ and $`e^2`$ are defined as usual. The operators $`\widehat{T}_a`$ are the generators of the $`SU(2)`$ group. In the coordinates $`(\alpha _T,\beta _T,\gamma _T)`$ they are parametrized as follows:
$`\widehat{T}^1`$ $`=`$ $`i\left(\mathrm{cos}\alpha _T\mathrm{cos}\beta _T{\displaystyle \frac{}{\alpha _T}}+\mathrm{sin}\alpha _T{\displaystyle \frac{}{\beta _T}}{\displaystyle \frac{\mathrm{cos}\alpha _T}{\mathrm{sin}\beta _T}}{\displaystyle \frac{}{\gamma _T}}\right),`$
$`\widehat{T}^2`$ $`=`$ $`i\left(\mathrm{sin}\alpha _T\mathrm{cot}\beta _T{\displaystyle \frac{}{\alpha _T}}\mathrm{cos}\alpha _T{\displaystyle \frac{}{\beta _T}}{\displaystyle \frac{\mathrm{sin}\alpha _T}{\mathrm{sin}\beta _T}}{\displaystyle \frac{}{\gamma _T}}\right),`$
$`\widehat{T}^3`$ $`=`$ $`i{\displaystyle \frac{}{\alpha _T}}.`$
The five-dimensional vectors $`\stackrel{}{A}^a`$ are given by the expressions
$`\stackrel{}{A}^1`$ $`=`$ $`{\displaystyle \frac{1}{r(r+x_0)}}(0,x_4,x_3,x_2,x_1),`$
$`\stackrel{}{A}^2`$ $`=`$ $`{\displaystyle \frac{1}{r(r+x_0)}}(0,x_3,x_4,x_1,x_2),`$
$`\stackrel{}{A}^3`$ $`=`$ $`{\displaystyle \frac{1}{r(r+x_0)}}(0,x_2,x_1,x_4,x_3).`$
Each term of the triplet $`A_j^a`$ coincides with the vector potential of a $`5D`$ Dirac monopole with a unit topological charge and with the line of singularity along the nonpositive $`x_0`$ semiaxis. The vectors $`A_j^a`$ are orthogonal to each other
$$A_j^aA_j^b=\frac{1}{r^2}\frac{rx_0}{r+x_0}\delta _{ab}$$
and to the vector $`\stackrel{}{x}=(x_0,x_1,x_2,x_3,x_4)`$ as well.
We see that equation (4) describes the charge–dyon system with $`SU(2)`$ monopoles which we call the Yang–Coulomb monopole (YCM). The YCM is defined as a five-dimensional system composed of the Yang monopole ($`A_j^a`$) of the topological charge $`+1`$ and the particle of the isospin $`(\widehat{T}_a)`$. Both the monopole and particle are also assumed to have electric charges of the opposite signs. Thus, the monopole–particle coupling is realized not only by the $`SU(2)`$ gauge field but also by the Coulomb interaction. At large distances the Coulomb structure becomes immaterial and YCM seems to be a pure Yang monopole. The YCM is a unique example of an integrable non-Abelian system. The $`SO(6)`$ group is a group of hidden symmetry of YCM which can be used for calculation of the energy spectrum of YCM by an algebraic method.
After quite complicated calculations, which are omitted here, we can reduce equation (30) to the form
$$\left(\mathrm{\Delta }_5\frac{4}{r(r+x_0)}\widehat{L}\widehat{T}\frac{2}{r(r+x_0)}\widehat{T}^2\right)\mathrm{\Psi }+\frac{2\mu }{\mathrm{}^2}\left(\epsilon +\frac{e^2}{r}\right)\mathrm{\Psi }=0,$$
(31)
where
$`\widehat{L}_1`$ $`=`$ $`{\displaystyle \frac{i}{2}}[D_{41}(x)+D_{32}(x)],`$
$`\widehat{L}_2`$ $`=`$ $`{\displaystyle \frac{i}{2}}[D_{13}(x)+D_{12}(x)],`$
$`\widehat{L}_3`$ $`=`$ $`{\displaystyle \frac{i}{2}}[D_{12}(x)+D_{34}(x)]`$
with
$$D_{ij}=x_i\frac{}{x_j}+x_j\frac{}{x_i}.$$
We see that equation (31) contains the LT-coupling term demonstrating that we have no way to separate the wave function dependence on $`\mathrm{I}\mathrm{R}^5`$ and $`𝐒^\mathrm{𝟑}`$.
In $`\mathrm{I}\mathrm{R}^5`$ we introduce the hyperspherical coordinates $`r[o,\mathrm{})`$, $`\theta [0,2\pi ]`$, $`\alpha [0,2\pi )`$, $`\beta [0,\pi ]`$ and $`\gamma [0,4\pi )`$ according to the relations
$`x_0`$ $`=`$ $`r\mathrm{cos}\theta ,`$
$`x_1+ix_2`$ $`=`$ $`r\mathrm{sin}\theta \mathrm{cos}{\displaystyle \frac{\beta }{2}}e^{i\frac{\alpha +\gamma }{2}},`$
$`x_3+ix_4`$ $`=`$ $`r\mathrm{sin}\theta \mathrm{sin}{\displaystyle \frac{\beta }{2}}e^{i\frac{\alpha \gamma }{2}},`$
and rewrite equation (31) as
$$\left(\mathrm{\Delta }_{r\theta }\frac{\widehat{L}^2}{r^2\mathrm{sin}^2(\theta /2)}\frac{\widehat{J}^2}{r^2\mathrm{cos}^2(\theta /2)}\right)\mathrm{\Psi }+\frac{2\mu }{\mathrm{}^2}\left(\epsilon +\frac{e^2}{r}\right)\mathrm{\Psi }=0,$$
(32)
where $`\widehat{J}_a=\widehat{L}_a+\widehat{T}_a`$ and
$$\mathrm{\Delta }_{r\theta }=\frac{1}{r^4}\frac{}{r}\left(r^4\frac{}{r}\right)+\frac{1}{r^2\mathrm{sin}^3\theta }\frac{}{\theta }\left(\mathrm{sin}^3\theta \frac{}{\theta }\right).$$
We introduce the separation ansatz
$$\mathrm{\Psi }=\mathrm{\Phi }(r,\theta )G(\alpha ,\beta ,\gamma ;\alpha _T,\beta _T,\gamma _T),$$
where $`G`$ are the eigenfunctions of $`\widehat{L}^2,\widehat{T}^2`$ and $`\widehat{J}^2`$ with the eigenvalues $`L(L+1)`$, $`T(T+1)`$ and $`J(J+1)`$.
Because of the LT–interaction, we seek the function $`G`$ in the form
$$G=\underset{M=m+t}{}(JM|L,m;T,t)D_{mm^{}}^L(\alpha ,\beta ,\gamma )D_{tt^{}}^T(\alpha _T,\beta _T,\gamma _T),$$
where $`(JM|L,m;T,t)`$ are the Clebsch–Gordan coefficients.
Let us take the function $`\mathrm{\Phi }(r,\theta )`$ in the form
$$\mathrm{\Phi }(r,\theta )=R(r)Z(\theta ).$$
Equation (32) is then separated into
$$\frac{1}{\mathrm{sin}^3\theta }\frac{d}{d\theta }\left(\mathrm{sin}^3\theta \frac{dZ}{d\theta }\right)\frac{2L(L+1)}{1\mathrm{cos}\theta }Z\frac{2J(J+1)}{1+\mathrm{cos}\theta }Z+\lambda (\lambda +3)Z=0$$
(33)
and a purely radial equation
$$\frac{1}{r^4}\frac{d}{dr}\left(r^4\frac{dR}{dr}\right)\frac{\lambda (\lambda +3)}{r^2}+\frac{2\mu }{\mathrm{}^2}\left(\epsilon +\frac{e^2}{r}\right)R=0$$
(34)
with the separation constant $`\lambda (\lambda +3)`$ being equal to the nonnegative eigenvalues of the global angular momentum.
In equation (33) it is convenient to change the variable as $`y=(1\mathrm{cos}\theta )/2`$ and set
$$Z(y)=y^L(1y)^JW(y).$$
Substituting this into equation (33) we obtain the hypergeometric equation
$$y(1y)\frac{d^2W}{dy^2}+[c(a+b+1)y]\frac{dW}{dy}abW=0,$$
where $`a=\lambda +L+J`$, $`b=\lambda +L+J+3`$, $`c=2L+2`$.
Thus, we find that
$$Z(\theta )=(1\mathrm{cos}\theta )^L(1+cos\theta )^LF(\lambda +J+L,\lambda +J+L+3,2L+2;\frac{1\mathrm{cos}\theta }{2}).$$
The solution behaves well at $`\theta =\pi `$ if the series $`F`$ terminates, that is
$$\lambda +J+L=n_\theta ,$$
with $`n_\theta =0,1,2,\mathrm{}`$.
Let us now consider the radial equation and introduce the function
$$f(r)=e^{kr}r^\lambda R(r).$$
It can easily be verified that the equation for $`f(r)`$ has the form of the confluent hypergeometric equation
$$z\frac{d^2f}{dz^2}+(cz)\frac{df}{dz}af=0,$$
where $`z=2kr,k=\sqrt{2\mu \epsilon /\mathrm{}^2}`$, $`c=2\lambda +4`$, $`a=\lambda +21/kr_0`$, $`r_0=\mathrm{}^2/me^2`$. For the bound state solutions ($`\epsilon <0`$), we have
$$\lambda +21/kr_0=n_r,n_r=0,1,2,\mathrm{};$$
therefore,
$$\epsilon _N^T=\frac{me^4}{2\mathrm{}^2(N/2+2)^2},$$
where $`N=2(n_r+\lambda )=2(n_r+n_\theta +J+L)`$.
For fixed $`T`$, the energy levels $`\epsilon _N^T`$ do not depend on $`L,J`$ and $`\lambda `$, i.e., they are degenerate. The total degeneracy is
$$g_N^T=(2T+1)\underset{\lambda }{}\underset{L}{}(2L+1)\underset{J}{}(2J+1).$$
After some tedious calculation we finally obtain
$$g_N^T=\frac{1}{12}(2T+1)^2\left(\frac{N}{2}T+1\right)\left(\frac{N}{2}T+2\right)\left\{\left(\frac{N}{2}T+2\right)\left(\frac{N}{2}T+3\right)+2T(N+5)\right\}.$$
For $`T=0`$ and $`N=2n`$ (even) the right-hand side of the last formula is equal to $`(n+1)(n+2)^2(n+3)/12`$ – that is, to the degeneracy of pure Coulomb levels. Further, we have $`T=0,1,\mathrm{},N/2`$ for even $`N`$ and $`T=1/2,3/2,\mathrm{},N/2`$ for odd $`N`$. Therefore,
$$g_N=\underset{T=0,\frac{1}{2}}{\overset{N/2}{}}g_N^T=\frac{(N+7)!}{7!N!}$$
i.e., we obtain the degeneracy of the energy levels for the 8D isotropic quantum oscillator.
Formulae (28) and (29) represent the duality transformation mapping the 8D quantum oscillator into charge–dyon system with the $`SU(2)`$ monopole.
## 9 Oscillator-like Systems
We have considered above the dyon–oscillator duality. This type of duality is valid not only for the $`1D,2D,4D`$ and $`8D`$ oscillators, but also for oscillator-like systems with the potentials
$$V(u^2)=C_0+C_2u^2+W(u^2),$$
where $`W(u^2)`$ has the form
$$W(u^2)=\underset{n=2}{\overset{\mathrm{}}{}}C_{2n}u^{2n}$$
For such modified potentials the ansatz (4) can be rewritten as
$$\epsilon =\frac{C_2}{4},e^2=\frac{EC_0}{4}$$
Thus, the value of the function $`V(u^2)`$ at $`u^2=0`$ contributes to the Coulomb coupling constant $`e^2`$. It is also easy to verify that the left-hand side of equation (34) develops the additional term $`(W(r)/4r)`$.
## 10 Exercises
$``$ Find the energy levels and normalized wave functions for states of the particle placed in the field $`V(x)=\alpha /|x|\mathrm{}^2\nu (\nu 1)/2\mu x^2`$, where $`x(\mathrm{},\mathrm{})`$ and $`\nu 0,1/2`$ (see Ref. ).
$``$ Prove that the 3D oscillator with coordinates confined by the 2D half-up cone for an angle of $`\pi /6`$ is dual to the 2D charge–dyon system obeying fractional statistics. Find the duality transformation (see Ref. ).
$``$ Calculate the length-element $`dl^2`$, metric tensor $`g_{\mu \nu }`$ and the Laplace operator $`^2/u_\mu ^2`$ in the coordinates (19).
$``$ Compute the integrals of motion for the 3-dimensional charge–dyon system, transforming for $`g=0`$ into the operator of orbital momentum and Runge–Lenz operator (see Ref. ).
$``$ Prove that the Goldhaber correction in the Hamiltonian of the 3-dimensional charge–dyon system is identical to the interaction $`\stackrel{}{\mu }\stackrel{}{B}`$ of the magnetic momentum $`\stackrel{}{\mu }`$ of a particle with the magnetic field $`\stackrel{}{B}`$ (see Ref. ).
$``$ Solve the Schrödinger equation for the 3-dimensional charge–dyon system in the parabolic coordinates $`x_1=\sqrt{\xi \eta }\mathrm{cos}\phi `$ $`x_2=\sqrt{\xi \eta }\mathrm{sin}\phi `$, $`x_3=(\xi \eta )/2`$ (see Ref. ).
$``$ Compute the expansion coefficients of the parabolic basis of the 3-dimensional charge–dyon system in terms of its spherical basis (see Ref. ).
$``$ Prove that the transformation (28) converts the Schrödinger equation for the $`8D`$ oscillator into equation (30).
$``$ Show that equation (30) can be transformed into (31).
$``$ Calculate the length-element $`dl^2`$, metric tensor $`g_{ij}`$ and Laplace operator $`^2/x_j^2`$ in the coordinates $`(r,\theta ,\alpha ,\beta ,\gamma )`$.
Acknowledgment
I would like to thank the organizers of the school for the occasion to give a talk for so young and active audience.
Also I am grateful to Yeranuhi Hakobyan for her help in preparing this lecture.
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# PRIMORDIAL FRACTAL DENSITY PERTURBATIONS AND STRUCTURE FORMATION IN THE UNIVERSE : 1-DIMENSIONAL COLLISIONLESS SHEET MODEL
## 1 INTRODUCTION
The present Universe shows a variety of structures. The galaxies are not distributed randomly in the Universe. Totsuji & Kihara (1969) and Peebles (1974) showed that the observed two-point correlation function $`\xi (r)`$ is given by a power law with respect to a distance $`r`$ as $`\xi (r)r^\gamma `$ with $`\gamma 1.8`$. The recent galaxy surveys also agree with this result, i.e. the power $`\gamma `$ is nearly equal to $`1.8`$ (CfA (Geller & Hachra (1989)), LCRS (Jing, Mo, & Börner (1998)), and ESP (Guzzo et al. (1998, 1999)). This may imply that the present distribution of galaxies is fractal. Sylos Labini, Montuori, & Pietronero (1998) have also claimed that all available data are consistent with a fractal structure with the dimension $`D2`$ up to the deepest observed scale ($`1000h^1`$\[Mpc\]). However, the observation of Cosmic Microwave Background Radiation (CMBR) has revealed that the Universe in the recombination era is homogeneous and isotropic at least in very large scale. Although CMBR observation seems to be more reliable, we should not decide yet whether the large scale structure of the Universe is really fractal up to the horizon scale or not. To answer this question more definitely, we should await forthcoming next galaxy survey projects (Colless (1995), Maddox (1997), Loveday & Pier (1998), Knapp et al. (1999)).
However, since it seems true that the galaxy distribution is really fractal up to a certain scale, one may ask how such a structure is formed in the evolution of the Universe. One of the most plausible explanations is that the nonlinear dynamics of the perturbations will provide such a scale-free structure during the evolution of the Universe. The pioneering work to explain the power-law behavior in nonlinear stage has been done by Davis & Peebles (1977). They assume a self-similar evolution of density fluctuation and some additional condition, i.e. a physical velocity $`\dot{r}`$ vanishes in nonlinear regime. Then they showed a relation between the power index $`\gamma `$ of two-point correlation function and that of initial power spectrum $`n`$ as $`\gamma =3(n+3)/(n+5)`$. If we have $`n=0`$, then we find that $`\gamma =1.8`$. Since their additional condition is not trivial and might not be appropriate, Padmanabhan (1996) and Yano & Gouda (1998a) extended their model to the case with non-vanishing $`\dot{r}`$. They found that the relation between $`n`$ and $`\gamma `$ is $`\gamma =[3h(n+3)]/[2+h(n+3)]`$, where, $`ha\dot{x}/\dot{a}x`$, which is a ratio of a peculiar velocity to the Hubble expansion. With this result, $`\gamma `$ can vary from 0 to 2 for $`n=1`$ (Harrison-Zel’dovich spectrum) and $`0h1`$ ($`h=1`$ corresponds to the Davis-Peebles solution). Since we do not know the stability of those solutions, in order to find which value of $`\gamma `$ is most likely, we should study the dynamics of density fluctuations in other methods, e.g. $`N`$-body simulation. Several groups in fact showed that a power-law behavior in two-point correlation function is obtained by $`N`$-body simulation with appropriate primordial density fluctuations (Miyoshi & Kihara (1975), Efstathiou (1979), Aarseth, Gott III, & Turner (1979), Frenk, White, & Davis (1983), Davis et al. (1985), Jing (1998)).
The question is whether those power-law behaviors mean that we have a fractal structure in the present Universe. Peebles (1985) and Couchman & Peebles (1998) showed how to proceed with a high resolution analysis in the $`N`$-body simulation using a kind of renormalization method. They have used Davis-Peebles solution as a scaling relation. Without such an ansatz, we do not know whether usual $`N`$-body simulation is suitable to discuss the formation of a fractal structure. With the present state of computers, it may not be possible to obtain high enough resolution to analyze a fractal structure.
As for a fractal structure in the Universe, one may ask another question. Did the Universe not have any non-trivial structure such as a fractal in the initial density fluctuations? In the conventional approaches, initial density perturbations are usually assumed to be given by a power-law (or a power-law-like) spectrum with random Gaussian phase. Although such initial conditions may provide the presently observed nonlinear scale-free structure via nonlinear dynamics, no one has shown whether such a structure is fractal or not, and if yes, what kind of fractal structure comes out. To provide a fractal structure in the present Universe, we may adopt an alternative scenario, in which primordial density fluctuations have already a fractal-like structure in the beginning. Note that a background spacetime is assumed to be a smooth universe, which is described by the Einstein-de Sitter universe, but not a fractal universe. The properties of an initial fractal may be preserved during the evolution of the Universe, then nonlinear fractal structure will be formed. In fact, De Gouveia Dal Pino et al. (1995) reported that the temperature fluctuation of CMBR has a fractal relation, and recently, Pando & Fang (1998) and Feng & Fang (2000) also reported that non-Gaussianity was detected in the distribution of Ly$`\alpha `$ forest lines in the QSO absorption spectra. In this scenario, several natural questions may arise. How does such a primordial fractal perturbation evolve into nonlinear regime? Will any properties of the initial fractal be preserved during the evolution of the Universe, or not? If not, what kind of nonlinear structure will come out at present? Is there any fundamental difference in the structure formation process between a conventional density perturbation and the present fractal one? In order to answer those questions, we study the time evolution of the initial density fluctuations with a fractal structure in Einstein de-Sitter universe.
Since we are interested in a fractal structure, a quite high resolution is required in our calculation. As we discussed, $`N`$-body simulation may not have enough resolution in the present state of computer development, unless we develop some skillful method. So, in this paper, we consider only a very simple toy model, which is a one-dimensional (1-D) sheet model, in order to get some insight into the questions raised in the above. To set up primordial fractal density perturbations, we distribute $`N`$ sheets initially by some systematic rule, i.e. we apply a Cantor set or random Cantor-type set (see below). Mathematically, in order to construct a Cantor set, the procedure must be repeated an infinite number of times, but it is not practically possible to set up such initial data. We therefore stop the procedure at a certain point, i.e. the initial set is given by several times removing line segments with a given ratio (Falconer (1990)). This could be justified because an infinite scale-free structure never exists in the real Universe. In order to construct the initial density perturbations, we set that the remaining segments have small positive density perturbations, while the removed ones correspond to small negative ones. Since we study a 1-D sheet-model, the motion of each sheet is described by an analytic solution (Zel’dovich (1970)), which guarantees enough resolution to analyze a fractal structure.
In $`\mathrm{\S }2`$, we present our formalism and initial setting. As for the initial data, we consider three cases: regular Cantor set, random Cantor-type set, and random white noise. Comparing those time evolutions, we show our results in $`\mathrm{\S }3`$. In $`\mathrm{\S }4`$, we focus particularly on the phase space. The conclusion and discussion follow in $`\mathrm{\S }5`$.
## 2 FORMALISM AND INITIAL DATA
### 2.1 Dynamical Equations
In order to study the structure formation of the expanding Universe, there are so far three approaches: $`N`$-body simulation, the Eulerian perturbation approach and the Lagrangian one. Although the final answer for structure formation would be obtained by $`N`$-body simulation, it may not be possible to so far answer the questions about a fractal structure. As for the perturbation approaches, these are just an approximation and will break down in nonlinear regime, although the Lagrangian approach would be better if we are interested in the density perturbations. This is just because a density fluctuation $`\delta `$ and a peculiar velocity $`𝐯`$ are perturbed quantities in the Eulerian approach (Peebles (1980)), while displacement of particles is assumed to be small in the Lagrangian approach (Zel’dovich (1970), Bouchet (1992, 1995), Coles & Lucchin (1995), Catelan (1995)). Its first order solution is the so-called Zel’dovich approximation (Zel’dovich (1970)). The Lagrangian approach is confirmed to be better than the Eulerian approach by comparison of these results in several cases (Munshi, Sahni, & Starobinsky (1994), Sahni & Shandarin (1996), Yoshisato, Matsubara, & Morikawa (1998)).
The perturbation variable $`𝐒`$ in the Lagrangian approach describes a displacement of dust particles from a uniform distribution and is defined as:
$`𝐱`$ $`=`$ $`𝐪+𝐒(q,t),`$ (1)
$`𝐱`$ $`:`$ $`\mathrm{Eulerian}\mathrm{comoving}\mathrm{coordinate},`$
$`𝐪`$ $`:`$ $`\mathrm{Lagrangian}\mathrm{comoving}\mathrm{coordinate}.`$
The density fluctuation is given by the Jacobian $`J`$ as
$`\delta (𝐪,t)`$ $`=`$ $`{\displaystyle \frac{1J}{J}}`$ (2)
where $`J\mathrm{det}\left(𝐱/𝐪\right)`$.
Even though the Lagrangian approach is better, it is still an approximation and then is not suitable to discuss a highly nonlinear structure formation. However, there is one exceptional case. If the distribution is plane symmetric, the system is one-dimensional, and then the Zel’dovich approximation turns out to be an exact solution. Hereafter, we discuss only one-dimensional problem, which Lagrangian perturbation is given by
$`x=q+S(q,t)`$ (3)
where $`x`$ and $`q`$ are the one-dimensional Eulerian and Lagrangian comoving coordinates, respectively. For Einstein de-Sitter model, the solution is given by (Gouda & Nakamura (1989), Bouchet et al. (1995)):
$$S(q,t)=a(t)S_1(q)+a(t)^{3/2}S_2(q),$$
(4)
where the scale factor $`a`$ changes as $`a=(t/t_0)^{2/3}`$. Then we find the position and its velocity of a dust particle at a scale factor $`a(t)`$ with respect to the Lagrangian coordinate $`q`$ as
$`x(q,a)`$ $`=`$ $`q+a(t)S_1(q)+a(t)^{3/2}S_2(q),`$ (5)
$`\stackrel{~}{v}(q,a)`$ $`=`$ $`S_1(q){\displaystyle \frac{3}{2}}a(t)^{5/2}S_2(q),`$ (6)
where we have introduced a new peculiar velocity $`\stackrel{~}{v}x/a`$. In what follows, we use a scale factor $`a(t)`$ as a time coordinate instead of the physical time $`t`$.
Although the Zel’dovich solution is exact, in studying a formation of nonlinear structure, a serious problem will soon arise. As a density fluctuation grows, we find a shell crossing. For a realistic matter fluid, a pressure may prevent such a singularity from forming. Then the solution will no longer describe the evolution of perturbations after a shell crossing. However, we may have another choice. If instead of a usual matter fluid, we have a collisionless particle such as some dark matter, we can go beyond the shell crossing. The particles, which are described by plane parallel sheets, will pass through each other without collision. Then after this crossing, we rediscover that the Zel’dovich solution is again exact. Therefore, we have a series of exact solutions, which is almost analytic (Gouda & Nakamura (1989), Yano & Gouda (1998b)). We shall call it the 1-D sheet model.
To be more precise, when a crossing by two sheets has occurred, those two sheets exchange their numbering as follows. Suppose there are two sheets $`q_1,q_2(q_1<q_2)`$ with the Eulerian coordinates $`x(q_1,a),x(q_2,a)(x(q_1,a)<x(q_2,a))`$ and with velocities $`\stackrel{~}{v}(q_1,a),\stackrel{~}{v}(q_2,a)`$, respectively. Assume those sheets cross over at $`a_{\mathrm{cross}}`$, i.e. $`x(q_1,a_{\mathrm{cross}})=x(q_2,a_{\mathrm{cross}})`$ with $`\stackrel{~}{v}(q_1,a_{\mathrm{cross}})>\stackrel{~}{v}(q_2,a_{\mathrm{cross}})`$. For the evolution after a shell crossing ($`a>a_{\mathrm{cross}}`$), exchanging their numbering ($`q_1\text{ }\stackrel{}{}q_2`$) as
$`x(q_1,a_{\mathrm{cross}})`$ $``$ $`x(q_2,a_{\mathrm{cross}}),\stackrel{~}{v}(q_1,a_{\mathrm{cross}})\stackrel{~}{v}(q_2,a_{\mathrm{cross}}),`$
$`x(q_2,a_{\mathrm{cross}})`$ $``$ $`x(q_1,a_{\mathrm{cross}}),\stackrel{~}{v}(q_2,a_{\mathrm{cross}})\stackrel{~}{v}(q_1,a_{\mathrm{cross}}),`$ (7)
we find again a natural ordering between the Lagrangian and Eulerian coordinates, i.e. $`x(q_1,a)<x(q_2,a)`$ for $`q_1<q_2`$. By this exchange, we obtain a new distribution of sheets ($`x(q,a),\stackrel{~}{v}(q,a)`$) just after a shell crossing. Using this distribution as an initial data, we find a next exact time evolution of the system by Zel’dovich solution. In order to fix the initial data, we have to determine $`S_1(q)`$ and $`S_2(q)`$ in (4) for the given distribution ($`x(q,a),\stackrel{~}{v}(q,a)`$). From (5,6) we find the solution is
$`S_1(q)`$ $`=`$ $`{\displaystyle \frac{3}{5a}}(xq)+{\displaystyle \frac{2}{5}}\stackrel{~}{v},`$
$`S_2(q)`$ $`=`$ $`{\displaystyle \frac{2}{5}}a^{3/2}\left\{(xq)a\stackrel{~}{v}\right\}.`$ (8)
The new exact solution (4) with (8) is valid until we encounter next shell crossing. another two sheets will cross over.
We repeat this prescription every time when we encounter a shell crossing. As a result, we obtain a series of the Zel’dovich’s exact solutions, which is regarded as an analytic solution for the 1-D collisionless sheet model. Note that this prescription is still valid in multi-stream region. Using this prescription, Gouda & Nakamura (1989) investigated a time evolution for the density perturbations with a scale-free initial power spectrum. They showed that the power spectrum will approach some characteristic value independent of the power index of the initial spectrum. This characteristic power index -1 is predicted by a catastrophe theory. Recently, Yano & Gouda (1998b) investigated a time evolution of the density perturbations for an initial power spectrum with a cutoff. In this case, they found that a self-similarity in all scales is no longer valid. The spectrum is classified into five ranges by its power index. Some spectra coincide with the above one for some scale ranges, but another power index, which is independent of the initial power spectrum, appears just beyond the cut-off scale. In $`\mathrm{\S }5`$, we show our result comparing with their result.
Since this 1-D sheet model is powerful enough to see the fine structure, we shall use this model to analyze a time evolution of the primordial fractal density perturbations.
### 2.2 Setting Up Initial Data
Since we are interested in initial density fluctuations with a fractal distribution, we have to construct such an initial data. For the sake of simplicity, we apply a Cantor set, or a random Cantor-type set (see below), in our construction. A Cantor set is given by the following procedure. We first divide a line with the length $`L`$ by some integer $`n_D`$ and then remove one line segment at the center. If $`n_D`$ is an even integer, not dividing a line by $`n_D`$, we just remove a line segment with the length $`L/n_D`$ from the center of the line. We then repeat this procedure for the remained line segments. In mathematical definition, the removal procedure must be repeated infinite times. However we believe that a fractal structure, even if it exists in the Universe, is not a mathematical one but its self-similarity may end at some scale. Then we also stop our procedure after a finite number of repetitions. We regard the remaining line segment as a region of positive density fluctuations ($`\delta _+`$ region), while the removed part corresponds to a region of negative density fluctuations ($`\delta _{}`$ region). $`\delta _+(>0)`$ and $`\delta _{}(<0)`$ are chosen to be uniform, i.e. both $`\delta _+`$ and $`\delta _{}`$ are some constants in all regions such that $`\delta _+|\delta _{}|10^3`$. Although this is very artificial and the realistic perturbations may depend on each scale just like conventional density perturbations, we analyze only this simplest case here (see Discussion). In addition to the Cantor set, in order to see the universality of our results, we also consider a random Cantor-type set as well as a distribution constructed by a white noise. We describe in more detail how to construct the initial data for each case.
#### 2.2.1 Regular Cantor set
In this case, we consider seven initial data. Each data is constructed by removing central line segments with a fixed ratio ($`1/n_D(n_D=3,6,8,10,12,15,20)`$) from the remained parts. Each density fluctuation has a different fractal dimension given by
$$D_0=\frac{\mathrm{log}2}{\mathrm{log}[2n_D/(n_D1)]}.$$
(9)
We assume that $`\delta _+|\delta _{}|`$. With this ansatz, we have fixed the repetition number of removal procedure ($`N_R`$). Although the number $`N_R`$ can be different for each initial data, in order to keep the same resolution for each model, (i.e. for the ratio of the smallest line segment $`\mathrm{}`$ to that of the whole region $`L`$ (our calculation space) to be almost same), we set $`N_R=57`$. We also fix $`\delta _+=10^3`$, which determines the negative density perturbation $`\delta _{}`$ such that the mean of fluctuations must vanish (see Table 1).
| $`n_D`$ | $`D_0`$ | $`N_R`$ | $`\mathrm{}/L`$ | $`\delta _{}`$ |
| --- | --- | --- | --- | --- |
| 3 | 0.631 | 5 | 1/243 | $`0.152\times 10^3`$ |
| 6 | 0.792 | 6 | 1/191.10 | $`0.504\times 10^3`$ |
| 8 | 0.838 | 7 | 1/325.95 | $`0.647\times 10^3`$ |
| 10 | 0.868 | 7 | 1/267.62 | $`0.917\times 10^3`$ |
| 12 | 0.888 | 7 | 1/235.36 | $`1.19\times 10^3`$ |
| 15 | 0.909 | 7 | 1/207.47 | $`1.61\times 10^3`$ |
| 20 | 0.931 | 7 | 1/183.29 | $`2.31\times 10^3`$ |
Table 1 :
The number of division $`n_D`$, its fractal dimension $`D_0`$, the repetition number of removal procedure $`N_R`$, the resolution (the ratio of the length of the shortest line segment to that of the whole system $`\mathrm{}/L`$), and each negative density perturbation $`\delta _{}`$.
We show one example of initial data for $`n_D=10`$ in Fig. 1. We put $`2^{17}`$ sheets in our calculation space $`L`$. By use of a box-counting method (the box length $`x`$ is chosen from $`2^{16}L(1.5\times 10^5L)`$ to $`10^1L`$), we have checked the fractal dimension of initial fluctuations. We calculate the number $`N(x)`$ of boxes with a length $`x`$ in which $`\delta _+`$-segments exist. The result is shown in Fig. 2. We find a power law in the $`x`$-$`N(x)`$ relation for the range of $`10^3Lx10^1L`$. From this curve, we estimate the fractal dimension as 0.868, which is the same value as that of the present Cantor set. Below $`x=10^3L`$, we have a small deviation from a power law relation, which corresponds to the limit of the resolution of our present model.
Since the Cantor set is quite systematically constructed, one may wonder that our results strongly depend on such a special setting and may not be universal. To answer such a question, we also analyze two different initial data settings: one is a random Cantor-type set and the other is just a white noise.
#### 2.2.2 Random Cantor-type set
The random Cantor-type set is defined so that the division number $`n_D`$ is fixed as that of a regular Cantor set, while its removal positions of line segments are determined by a random number. According to the box-counting method, we find that the fractal dimension of initial fluctuations is almost the same as that of a regular Cantor set for the same $`n_D`$. In this paper, we analyze 3 random Cantor-type set models: two models with $`n_D=10`$ and one with $`n_D=12`$. One initial distribution for $`n_D=10`$ is shown in Fig. 3.
#### 2.2.3 White Noise
In order to see whether or not our primordial fractal fluctuations will play an important role in a structure formation, in particular a formation of fractal structure in a nonlinear regime, we also study a time evolution of primordial fluctuations with white noise. The distribution of initial density fluctuations is given by a random number between $`10^3\delta 10^3`$. We analyze 2 models, for one of which the initial data is shown in Fig. 4.
## 3 TIME EVOLUTION OF PRIMORDIAL FRACTAL DENSITY PERTURBATIONS
In order to see how structures are formed, we have to describe the distribution by the Eulerian coordinate. It is convenient to compare the distribution at each time by use of the comoving coordinate $`x`$. We set the initial scale factor $`a_0=1`$. Since our initial density fluctuation is $`10^3`$, we find a first shell crossing at $`a=a_{\mathrm{cross}}10^3`$. We perform our calculation until $`a=(23)\times 10^4`$.
In order to see the detail to resolve a fractal structure, we put $`2^{17}`$ sheets in our calculation scale $`L`$. Using a box-counting method, i.e, counting the number of boxes which contain the region with density perturbation $`\delta `$ larger than $`1`$, we determine the fractal dimension of the nonlinear structures. The size of the box ranges from $`2^{16}L`$ to $`10^1L`$
We shall present our results for three types of initial data in order.
### 3.1 Regular Cantor Set
First, we show the results for the regular Cantor set with $`n_D=10`$. The time evolution is depicted in Fig 5. Because we set $`\delta _+|\delta _{}|10^3`$, a nonlinear structure appears at $`a5\times 10^2`$ (Fig. 5(a)). Before a shell crossing, the pattern of density fluctuations remains similar to that of initial distributions, although each separation is going to change through the gravitational interaction. We then find a shell crossing at $`a=a_{\mathrm{cross}}10^3`$ (Fig. 5(b)). After the shell crossing, the trace of initial Cantor set gradually disappears, because of the exchange of the shells by a crossing (Fig. 5(b)-(f)). Although many sheets cross each other, the peculiar velocities do not vanish immediately because of collisionless sheets, and then the pattern of nonlinear structure will change continuously. After sufficient evolution of the nonlinear structure, we find a self-similarity in the structure, which seems to be fractal. In fact, enlarging some regions, we find similar density distributions (Fig. 6).
In order to judge whether such a structure is fractal or not, we use a box-counting method, which gives a fractal dimension $`D_F`$. Before a shell crossing, we find that the dimension $`D_F`$ decreases in time from the initial value $`D_0=0.868`$, but the error in the estimation increases with time (Fig. 7(a)). This is because although the pattern of initial density fluctuations remains even in a nonlinear stage before a shell crossing, change of each separation breaks the initial fractal distribution. Then, the initial fractal distribution seems to disappear.
However, after a shell crossing, the fractal dimension starts to increase again and the error in the estimation becomes much smaller. The fractal structure seems to recover. More surprisingly, the dimension $`D_F`$ approaches some constant ($`D_{\mathrm{asym}}0.9`$) after $`a1.5\times 10^4`$, which is a little bit different from the initial fractal dimension $`D_0=`$0.868 (Fig. 7.(b)). In fact, $`D_F=0.889\pm 0.009`$ at $`a=1.5\times 10^4`$ and $`D_F=0.890\pm 0.002`$ at $`a=2\times 10^4`$. Although $`D_0=0.868`$ is out of the error bar of $`D_F`$, the difference is very small. However we will see that $`D_F`$ is really independent of the initial fractal dimension $`D_0`$ later in the case of different initial distributions (see Fig. 10).
We also calculate the two-point spatial correlation function of nonlinear structures. If the structure is fractal, we have a relation between a fractal dimension and a power index of the correlation function (Falconer (1990)). The correlation function here is evaluated for the nonlinear regions where the density fluctuation $`\delta `$ becomes larger than $`1`$. At $`a=10^3`$, just after a first shell crossing, the correlation function shows a rapid oscillation between a positive and a negative values because of the periodic pattern of the Cantor set. After a shell crossing, the pattern of Cantor set disappears, and then such an oscillation also vanishes. After enough time passes, i.e. when the stable fractal structure is found, the correlation function becomes also stable (Fig. 8). The function is positive for the distance of $`x5\times 10^2L`$, beyond which it becomes negative because a shell crossing does not yet occur for such a scale and the largest nonlinear structure is about $`5\times 10^2L`$. The trace of initial Cantor set still remains. If the structure is fractal, the correlation function must show a power-law behavior. In fact, from Fig. 8, we find a power law relation in the range of $`10^4L<x<10^2L`$ as $`\xi x^\gamma `$ with $`\gamma =0.130\pm 0.005`$. When the fractal dimension is $`D_F`$, and the correlation function is given as
$$\xi (x)x^{(1D_F)},$$
(10)
then we find $`\gamma =1D_F`$. The evaluated fractal dimension from two point correlation is $`D_F=0.870\pm 0.005`$ at $`a=10000`$, which is close to the value obtained by a box-counting $`D_F=0.862\pm 0.009`$. We show the time evolution of two ”fractal” dimensions $`D_F`$ obtained from two point correlation and by a box-counting method in Fig. 9. From Fig. 9, we confirm that the power index $`\gamma `$ is well correlated to the fractal dimension $`D_F`$.
We find a fractal structure after a shell crossing. The fractal dimension is close to the fractal dimension of the initial density distributions. Does this dimension reflect that of the initial distributions? If this is so, why does it disappear once around a first shell crossing time and recover at very late time? In order to answer these question, we have looked for other initial conditions with different $`n_D`$s, i.e. $`n_D=3,6,8,12,15,20`$. The fractal dimensions $`D_0`$ of the initial density distributions are given in Table 1. As for the time evolution, we find similar behaviour for all models. The evolution of ”fractal” dimension $`D_F`$ is shown in Fig. 10. Surprisingly, for all models, all $`D_F`$ approach about $`0.9`$ at $`a=(23)\times 10^4`$, which is the end of our calculation. To confirm our result, we have checked the size-number ($`x`$-$`N(x)`$) relation by a box-counting method, which shows almost a straight line as in Fig. 2. Since the initial dimensions of primordial fluctuations were different, we would conclude that the fractal dimension obtained after nonlinear evolution is universal within our numerical accuracy.
### 3.2 Random Cantor-type Set
Since the Cantor set is highly systematically constructed, one may ask whether the present result strongly depends on such a special initial setting. Is the universal dimension of nonlinear fractal structures due to the primordial density fluctuations defined by the regular Cantor set? In order to answer this question, we shall analyze a different model with randomness, which we call a random Cantor-type set defined in the previous subsection.
We analyze 3 models: 2 models with $`n_D=10`$ (model 1 and model 2) and 1 model with $`n_D=12`$. We find that just as in the case of regular Cantor sets, the fractal dimension for nonlinear structures always approaches about $`0.9`$ (Fig. 11). We have also checked the size-number ($`x`$-$`N(x)`$) relation in a box-counting method, finding the same result as in the case of a regular Cantor set. Because we remove a line segment at a random position, we usually expect that the smallest segment is smaller than that of the regular Cantor set as we show in the previous subsection. As a result, the stable fractal structure will be formed later compared with the case with a regular Cantor set. In fact, in the case of $`n_D=12`$, we find the stable dimension ($`0.9`$) around $`a2.3\times 10^4`$ (Fig.11(c)).
### 3.3 White noise case
Then one may have another question. Is the universal fractal dimension obtained above via nonlinear dynamics and independent of the initial distribution? In order to answer this question, we also analyze a model with white noise fluctuation. We analyze two models. Both models do not show up the above universal fractal dimension, although we find some different stable asymptotic dimension ($`0.7`$) (Fig. 12). The error in estimating the dimension is larger than that of the Cantor set model, and the box-counting shows some deviation from the power law relation (Fig. 13). Hence, it would not be a fractal. We will discuss this in out Discussion.
## 4 ANALYSIS OF PHASE SPACE
Although the analysis by a box-counting suggests that a nonlinear fractal structure with a universal dimension appears from primordial fractal fluctuations, we may get more information from a detailed study of the nonlinear structures obtained. For this purpose, we analyze our result in a phase space.
We show the time evolution of the structures for the regular Cantor set with $`n_D=10`$ in a phase space. Initially, the sheets distribution in a phase space is given by a notched curve because $`\delta _+`$ and $`\delta _{}`$ are constant (Fig. 14). If we enlarge the pictures, we find a similar notched curve because of the present initial setting. These notches reflect a self-similarity in the initial Cantor set. This behavior does not change before a shell crossing (Figs. 14 (a), (b)). Only the slopes of the line segments become steeper because of the concentration of sheets. After a shell crossing, the curve in a phase space shows very complicated behavior. Two sheets exchanged by a shell crossing are decelerated by the mutual gravitational interaction, and then the curve will swirl (Gouda & Nakamura 1989). As the structure evolves, some vortices are combined and form a larger vortex (Fig. 14). When the ”fractal” dimension becomes stable around 0.9, we find that the large vortex consists of some similar small vortices. These small vortices also consist of similar but much smaller vortices (Fig. 15). This discrete self-similarity in a phase space is found in all models. As the structure evolves, some vortices are combined and form a larger vortex (Fig. 14). Although the initial fractal distribution seems to disappear, some trace remains in the phase space.
One may wonder why the centers of the vortices appear in $`\stackrel{~}{v}>0`$ region ( or $`\stackrel{~}{v}<0`$ region). The formation of a vortex in a phase space can be easily understood in the case of a single wave ($`S_1=ϵ\mathrm{sin}q`$) (see Gouda & Nakamura (1989)). In that case, the velocity at the center of the vortex vanishes. Then the appearance of the vortex in $`\stackrel{~}{v}>0`$ region ( or $`\stackrel{~}{v}>0`$ region) seems inconsistent. If we pursue each particle motion, nothing strange happens. They move without swirling except at $`\stackrel{~}{v}=0`$ point as shown in Fig. 16.
As was the case with a regular Cantor-type set, we also find similar results in the phase space (Fig. 17). However, a discrete self-similarity is hard to find, although the larger vortices contain smaller vortices as in the case of the regular Cantor-type set. In the case with white noise fluctuation, we cannot find any hierarchical vortex structure in a phase space (Fig. 18).
## 5 CONCLUSIONS AND DISCUSSION
We have studied the nonlinear evolution of primordial fractal fluctuations by using a 1-D sheet model. We have analyzed 7 models with initial fluctuations constructed by a regular Cantor set, 3 models with initial fluctuations constructed by a random Cantor-type set, and 2 models with white noise fluctuations. For all models except for the case with white noise, we find a kind of attractor with a universal fractal dimension ($`0.9`$) as the fluctuations evolve into nonlinear regime. In the case with white noise fluctuations, the estimated dimension becomes stable around $`0.7`$, but the error in the estimation is larger than the other cases and the power-law behavior in a box-counting is also not completely fitted. Then, it may not contain a fractal structure. From the phase space analysis, we find a hierarchical structure, that is, the large vortex consists of some similar small vortices, and such small vortices again consist of similar but much smaller vortices. In particular, we find a discrete self-similarity for the model with a regular Cantor set.
Why is the fractal dimension close to $`0.9`$? Is it really universal? Is the present fractal structure really an attractor? Although we need more analysis to answer this question, we have some hints in previous work. Gouda and Nakamura studied the present 1-D sheet model for the initial power law spectrum. They found two types of generic singularities when we have a shell crossing (Gouda & Nakamura 1988, 1989). When a first shell crossing appears, the relation between Eulerian and Lagrangian coordinates must be
$$x=q_c+\beta (qq_c)^3+\mathrm{},$$
(11)
while that after a shell crossing turns out to be
$$x=q_c+\beta (qq_c)^2+\mathrm{}.$$
(12)
Following Arnold’s classification, the former and latter cases are classified into A3 and A2, respectively. A3 is structurally unstable and may appear transiently in the expanding Universe. A2 is structurally stable and appears universally for the initial power-law spectrum. The latter case gives
$`\delta _k`$ $`=`$ $`{\displaystyle \delta (x)e^{ikx}𝑑x}`$ (13)
$``$ $`{\displaystyle (\beta x)^{1/2}e^{ikx}𝑑x}`$
$`=`$ $`(\beta k)^{1/2}{\displaystyle \eta ^{1/2}e^{i\eta }𝑑\eta },`$
i.e. $`P(k)k^1`$. This predicts $`\gamma =0`$, i.e. $`D_F=1`$, which is rather close to our “universal” dimension 0.9. Although one may wonder that these are essentially the same, we have another result which suggests that there seems to exist a new type of stable phase. Recently, Yano and Gouda analyzed a more realistic case, i.e. the initial power law spectrum with a cut-off and found 5 characteristic regions in Fourier space (Yano & Gouda 1998a, b). The regime 1 is the linear one and then it is just an initial power spectrum. In the regime 2, they found $`P(k)k^1`$, which is the single-caustic regime (Gouda & Nakamura (1989)). The regime 3 is called the multi-caustic regime, in which the power spectrum depends on the initial power-law index. Beyond the cut-off scale, two regimes appear, one gives $`P(k)k^1`$ (regime 5), which may correspond to A2 type stable solution. In the intermediate wave number $`k`$ between the regime 3 and regime 5, they find $`k^\nu `$, which $`\nu `$ is independent of initial power index and close to 1, but a little less. They called it the virialized regime. This seems to be a new transient region, which may appear in some specific initial conditions. We would conjecture that the fractal structure with a universal dimension 0.9 corresponds to this virialized regime (regime 4) and the dimension 0.7 found in the case with white noise would be the regime 3. By reanalyzing the Yano-Gauda model in the case of $`k=0`$, we have confirmed that $`\nu =0.9`$. We also find a small tail with index 0.7 in the size-number relation in the Cantor set model with $`n_D=15,20`$ (Fig. 19).
This conjecture is also supported by the analysis for a self-gravitating 1-D sheet model without the background expansion of the Universe (Tsuchiya, Konishi & Gouda 1994). They found two time scales; one is a micro relaxation time ($`t_{\mathrm{micro}}=Nt_c`$) and the other is a global relaxation time ($`t_{\mathrm{global}}=4\times 10^4Nt_c`$), where $`t_c=\sqrt{L/4\pi GNm}`$ is a crossing time. After $`t_{\mathrm{micro}}`$, some equilibrium state is reached by exchanging particle energy, but the global relaxation is not achieved, i.e. the partition function is not yet described by an equilibrium state such as an ergodic state (Tsuchiya, Konishi & Gouda 1994). In the present model, we can speculate that the fractal structure is obtained after this micro relaxation time but before the global relaxation time. In fact, if we estimate the above time scales in the present models, we find that $`t_{\mathrm{micro}}`$ corresponds to $`a=5\times 10^3`$, while $`t_{\mathrm{global}}`$ corresponds to $`a=5\times 10^6`$. The time when we find a stable fractal structure ($`a=(13)\times 10^4`$) is between those two time scales. If this speculation is true, our fractal structure is temporal. In the future of the Universe, it will evolve into more relaxed and ergodic state.
Since we analyze the simplest case, we have to extend our analysis to more generic cases. First, we should study different types of fractal in order to check whether the present results are universal for any fractal distributions or not. Secondly, we need to analyze the case with scale-dependent fluctuations. In the present analysis, we set $`\delta _+=`$ constant and $`\delta _{}=`$ constant. In the realistic case, there must be a scale dependence to the fluctuations. In the conventional perturbations, we usually assume a power-law spectrum with some cutoff. Even if the primordial fluctuations contain a fractal structure, their amplitude may depend on the scale. Its dependence may change the present results. In particular, in the present model, the scenario of structure formation could be different from either TOP-DOWN or BOTTOM-UP for some scale-dependence. The primordial fractal fluctuations will evolve directly into a hierarchical nonlinear structure. But, it will definitely depend on the scale dependence of the fluctuations. Secondly, we need to extend the present analysis to other cosmological models, i.e. the open Universe model and $`\mathrm{\Lambda }0`$ flat Universe model. For the 1-D sheet model, the solutions are still exact, and the growth and decay rates in these models are different from those in the Einstein-de Sitter Universe model. We expect that the structure formation after a shell crossing is not the same as that in the present cosmological model, and then the fractal dimension would be different.
For more realistic cases, we must study either the 2-D or 3-D model. Since the Zel’dovich solution is no longer exact, we have to explore a new method. In order to preserve a high resolution, we may develop a kind of renormalization method in $`N`$-body simulation as Couchman & Peebles (1998).
Finally, it would also be interesting to look for the origin of such a primordial fractal density perturbation. The inflationary scenario may provide the origin of primordial fluctuations. One may wonder whether such a fractal primordial fluctuation is expected in some inflationary models. If we have more than two scalar fields, then the system is not integrable and may show a chaotic behavior or a fractal property (Easther & Maeda (1999)). Such a model might show up a kind of fractal density perturbation.
We would like to thank N. Gouda, P. Haines, O. Iguchi, T. Kurokawa, V. Lukash, M. Morikawa, A. Nakamichi, Y. Sota, T. Yano, and A. Yoshisato for many useful discussions. Part of this work was done while KM was participating the program, “Structure Formation in the Universe”, at the Newton Institute, University of Cambridge. KM is grateful to the Newton Institute for their hospitality. Our numerical computation was carried out by Yukawa Institute Computer Faculty. This work was supported partially by a Grant-in-Aid for Scientific Research Fund of the Ministry of Education, Science and Culture (Specially Promoted Research No. 08102010), and by the Waseda University Grant for Special Research Projects.
Figure Captions:
Fig.1. The initial fractal density fluctuation for the regular Cantor set with $`n_D=10`$. Enlarging the picture, we find the same pattern up to the repetition number of removal procedure $`N_R=7`$.
Fig. 2. The size-number ($`x`$-$`N(x)`$) relation in a box-counting method for the model in Fig. 1. The dotted line shows $`N(x)x^{0.868}`$. We find that the similarity exists in the range between $`10^3L`$ and $`10^1L`$. Below $`x=10^3L`$, we have a deviation from the power law relation, which corresponds to our resolution limit ($`\mathrm{}/L1/267.62`$).
Fig. 3.(a) The initial fractal density fluctuation for the random Cantor-type set with $`n_D=10`$. The random Cantor-type set is defined so that the division number $`n_D`$ is fixed as that of a regular Cantor set, while its removal positions of line segments are determined by a random number. This is called the model 1 in the text. Enlarging the picture, we find the similar pattern up to the repetition number of removal procedure $`N_R=7`$. The fractal dimension of this initial fluctuation is $`0.868`$, which is the same as that of the regular Cantor set with $`n_D=10`$.
(b) Another initial fractal density fluctuation set up by the same prescription as (a) (model 2).
Fig. 4.(a) The initial density fluctuation created by a white noise. This is called the model 1 in the text.
(b) Another initial density fluctuation set up by the same prescription as (a) (model 2).
Fig. 5.(a) Time evolution of density fluctuation for the case with $`n_D=10`$. This snap shot is at $`a=5\times 10^2`$, when the positive fluctuation $`\delta _+`$ grows just to a nonlinear scale ($`\delta _+=1`$).
(b) The same as (a) but at $`aa_{\mathrm{cross}}10^3`$, which is just after a first shell crossing.
(c) The same as (a) but at $`a=6\times 10^3`$.
(d) The same as (a) but at $`a=10^4`$.
(e) The same as (a) but at $`a=1.5\times 10^4`$.
(f) The same as (a) but at $`a=2\times 10^4`$.
Fig. 6. Hierarchical nonlinear structure at $`a=10^4`$. Enlarging the picture, we find the self-similar structure.
Fig. 7.(a) The evolution of the ”fractal” dimension $`D_F`$ of nonlinear structure for the case of $`n_D=10`$. Before a shell crossing ($`a<a_{\mathrm{cross}}10^3`$).
(b) The same as (a) but after a shell crossing ($`aa_{\mathrm{cross}}10^3`$).
Fig. 8. Two-point correlation function of nonlinear structures at $`a=10^4`$. The dotted line is $`\xi x^{0.130}`$.
Fig. 9. The time evolution of the ”fractal” dimensions $`D_F`$ obtained from two point correlation and by a box-counting method The circles denote $`D_F`$ by two-point correlation ($`D_F=1\gamma `$), while the squares are those produced by a box-counting. Two methods to determine the fractal dimension agree well each other.
Fig. 10.(a) The time evolution of the ”fractal” dimension $`D_F`$ of nonlinear structure for $`n_D=3`$.
(b) The same as (a) but $`n_D=6`$.
(c) The same as (a) but $`n_D=8`$.
(d) The same as (a) but $`n_D=12`$.
(e) The same as (a) but $`n_D=15`$.
(f) The same as (a) but $`n_D=20`$.
Fig. 11.(a) The evolution of the ”fractal” dimension $`D_F`$ of nonlinear structure in the case of the random Cantor-type initial fluctuations for $`n_D=10`$ (model 1). The dimension approaches about $`0.9`$ at $`a=2\times 10^4`$.
(b) The same as (a) ($`n_D=10`$, model 2). The dimension approaches about $`0.9`$ at $`a=2.8\times 10^4`$.
(c) The same as (a) ($`n_D=12`$). The dimension approaches about $`0.9`$ after $`a=2.3\times 10^4`$.
Fig. 12.(a) The evolution of the ”fractal” dimension $`D_F`$ of nonlinear structure in the case with white noise fluctuations (model 1). The dimension approaches about $`0.7`$.
(b) The same as (a) but of model 2. We also find the same stable dimension $`0.7`$.
Fig. 13. The size-number ($`x`$-$`N(x)`$) relations in a box-counting method for the models in Fig. 12. There is some small deviation from a power-law behaviour.
Fig.14.(a) Structure Formation in phase space ($`n_D=10`$). At $`a=1`$. The distribution draws a self-similar graph, because initial density fluctuation was given by a regular Cantor set.
(b) The same as (a) but at $`a=10^3`$. At immediately after shell crossing, the sheets just began production of vortices.
(c) The same as (a) but at $`a=3\times 10^3`$. The self similarity of vortices exists during four steps.
(d) The same as (a) but at $`a=5\times 10^3`$. The self similarity of vortices exists during five steps.
(e) The same as (a) but at $`a=10^4`$. Two of $`L`$ scale structure was begun to combined.
(f) The same as (a) but at $`a=1.5\times 10^4`$.
(g) The same as (a) but at $`a=2\times 10^4`$.
Fig. 15. A sheet distribution at $`a=2\times 10^4`$ in phase space ($`n_D=10`$). The dotted squared regions are enlarged in the next figures shown by arrows. We find a discrete self similarity, i.e. a large vortex consists of some similar small vortices, and those small vortices also consist of similar but much smaller vortices.
Fig. 16. The orbits of 7 particles in phase space ($`n_D=10`$, 6000 $`a20000`$). Each particle did not swirl, however distribution of particles become to set of vortices.
Fig. 17. A sheet distribution at $`a=30000`$ in phase space for the model with a random Cantor-type set ($`n_D=12`$). In this figure, the larger vortices contain smaller vortices as the case of regular Cantor set, but a discrete self-similarity is hard to be found.
Fig. 18. A sheet distribution $`a=30000`$ in phase space for the model with white noise fluctuations (model 1). The vortices were widely spread and a nest of vortices is rarely found.
Fig. 19. The size-number ($`x`$-$`N(x)`$) relation in a box-counting method for the regular Cantor set with $`n_D=15`$ at $`a=30000`$. The dotted line denotes $`N(x)x^{0.9}`$, while the dotted-dash line is $`N(x)x^{0.7}`$.
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# THE STRUCTURE OF THE NUCLEONaafootnote aSupported by the Deutsche Forschungsgemeinschaft (SFB 443)
## 1 Introduction
Nucleons are composite systems with many internal degrees of freedom. The constituents are quarks and gluons, which are bound by increasingly strong forces if the momentum transfer decreases towards the GeV region. The “running” coupling constant of the strong interaction, $`\alpha _s(Q^2)`$ in fact diverges if $`Q^2`$ approaches $`\mathrm{\Lambda }_{QCD}^2(200`$MeV/c)<sup>2</sup> corresponding to a scale in space of about 1 fm. This is the realm of nonperturbative quantum chromodynamics (QCD), where confinement plays a major role, and quarks and gluons cluster in color neutral objects. Such correlations between the constituents have the consequence that nucleons in their natural habitat, i.e. at the confinement scale, have to be described by hadronic degrees of freedom rather than quarks and gluons.
QCD is a nonlinear gauge theory developed on the basis of massless quarks and gluons $`^\mathrm{?}`$. The interaction among the gluons gives rise to the nonlinearity, and the interaction among the quarks is mediated by the exchange of gluons whose chromodynamic vector potential couples to the vector current of the quarks. If massless particles interact by their vector current, their helicity remains unchanged. In practice one has to restrict this discussion to $`u`$, $`d`$ and $`s`$ quarks with masses $`m_u5`$ MeV, $`m_d9`$ MeV and $`m_s175`$ MeV, which are all small at the mass scale of the nucleon. These quarks can be described by SU(3)$`{}_{R}{}^{}`$ SU(3)<sub>L</sub> as long as right and left handed particles do not interact, which is what happens if the helicity is conserved. By combining right and left handed currents, one obtains the vector currents $`J_\mu ^a`$ and the axial vector currents $`J_{5\mu }^a`$,
$$J_\mu ^a=\overline{q}\gamma _\mu \frac{\lambda ^a}{2}q,J_{5\mu }^a=\overline{q}\gamma _\mu \gamma _5\frac{\lambda ^a}{2}q,$$
(1)
where $`q`$ are Dirac spinors of the massless and point-like light quarks and $`\gamma _\mu ,\gamma _5`$ the appropriate Dirac matrices. The quantities $`\lambda ^a`$, $`a=1\mathrm{}8`$ denote the Gell-Mann matrices of SU(3) describing the flavor structure of the 3 light quarks. It is often convenient to introduce the unit matrix $`\lambda ^0`$ in addition to these matrices.
In the context of these lectures we shall only need the “neutral” currents corresponding to $`\lambda =3`$, 8 and 0, which have a diagonal form in the standard representation. The photon couples to quarks by the electromagnetic vector current $`J_\mu ^{em}J_\mu ^{(3)}+\frac{1}{\sqrt{3}}J_\mu ^{(8)}`$, corresponding to isovector and isoscalar interactions respectively. The weak neutral current mediated by the $`Z^0`$ boson couples to the 3rd, 8th and 0th components of both vector and axial currents. While the electromagnetic current is always conserved, $`^\mu J_\mu ^{em}=0`$, the axial current is only conserved in the limit of massless quarks. In this limit there exist conserved charges $`Q^a`$ and axial charges $`Q_5^a`$, which are connected by current algebra,
$$[Q^a,Q^b]=if^{abc}Q^c,[Q_5^a,Q_5^b]=if^{abc}Q^c,[Q_5^a,Q^b]=if^{abc}Q_5^c,$$
(2)
with $`f^{abc}`$ the structure constants of SU(3). Such relations were an important basis of low energy theorems (LET), which govern the low energy behavior of (nearly) massless particles.
The puzzle we encounter is the following: The massless quarks appearing in the QCD Lagrangian conserve the axial currents but the nucleons as their physical realizations are massive and therefore do not conserve the axial currents. The puzzle was solved by Goldstone’s theorem. At the same time as the “3 quark system” nucleon becomes massive by means of the QCD interaction, the vacuum develops a nontrivial structure due to finite expectation values of quark-antiquark pairs (condensates $`\overline{q}q`$), and so-called Goldstone bosons are created, $`\overline{q}q`$ pairs with the quantum numbers of pseudoscalar mesons. These Goldstone bosons are massless, and together with the massive nucleons they act such that chirality is locally conserved. This mechanism can be compared to the local gauge symmetry of quantum electrodynamics, which is based on the fact that both (massless) photon and (massive) matter fields have to be gauge transformed.
In QCD the chiral symmetry is definitely broken by the small but finite quark masses. As a consequence also the physical “Goldstone bosons”, in particular the pions, acquire a finite mass $`m_\pi `$, which is generally assumed (though not proven) to follow the Gell-Mann-Oakes-Renner relation
$$m_\pi ^2f_\pi =(m_u+m_d)\overline{q}q,$$
(3)
with the condensate $`\overline{q}q(225\text{MeV})^3`$, and $`f_\pi 93`$MeV the pion decay constant. Since the pions are now massive, the corresponding axial currents are no longer conserved and the 4-divergence of the axial current becomes
$$^\mu J_{5\mu }^af_\pi m_\pi ^2\varphi _\pi ^a,$$
(4)
where $`\varphi _\pi ^a`$ describes the local field of charged pions $`(a=1`$ and 2). In other words the weak decays
$$\pi ^+\mu ^++\nu _\mu \text{and}\pi ^{}\mu ^{}+\overline{\nu }_\mu $$
(5)
proceed via coupling to the axial current (Fig. 1).The pion and its axial current disppear from the hadronic world and leave the (hadronic) vacuum behind. In particular we note that a finite value of the divergence of Eq. (4) has 3 requirements: the decay of the pion can take place, the pion mass is finite, and a local pion field exists.
While the charged pions decay weakly with a life-time of $`2.610^8`$ sec, the neutral pion decays much faster, in $`8.410^{17}`$ sec, by means of the electromagnetic interaction,
$$\pi ^0\gamma +\gamma .$$
(6)
Again axial current disappears, corresponding to
$$^\mu J_{5\mu }^3=\frac{\alpha _{fs}}{\pi }\stackrel{}{E}\stackrel{}{B},$$
(7)
where $`\alpha _{fs}=e^2/4\pi `$ is the fine structure constant, and $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ are the electromagnetic fields. We note that two electromagnetic fields have to participate, because two photons are created, and that they have to be combined as a pseudoscalar, because the pseudoscalar pion disappears. The transition of Eq. (7) can be visualized by the intermediate quark triangle of Fig. 1. It is called the “triangle anomaly”, because such transitions cannot exist in classical theories but only occur in quantum field theories via the renormalization procedure. Such terms are also predicted on general grounds (Wess-Zumino-Witten term). We note in passing that a similar anomaly is obtained in QCD by replacing the electromagnetic fields by the corresponding color fields, $`\stackrel{}{E}_c`$ and $`\stackrel{}{B}_c`$, $`\alpha _{fs}`$ by the strong coupling $`\alpha _s`$, and by an additional factor 3 for u, d, and s quarks,
$$^\mu J_{5\mu }^0=3\frac{\alpha _s}{\pi }\stackrel{}{E}_c\stackrel{}{B}_c.$$
(8)
As a consequence the component $`J_{5\mu }^0`$ is not conserved, not even in the case of massless quarks (“$`U_A(1)`$ anomaly”).
Unfortunately, no ab-initio calculation can yet describe the interesting but complicated world of the confinement region. In principle, lattice gauge theory should have the potential to describe QCD directly from the underlying Lagrangian. However, these calculations have yet to be restricted to the “quenched approximation”, i.e. initial configurations of 3 valence quarks. This is a bad approximation for light quarks, because the Goldstone mechanism creates plenty of sea quarks, and therefore the calculations are typically performed for massive quarks, $`m_q100`$ MeV, and then extrapolated to the small $`u`$ and $`d`$ quark masses. In this way one obtains reasonable values for mass ratios of hadrons and qualitative predictions for electromagnetic properties. However, some doubt may be in order whether such procedure will describe the typical threshold behavior of pionic reactions originating from the Goldstone mechanism.
A further “ab initio” calculation is chiral perturbation theory (ChPT), which has been established by Weinberg in the framework of effective Lagrangians and put into a systematic perturbation theory by Gasser and Leut-wyler $`^\mathrm{?}`$. Based on the Goldstone mechanism, the threshold interaction of pions is weak both among the pions and with nucleons, and furthermore the pion mass is small and related to the small quark masses $`m_u`$ and $`m_d`$ by Eq. (3). As a consequence ChPT is a perturbation in a parameter $`p:=(p_1,p_2,\mathrm{};m_u,m_d)`$, where $`p_i`$ are the external 4-momenta of a particular (Feynman) diagram. (Note that also the time-like component, the energy, is small at threshold because of the small mass!). This theory has been applied to photoinduced reactions by Bernard, Kaiser, Meißner and others $`^\mathrm{?}`$ over the past decade. As a result several puzzles have been solved and considerable insight has been gained. There exists, however, the problem that ChPT cannot be renormalized in the “classical” way by adjusting a few parameters to the observables. Instead the renormalization has to be performed order by order, the appearing infinities being removed by counter terms. This procedure gives rise to a growing number of low energy constants (LECs) describing the strength of all possible effective Lagrangians consistent with QCD, at any given order of the perturbation series. These LECs, however, cannot (yet) be derived from QCD but have to be fitted to the data, which leads to a considerable loss of predictive powers in the higher orders of perturbation. A further problem arises in the nucleonic sector due to the nucleon’s mass $`M`$, which is of course not a small expansion parameter. The latter problem has been overcome by heavy baryon ChPT (HBChPT), a kind of Foldy-Wouthuysen expansion in $`M^1`$. The solution is achieved, however, at the expense of going from an explicitly relativistic field theory to a nonrelativistic scheme.
Beside lattice gauge theory and ChPT, which are in principle directly based on QCD, there exists a host of QCD inspired models, which we shall not discuss at this point but occasionally refer to at later stages.
## 2 KINEMATICS
Let us consider the kinematics of the reaction
$$e(k_1)+N(p_1)e(k_2)+N(p_2),$$
(9)
with $`k_1=(\omega _1,\stackrel{}{k}_1)`$ and $`p_1=(E_1,\stackrel{}{p}_1)`$ denoting the four-momenta of an electron $`e`$ and a nucleon $`N`$ in the initial state, and corresponding definitions for the final state (Fig. 2). These momenta fulfil the on-shell conditions
$$p_1^2=p_2^2=M^2,k_1^2=k_2^2=m^2,$$
(10)
and furthermore conserve total energy and momentum,
$$k_1+p_1=k_2+p_2.$$
(11)
If we also assume parity conservation, the scattering amplitudes should be Lorentz invariants depending on the Lorentz scalars that can be constructed from the four-momenta. By use of Eqs. (10) and (11) it can be shown that there exist only two independent Lorentz scalars, corresponding to the fact that the kinematics of Eq. (9) is completely described by, e.g., the lab energy of the incident electron, $`\omega _L`$, and the scattering angle $`\mathrm{\Theta }_L`$. In order to embed relativity explicitly, it is useful to express the amplitudes in terms of the 3 Mandelstam variables
$$s=(k_1+p_1)^2,t=(k_2k_1)^2,u=(p_2k_1)^2.$$
(12)
Since only two independent Lorentz scalars exist, these variables have to fulfil an auxiliary condition, which is
$$s+t+u=2(m^2+M^2).$$
(13)
In the $`cm`$ frame, the 3-momenta of the particles cancel and $`s=(\omega _{cm}+E_{cm})^2=W_s^2=W^2`$, where $`W`$ is the total energy in that frame. Furthermore, the initial and final energies of each particle are equal, hence $`t=(\stackrel{}{k}_2\stackrel{}{k}_1)_{cm}^2=\stackrel{}{q}_{cm}^2`$, where $`\stackrel{}{q}_{cm}`$ is the 3-momentum transfer in the $`cm`$ system. From these definitions it follows that $`s(m+M)^2`$ and $`t0`$ in the physical region. Since $`s`$ is Lorentz invariant, the threshold energy $`\omega _{lab}`$ can be obtained by comparing $`s`$ as expressed in the $`lab`$ and $`cm`$ frames. Moreover, in a general frame $`t=(k_2k_1)^2=q^2<0`$ describes the square of 4-momentum of the virtual photon $`\gamma ^{}`$, exchanged in the scattering process (“space-like photon”). Since $`t`$ is negative in the physical region of electron scattering, we shall define the positive number $`Q^2=q^2`$ for further use. We also note that in pair annihilation, $`e^+e^{}\gamma ^{}`$, the square of 4-momentum is positive, $`q^2=m_\gamma ^{}^2>0`$ (“time-like photon”).
The above equations can be easily applied to Compton scattering,
$$\gamma (k_1)+N(p_1)\gamma (k_2)+N(p_2),$$
(14)
by replacing $`m`$ by zero, the mass of a real photon, and to virtual Compton scattering (VCS),
$$\gamma (k_1)+N(p_1)\gamma (k_2)+N(p_2),$$
(15)
by replacing $`m^2k_1^2=q^2<0`$. Due to the spins of photon and nucleon, several Lorentz structures appear in the scattering amplitude, and each of these structures has to be multiplied by a scalar function depending in the most general case on 3 variables, $`F=F(s,t,Q^2)`$.
Another generalization occurs if the nucleon is excited in the scattering process, in which case $`p_2^2=(M^{})^2>M^2`$ becomes an additional variable. Introducing the Bjorken variable $`x=Q^2/2p_1q`$ we find that $`x=1`$ corresponds to elastic scattering, while inelastic scattering is described by values $`0x<1`$.
For further use we shall acquaint ourselves with the Mandelstam plane for (real) Compton scattering, as shown by Fig. 3. Due to the symmetry of the Mandelstam variables, the figure can be constructed on the basis of a triangle with equal sides and heights equal to $`2M^2`$ according to Eq. (13) for $`m=m_\gamma =0`$. The axes $`s=0`$, $`t=0`$, and $`u=0`$ are then obtained by drawing straight lines through the sides of the triangle. The physically allowed region for $`k_1+p_1k_2+p_2`$ is given by the horizontally hatched area called “s channel” with $`sM^2`$ and $`t0`$. If we replace $`p_1p_1`$ and $`p_2p_2`$ in Eq. (11), we obtain the “u channel” reaction $`k_1+p_2k_2+p_1`$ given by the horizontally hatched area to the left.
Finally, if we look at Fig. 2 from the left side, we obtain the $`t`$ channel reaction $`\gamma (k_1)+\gamma (k_2)N(p_1)+\overline{N}(p_2)`$, which corresponds to the replacements $`k_2k_2`$ and $`p_1p_1`$ and is physically observable for $`t>2M^2`$ (hatched area at top of Fig. 3). Referring again to the $`s`$ channel, the boundaries of the physical region correspond to the scattering angles 0 and 180. The former case leads to zero momentum transfer, i.e. the line $`t=0`$, the latter case to the hyperbolic boundary of the region at negative $`t`$ values. The u-channel region is then simply obtained by a reflection of the figure at the line $`s=u`$ given by the $`t`$ axis. Finally, the boundary of the t-channel region is given by the upper branch of the hyperbola, separated from the lower one by $`4M^2`$.
Still in the context of Compton scattering, Fig. 4 shows the Born diagrams (tree graphs) contributing to the reaction. In order of appearance on the $`rhs`$, we find the direct, the crossed, and the $`\pi ^0`$ pole terms, exhibiting pole structures as $`(sM^2)^1`$, $`(uM^2)^1`$ and $`(tm_\pi ^2)^1`$ respectively. Except for the origin at $`s=u=M^2`$ (“scattering” of photons with zero momentum), these poles are situated on straight lines outside of the physical regions. However, photon scattering at small energies is obviously dominated by the poles at $`s=M^2`$ and $`u=M^2`$. The “low-energy theorem” asserts that for a particle with charge $`e`$ and mass $`M`$, the scattering amplitude behaves as $`T=\frac{e^2}{M}+O(\omega _{cm}^2)`$, where $`e^2/4\pi 1/137`$. It is derived on the basis that (i) only the Born terms have pole singularities for $`\omega _{cm}0`$, which results in the Thomson amplitude $`(e^2/M)`$, and (II) gauge invariance or current conservation, which allows one to express the next-to-leading-order terms in $`\omega _{cm}`$ by the Born contributions. Therefore, the internal structure (polarizability) of the system enters only in terms of relative order $`\omega _{cm}^2`$, i.e. is largely suppressed near threshold.
If the energy of the photon is sufficient to produce a pion, $`\sqrt{s}>M+m_\pi `$, Compton scattering competes with the much stronger hadronic reactions and becomes complex. The same is true in the $`t`$ channel, whenever the two photons carry more energy than $`\sqrt{t}=2m_\pi `$. Therefore the Compton amplitudes are only real in an area around the origin $`(s=u=M^2,t=0)`$, i.e. in the triangle shaped by the dashed lines in Fig. 3. Due to this reality relation, however, the Compton amplitudes can be analytically continued into the unphysical region, and information from the different physical regions can be combined to construct a common amplitude for the whole Mandelstam plane. Summarizing the role of the singularities for the specific reaction of Compton scattering we find: (I) The nucleon poles in the direct and crossed Born graphs, at $`s=M^2`$ and $`u=M^2`$, which are close to and therefore important near threshold, (II) the pion pole term at $`t=m_\pi ^2`$ and a branch cut starting at $`t=4m_\pi ^2`$ due to the opening of the $`2\pi `$ continuum, which affect the forward amplitude at any energy and (III) the opening of hadronic channels at $`s,u>(M+m_\pi )^2`$, which lead to a complex amplitude and a much enhanced Compton cross section, particularly near resonances at $`s=M_{\text{res}}^2`$.
Let us finally consider the spin degrees of freedom of the involved particles. A virtual photon with momentum $`\stackrel{}{q}`$ carries a polarization described by the vector potential $`\stackrel{}{A}`$, which has both a transverse part, $`\stackrel{}{A}_T\stackrel{}{q}`$, as in the case of a real photon, and a longitudinal component $`\widehat{q}\stackrel{}{A}`$, which is related to the time-like component $`A_0`$ by current conservation, $`qA=q_0A_0\stackrel{}{q}\stackrel{}{A}=0`$. As a consequence the cross section for the reaction of Eq. (9) takes the (somewhat symbolical) form
$$\frac{d\sigma }{d\mathrm{\Omega }}=\mathrm{\Gamma }(\sigma _T+\epsilon \sigma _L),$$
(16)
where $`\mathrm{\Gamma }`$ describes the flux of the virtual photon spectrum, and $`\sigma _T`$ and $`\sigma _L`$ the transverse and longitudinal cross sections respectively. The so-called transverse polarization $`\epsilon `$ of the virtual photon field is given by kinematical quantities only, which can be varied such that the partial cross sections remain constant. In this way the two partial cross sections can be separated by means of a “Rosenbluth plot”.
Concerning the electron, we shall assume that it is highly relativistic, hence its spin degree of freedom will be described by the helicity $`h=\stackrel{}{s}\widehat{k}=\pm \frac{1}{2}`$, the projection of the spin $`\stackrel{}{s}`$ on the momentum vector $`\stackrel{}{k}`$. As long as the interaction is purely electromagnetic, a polarization of the electron alone does not change the structure of the cross section, Eq. (16). However, new structures appear if both electron and nucleon are polarized. In particular the reaction $`\stackrel{}{e}+\stackrel{}{N}`$ anything is described by the cross section $`^\mathrm{?}`$
$$\frac{d\sigma }{d\mathrm{\Omega }}=\mathrm{\Gamma }[\sigma _T+\epsilon \sigma _L+P_eP_x\sqrt{2\epsilon (1\epsilon )}\sigma _{LT}^{}+P_eP_z\sqrt{1\epsilon ^2}\sigma _{TT}^{}],$$
(17)
where $`P_e=2h=\pm 1`$ refers to the helicity of the electron, and $`P_z`$ and $`P_x`$ are the longitudinal and transverse polarizations of the nucleon defined by the momentum of the virtual photon and an axis perpendicular to that direction (note: $`P_x`$ lies in the scattering plane of the electron and takes positive values on the side of the scattered electron).
In a more general experiment with production of pseudoscalar mesons, e.g. pions,
$$\stackrel{}{e}+\stackrel{}{N}e^{}+N^{}+\pi ,$$
(18)
up to 18 structure functions can be defined $`^\mathrm{?}`$, and this number increases further when higher spins are involved, e.g. if the electron is scattered on a deuteron target or if a vector particle (real photon, $`\rho `$ or $`\omega `$ meson etc.) is produced.
## 3 FORM FACTORS
Consider the absorption of a virtual photon with four-momentum q at an hadronic vertex. If the hadron stays intact after this process, i.e. in the case of elastic lepton scattering, the photon probes the expectation value of the hadronic vector current. If moreover the hadron is a scalar or pseudoscalar particle, the vector current has to be proportional to the two independent combinations of the 3 external four-momenta. Choosing $`q=p_2p_1`$ and $`P=(p_1+p_2)/2`$ as the independent vectors,
$$J_\mu :=p_2J_\mu p_1=F_1\frac{P_\mu }{m}+F_2\frac{q_\mu }{m}.$$
(19)
In this way we define two form factors, $`F_1`$ and $`F_2`$, which have to be scalars and as such may be expressed by functions of the independent Lorentz scalars that can be constructed. It is again a simple exercise to show that there exists only one independent scalar, e.g. $`Q^2=q^2`$, because $`Pq=0`$ and $`P^2=m^2\frac{1}{4}q^2`$ in the case of elastic scattering off a particle with mass $`m`$.
Next we can exploit the fact that the vector current of Eq. (19) is conserved, which follows from gauge invariance. The result is
$$0=q_\mu J^\mu =F_1\frac{p_2^2p_1^2}{2m}+F_2\frac{q^2}{m}.$$
(20)
Since $`p_1^2=p_2^2=m^2`$ for on-shell particles, the first term is zero and hence $`F_2`$ has to vanish identically. Therefore the vector current of, e.g., an on-shell pion has to take the form
$$J^\mu (\pi )=\frac{p_1^\mu +p_2^\mu }{2m_\pi }F_\pi (Q^2).$$
(21)
The form factor is normalized to $`F_\pi (0)=e_\pi `$, here and in the following in units of the elementary charge e. In this way we obtain, in the static limit $`q_\mu 0`$ and $`p_{2\mu }p_{1\mu }(m_\pi ,\stackrel{}{0})`$, the result $`J_\mu (e_\pi ,\stackrel{}{0})`$ for a charge $`e_\pi `$ at rest.
The situation is more complicated in the case of a particle with a spin like the nucleon, because now the independent momenta $`q`$ and $`P`$ can be combined with the familiar 16 independent $`4\times 4`$ matrices of Dirac’s theory: 1 (scalar), $`\gamma _5`$ (pseudoscalar), $`\gamma _\mu `$ (vector), $`i\gamma _5\gamma _\mu `$ (axial vector), and $`\sigma _{\mu \nu }`$ (antisymmetrical tensor). It is straightforward but somewhat tedious to show that the most general vector current of a spin-1/2 particle has to take the form
$$J_\mu :=p_2J_\mu p_1=\overline{u}_{p_2}\left(F_1\gamma _\mu +i\frac{F_2}{2m}\sigma _{\mu \nu }q^\nu \right)u_{p_1},$$
(22)
where $`u_{p1}`$ and $`u_{p2}`$ are the 4-spinors of the nucleon in the initial and final states respectively. The first structure on the $`rhs`$ is the Dirac current, which appears with the Dirac form factor $`F_1`$. The second term reflects the fact that due to its internal structure the particle acquires an anomalous magnetic moment $`\kappa `$, which appears with the Pauli form factor $`F_2`$. These form factors are normalized to $`F_1^p(0)=1`$, $`F_2^p(0)=\kappa _p=1.79`$ and $`F_1^n(0)=0`$, $`F_2^n(0)=\kappa _n=1.91`$ for proton and neutron respectively.
From the analogy with nonrelativistic physics, it is seducing to associate the form factors with the Fourier transforms of charge and magnetization densities. The problem is that a calculation of the charge distribution $`\rho (\stackrel{}{r})`$ involves a 3-dimensional Fourier transform of the form factor as function of $`\stackrel{}{q}`$, while in general the form factors are functions of $`Q^2=\stackrel{}{q}^2\omega ^2`$. However, there exists a special Lorentz frame, the Breit or brickwall frame, in which the energy of the virtual photon vanishes. This can be realized by choosing, e.g., $`\stackrel{}{p}_1=\stackrel{}{q}/2`$ and $`\stackrel{}{p}_2=+\stackrel{}{q}/2`$ leading to $`E_1=E_2=(m^2+\stackrel{}{q}^2/4)^{1/2}`$, $`\omega =0`$, and $`Q^2=\stackrel{}{q}^2`$.
In that frame the vector current takes the form
$$J_\mu =(G_E(Q^2),i\frac{\stackrel{}{\sigma }\times \stackrel{}{q}}{2m}G_M(Q^2)),$$
(23)
where $`G_E`$ stands for the time-like component of $`J_\mu `$ and hence is identified with the Fourier transform of the electric charge distribution, while $`G_M`$ appears with a structure typical for a static magnetic moment and hence is interpreted as Fourier transform of the magnetization density. The two “Sachs form factors” $`G_E`$ and $`G_M`$ are related to the Dirac form factors by $`^\mathrm{?}`$
$$G_E(Q^2)=F_1(Q^2)\tau F_2(Q^2),G_M(Q^2)=F_1(Q^2)+F_2(Q^2),$$
(24)
where $`\tau =Q^2/4m^2`$ is a measure of relativistic (recoil) effects. While Eq. (24) is taken as a general, covariant definition, the Sachs form factors can only be Fourier transformed in a special frame, namely the Breit frame, with the result
$`G_E(\stackrel{}{q}^2)`$ $`=`$ $`{\displaystyle \rho (\stackrel{}{r})e^{i\stackrel{}{q}\stackrel{}{r}}d^3\stackrel{}{r}}`$ (25)
$`=`$ $`{\displaystyle \rho (\stackrel{}{r})d^3\stackrel{}{r}}{\displaystyle \frac{\stackrel{}{q}^2}{6}}{\displaystyle \rho (\stackrel{}{r})\stackrel{}{r}^2d^3\stackrel{}{r}}+\mathrm{},`$
where the first integral yields the total charge in units of $`e`$, i.e. 1 for the proton and 0 for the neutron, and the second integral defines the square of the electric $`rms`$ radius, $`r^2_E:=r_E^2`$ of the particle. The interpretation of $`G_E`$ in terms of the charge distribution has recently been discussed again $`^\mathrm{?}`$.
We note that each value of $`Q^2`$ requires a particular Breit frame. Therefore, information has to be compiled from an infinity of different frames, which is then used as input for the Fourier integral for $`\rho (\stackrel{}{r})`$ in terms of $`G_E(\stackrel{}{q}^2)`$. Therefore, the density $`\rho (\stackrel{}{r})`$ is not an observable that we can “see” in any particular Lorentz frame but only a mathematical construct in analogy to a “classical” charge distribution. The problem is, of course, that due to the small mass of an “elementary” particle, recoil effects (measured by $`\tau `$) and size effects (measured by $`r^2`$) become comparable and cannot be separated in a unique way. This situation is numerically quite different in the case of a heavy nucleus for which the size effects dominate the recoil effects by many orders of magnitude!
The two Sachs form factors may be determined from the differential cross section
$$\frac{d\sigma }{d\mathrm{\Omega }}=\sigma _{\text{Mott}}\left(\frac{G_E^2+\tau G_M^2}{1+\tau }+2\tau \mathrm{tan}^2\frac{\theta }{2}G_M^2\right)$$
(26)
by means of a “Rosenbluth plot”, showing the cross section as function of $`\mathrm{tan}^2\frac{\theta }{2}`$ for constant $`Q^2`$. The data should lie on a straight line with a slope $`2\tau G_M^2`$, and the extrapolation to $`\tau =0`$ will determine the electric form factor $`G_E`$. Unfortunately, the Rosenbluth plot has a limited range of applicability. For decreasing $`Q^2`$, also $`\tau `$ and the slope become small and the error bars on $`G_M^2`$ increase. Large $`Q^2`$, on the other hand side, lead to a small electric contribution $`G_E^2/\tau `$ with large errors for the electric form factor.
In the case of the proton, the Rosenbluth plot was evaluated up to $`Q^2=8.8`$ GeV<sup>2</sup> at SLAC $`^\mathrm{?}`$. The results are shown in Fig. 5. Additional and more precise information can be obtained at the new electron accelerators by double-polarization experiments, in particular by target polarization $`\stackrel{}{p}(\stackrel{}{e},e^{})p`$ and recoil polarization, $`p(\stackrel{}{e},e^{})\stackrel{}{p}`$. The asymmetry $`A`$ measured by such an experiment is given by $`^\mathrm{?}`$
$$A=P_e\frac{\sqrt{2\tau \epsilon (1\epsilon )}G_EG_MP_x+\tau \sqrt{1\epsilon ^2}G_M^2P_z}{\epsilon G_E^2+\tau G_M^2},$$
(27)
where $`P_e`$ is the (longitudinal) polarization of the incident electron, and $`P_x`$ and $`P_z`$ are the transverse and longitudinal polarization components of the nucleon as defined in Eq. (17). In particular we find that the longitudinal-transverse interference term, appearing if the nucleon is polarized perpendicularly (sideways) to $`\stackrel{}{q}`$, will be proportional to $`G_EG_M`$, while the transverse-transverse interference term, appearing for polarization in the $`\stackrel{}{q}`$ direction, will be proportional to $`G_M^2`$. The ratio of both measurements then determines $`G_E/G_M`$ with high precision, because most normalization and efficiency factors will cancel.
Within the large error bars of the experiments, the older data followed surprisingly close the so-called “dipole fit” for the Sachs form factors,
$`G_E^p`$ $`=`$ $`G_M^p/\mu _p=G_M^n/\mu _n=(1+Q^2/M_V^2)^2:=G_D`$
$`G_E^n/\mu _n`$ $`=`$ $`\tau (1+Q^2/M_V^{}^2)^1(1+Q^2/M_V^2)^2:=G_P,`$ (28)
with $`\mu _p=2.79,\mu _n=1.91`$, $`M_V=840`$ MeV and $`M_V^{}=790`$ MeV. Since $`\tau =Q^2/4m^2,G_E^n(0)`$ vanishes, while $`G_E^p(0)=1`$ and the magnetic form factors approach the total magnetic moments for $`Q^20`$. In the asymptotic region $`Q^2\mathrm{}`$, all Sachs form factors should have a $`Q^4`$ behavior according to perturbative QCD. Inverting Eq. (24) we also find the asymptotic behavior of the Dirac form factors as required by pQCD, $`F_1Q^4`$ and $`F_2Q^6`$.
Already the SLAC experiments showed, however, that $`G_M^p/\mu _pG_D`$ falls much below unity at the higher momentum transfers $`^\mathrm{?}`$, reaching values of about 0.65 at $`Q^2=20`$ (GeV/c)<sup>2</sup>. For the reason pointed out before, $`G_E^p`$ was not well determined by these experiments. This situation has changed dramatically by the recent results from Jefferson Lab, which were obtained by scattering polarized electrons in coincidence with the polarization of the recoiling protons $`^\mathrm{?}`$. In this way it was possible to separate the form factors up to $`Q^2=3.5`$(GeV/c)<sup>2</sup>, where $`G_E^p/G_M^p`$ reaches the surprisingly low value of about 0.55, i.e. $`G_E^p`$ falls below $`G_D`$ even faster than $`G_M^p`$.
However, the situation is even more complex in the case of the neutron. The only exact information used to be the electric neutron radius, $`r^2_E^n0.11`$ fm<sup>2</sup>, which was obtained by scattering low energy neutrons off a $`{}_{}{}^{208}Pb`$ target $`^\mathrm{?}`$. Since there is no free neutron target, electron scattering data have to be obtained from light nuclei such as <sup>2</sup>H or <sup>3</sup>He making appropriate corrections for binding effects. This is a particularly difficult task for $`G_E^n`$, because it is smaller than the other form factors by a factor 10-20. In the past, results were obtained by either deuteron breakup in quasifree (neutron) kinematics $`^\mathrm{?}`$ or elastic scattering off the deuteron $`^\mathrm{?}`$, assuming that all other form factors and wave function corrections were well under control. Though the data reached a remarkable statistical accuracy, large systematical errors remained, particularly with regard to the nucleon-nucleon potential. While it had been pointed out long ago that double-polarization experiments should be much less model-dependent, such data were only taken very recently $`^\mathrm{?}`$. As shown in Fig. 6 the electric form factor of the neutron seems to be much larger than previously thought of. With the exception of the <sup>3</sup>He point at $`Q^20.35`$ (GeV/c)<sup>2</sup>, the new data follow the full line (“Mainz fit”) as opposed to the dashed line (“Saclay fit”, obtained from elastic $`ed`$ scattering). It is remarkable that the <sup>2</sup>H data point at the lowest $`Q^2`$ has moved upward by nearly a factor of 2 by taking account of final state interactions $`^\mathrm{?}`$, while these corrections are only at the percent level for the higher $`Q^2`$. This observation is at variance with the earlier assumption that final state interactions would not play any role in this kind of experiment. In view of this lesson from the deuteron it may be assumed that also the lowest <sup>3</sup>He data point will move once a complete calculation of final state and meson exchange effects exists.
The following Fig. 7 compares the neutron charge density obtained by Fourier transforming the older and the more recent data. Both results are in qualitative agreement with our expectation that the neutron charge density should have a positive core surrounded by a negative cloud $`^\mathrm{?}`$. The remarkable facts are, however, that the new data lead to a lower zero-crossing at r=0.7 fm in comparison with the older results (r=0.9 fm), and that both maximum and minimum become more pronounced. If one naively interprets the total negative charge as the pion cloud, one finds a probability of about 60 % that the neutron has a proton core surrounded by a $`\pi ^{}`$ cloud. Such an idea is quite natural for models of pions and nucleons, in particular for chiral bag models. It is interesting to note that a similar density is also predicted by the constituent quark model. The hyperfine interaction leading to the $`\mathrm{\Delta }`$-nucleon mass splitting predicts, at the same time, a stronger repulsion of quarks with equal flavor. Therefore the two $`d`$ quarks with total charge $`2/3`$ will move to the bag surface while the up quark goes to the center.
A final remark is in order concerning the proton radius. The experimental situation for $`r_E^p`$ is shown in Table 1. The large data spread for this very elementary quantity is truly surprising. The recent optical and radio frequency experiments had in mind, of course, to search for the limits of quantum electrodynamics by measuring Lamb shifts and hyperfine structures. In spite of an astounding accuracy of about 12 decimals, the analysis was stopped at about the 7th decimal by the existing uncertainties in the proton radius. If all deviations from theory are attributed to size effects, considerably larger radii are obtained than in the case of electron scattering.
The size of the nucleon is not just an academical question, but of tremendous consequence for our understanding of hadronic matter. The Table also shows the radius of an equivalent, homogeneous charge distribution, $`R_{eq}=(\frac{5}{3}r_E^2)^{1/2}`$, and the resulting “volume” of a nucleon. Obviously the volume grows by nearly 50 % by going from the Stanford value to the more recent results. Hence nucleons in a nucleus may get into a very uncomfortable environment: they need more space than is actually available. This situation is of course quite different from most models of nuclei and nuclear matter, which are based on effective interactions between point particles.
Note added in proof: In a recent paper Rosenfelder pointed out that Coulomb corrections will increase the proton radius, as measured by electron scattering, to $`r_E=(0.880\pm 0.015)`$ fm, with an error bar depending on the fit strategy $`^\mathrm{?}`$.
## 4 STRANGENESS
The strangeness content of the nucleon manifests itself by matrix elements $`<N\overline{s}\mathrm{\Gamma }sN>`$, with $`\mathrm{\Gamma }`$ any of the 5 Dirac structures of Section 3. Though observation of these matrix elements is necessarily proof for the existence of $`s`$ quarks in the nucleon, the strength of the 5 matrix elements may well be different. Since there exists no net strangeness in the nucleon, these observables open, in principle, a clear window on sea degrees of freedom. Information on the strange quark content comes essentially from three sources:
1. Deep inelastic lepton scattering.
The experiments clearly indicate a break-down of the Ellis-Jaffe sum rule based on SU(3) symmetry $`^\mathrm{?}`$. Such experiments have led to the so-called “spin crisis” of the nucleon, which was eventually explained by sea quark and gluon contributions to the spin of the nucleon. The observed symmetry breaking is proportional to the axial vector current carried by the $`s`$ quarks $`^\mathrm{?}`$.
2. Pion-nucleon scattering.
Dispersion analysis allows one to extrapolate $`\pi N`$ scattering to the (unphysical) Cheng-Dashen point at $`s=u`$ and $`t=2m_\pi ^2`$. The scattering amplitude at this point is essentially given by the $`\sigma `$ term $`^\mathrm{?}`$,
$$\sigma _{\pi N}=\frac{m_u+m_d}{2}<N\overline{u}u+\overline{d}dN>.$$
(29)
In combination with similar information on $`KN`$ scattering and approximate SU(3), the scalar $`\overline{s}s`$ condensate can be determined. The size of this effect is, unfortunately, not well established. The prediction of ChPT is $`^\mathrm{?}`$
$$\sigma _{\pi N}=58.3(10.56+0.33)\text{MeV}=45\text{MeV},$$
(30)
the 3 terms in this equation indicating the (slow) convergence of the perturbation series, while typical phenomenological analyses are in the range of $`\sigma _{\pi N}=(60\pm 20)`$ MeV. It is obvious that this uncertainty will also affect the value of the strangeness contribution, which therefore carries large error bars,
$$y=\frac{2<\overline{s}s>}{<\overline{u}u+\overline{d}d>}=0.21\pm 0.20.$$
(31)
3. Parity-violating lepton scattering.
The interest in this experiment stems from the observation that the photon and the $`Z^0`$ gauge boson couple differently to the vector currents of the quarks. An interference term of the electromagnetic and the weak neutral current is parity-violating (PV) and thus can be determined by a PV asymmetry. This presents the opportunity to measure a third form factor, in addition to the electromagnetic form factors of neutron and proton. If these 3 form factors would be known to sufficient precision, the density distributions of $`u,d`$ and $`s`$ quarks could be determined $`^\mathrm{?}`$.
The strange vector current $`<N^{}\overline{s}\gamma _\mu sN>`$ takes the general form of Eq. (22) with Dirac $`(F_1^s)`$ and Pauli $`(F_2^s)`$ strangeness form factors. Since the nucleon has no net strangeness, $`F_1^s(Q^2=0)=0`$. It follows from Eq. (24) that
$`G_E^s(Q^2)`$ $`=`$ $`{\displaystyle \frac{1}{6}}Q^2<r^2>_E^s+[Q^4],`$
$`G_M^s(Q^2)`$ $`=`$ $`\mu ^s+[Q^2],`$ (32)
with $`<r^2>_E^s`$ the square of the electric $`rms`$ radius and $`\mu ^s=\kappa ^s=F_2^s(0)`$ the (anomalous) magnetic moment due to the strange quark sea. Instead of the radius one often finds the dimensionless quantity $`\varrho ^s=dG/d\tau `$, the derivative of a particular form factor with regard to the quantity $`\tau `$, which is related to
$$<r^2>^s=0.066\text{fm}^2\varrho ^s.$$
(33)
The new information from parity-violating $`\stackrel{}{e}+Ne^{}+N^{}`$ can be obtained from the asymmetry $`A=(d\sigma ^+d\sigma ^{})/(d\sigma ^++d\sigma ^{})`$, where $`d\sigma ^+`$ and $`d\sigma ^{}`$ denote the cross sections for positive and negative helicities of the incident electron. While such an asymmetry must vanish in the purely electromagnetic case, it can appear by an interference between the leading electromagnetic and the much smaller (parity violating) weak interaction. Of course, the leading term is obtained by the absolute square of the amplitudes for photon exchange, resulting in a contribution $`O(e^4)`$, while the subleading term is given by the interference of photon exchange and $`Z^{}`$ exchange, which is $`O(e^2G_F^2)`$, with $`G_F`$ Fermi’s constant of weak interactions. While the photon couples only via the vector current, the $`Z^{}`$ can couple both to vector and axial currents. The interesting, parity violating interference term appears if the $`Z^{}`$ couples with the vector current to the nucleon and with the axial current to the electron or vice versa.
Table 2 shows, in standard notation, the vertices for the coupling of some leptons and quarks to photon and $`Z^{}`$,
where $`g^{}=e/4sc`$, $`s=\mathrm{sin}\theta _W`$, $`c=\mathrm{cos}\theta _W`$, and $`\theta _W`$ is the Weinberg angle, given by $`\mathrm{sin}^2\theta _W=0.2319\pm 0.0005`$. We observe that the coupling of the electron to the $`Z^{}`$ is dominated by the axial vector, because the vector part is suppressed by $`4\mathrm{sin}^2\theta _W1`$. By the same fact the vector currents of the quarks couple quite differently to photons and $`Z^{}`$ bosons, in particular the ratio of $`u`$ to $`d`$ or $`s`$ quark couplings reverses from -2 about $`\frac{1}{2}`$ if going from electromagnetic to weak neutral interactions.
The experimental information on strangeness is given by the asymmetry, which in the case of the nucleon takes the form
$`𝒜`$ $`=`$ $`{\displaystyle \frac{d\sigma ^+d\sigma ^{}}{d\sigma ^++d\sigma ^{}}}`$ (34)
$`=`$ $`{\displaystyle \frac{\frac{G_E\stackrel{~}{G}_E+\tau G_M\stackrel{~}{G}_M}{1+\tau }+2\tau G_M\stackrel{~}{G}_M\mathrm{tan}^2\frac{\theta }{2}+\mathrm{}(14s^2)G_M\stackrel{~}{G}_A}{\frac{G_E^2+\tau G_M^2}{1+\tau }+2\tau G_M^2\mathrm{tan}^2\frac{\theta }{2}}}`$
$`=`$ $`𝒜^E(\stackrel{~}{G}_E)+𝒜^M(\stackrel{~}{G}_M)+𝒜^A(\stackrel{~}{G}_A).`$
Though the 3 form factors $`\stackrel{~}{G}_E,\stackrel{~}{G}_M,\stackrel{~}{G}_A`$ can in principle be separated by a super Rosenbluth plot, definite results will take some time. The total asymmetry in a typical experiment is $`𝒜10^4Q^2/`$ (GeV/c)<sup>2</sup>, and only a small fraction of $`𝒜`$ is due to the expected effects of the strange quarks. According to Table 2 these effects can be obtained by the quark currents
$$J_\mu ^{(\gamma )}=e_q\overline{q}\gamma _\mu q,J_\mu ^{(Z_0)}=\stackrel{~}{e}_q\overline{q}\gamma _\mu q,$$
(35)
with $`e_u=2/3`$, $`e_d=e_s=1/3`$, $`\stackrel{~}{e}_u=1+8s^2/3`$ and $`\stackrel{~}{e}_d=\stackrel{~}{e}_s=14s^2/3`$, where $`s^2=\mathrm{sin}^2\theta _W`$. The matrix element of these quark currents between nucleon states can be parametrized by form factors describing the quark structure, e.g.
$$p^{}\overline{s}\gamma _\mu sp=G_s(Q^2)\overline{u}_p^{}\gamma _\mu u_p+\text{magnetic terms}.$$
(36)
The sum of the $`u`$, $`d`$, and $`s`$ quark contributions must equal the form factor of the nucleon,
$`G^p`$ $`=`$ $`{\displaystyle \frac{2}{3}}G_u{\displaystyle \frac{1}{3}}(G_d+G_s),`$
$`G^n`$ $`=`$ $`{\displaystyle \frac{2}{3}}G_d{\displaystyle \frac{1}{3}}(G_u+G_s),`$
$`\stackrel{~}{G}^p`$ $`=`$ $`\left(1+{\displaystyle \frac{8}{3}}s^2\right)G_u+\left(1{\displaystyle \frac{4}{3}}s^2\right)(G_d+G_s).`$ (37)
In these equations, $`G_{u/d/s}`$ are the quark distributions in the proton, and those of the neutron have been assumed to follow from isospin symmetry. If the 3 form factors on the $`lhs`$ of Eq. (37) have been measured, the strange quark contribution can be determined from
$$\stackrel{~}{G}^p=(14\mathrm{sin}^2\theta _W)G^p+G^n+G^s.$$
(38)
A particularly simple formula may be obtained for PV scattering $`^\mathrm{?}`$ off <sup>4</sup>He. Since this nucleus has spin zero, there exists only a charge monopole form factor. Furthermore <sup>4</sup>He is well described by an isoscalar system of nucleons having the same spatial wave functions. Under these assumptions the asymmetry may be cast into the form
$$𝒜(^4\text{He})=\frac{G_FQ^2}{\pi \sqrt{2}\alpha _{fs}}(\mathrm{sin}^2\theta _W+\frac{G_E^s}{2(G_E^p+G_E^n)}).$$
(39)
With $`G_F`$ the Fermi constant and $`\alpha _{fs}`$ the fine structure constant, the factor in front of the bracket is about $`410^4Q^2/`$ (GeV/c)<sup>2</sup>, and with the value of $`\theta _W`$ the “non strange” asymmetry is about $`10^4Q^2/`$ (GeV/c)<sup>2</sup>, which sets the scale for this difficult experiment. While earlier experiments on PV electron scattering $`^\mathrm{?}`$ were performed in order to determine $`\theta _W`$, which of course required that $`G_E^s`$ be negligible, the Weinberg angle is now known to 3 digits and today the motivation is to determine the strange quark contribution.
The simplest model for the strangeness contribution is, say, a proton that part of the time contains a strange pair,
$$p=u^2d+u^2ds\overline{s}+\mathrm{},$$
(40)
with the ellipse standing for $`u`$ and $`d`$ pairs and higher configurations. As long as $`s`$ and $`\overline{s}`$ quarks have the same spatial wave function, their charges cannot be seen by the electron. In order to separate the quarks in space however, the wave functions have to be correlated, the simplest long-range correlation being the clustering of the second component in Eq. (40) in the form of $`\mathrm{\Lambda }(uds)K^+(u\overline{s})`$. This model will therefore predict, as contribution of the strange sea, a positively charged cloud $`(K^+)`$ and a negative core (the neutral $`\mathrm{\Lambda }`$ relative to the charged $`p`$$`^\mathrm{?}`$. As a result both the anomalous magnetic moment of the proton, $`\kappa _p`$, and the value of $`r^2_E^p`$ will be increased. Since the $`s`$ quark has negative charge, this model predicts $`\mu ^s=\kappa ^s<0`$ and $`r^2_E^s<0`$ for the quantities introduced in Eq. (32).
A second model is based on dispersion relations, which tend to predict a strong contribution of the $`\mathrm{\Phi }`$(1020) in order to combine with the $`\omega `$(780) to an approximate dipole form of the isoscalar form factors $`^\mathrm{?}`$. Since the $`\mathrm{\Phi }`$ is practically an $`s\overline{s}`$ configuration, its appearance is related with strangeness in the nucleon. Other calculations have been performed in Skyrme, chiral quark-soliton and constituent quark models, and in the framework of lattice QCD and ChPT. Such calculations generally result in negative values for $`\mu ^s`$ with a range of $`.3\mu ^s.7`$, while $`r^2_E^s0.15`$ fm<sup>2</sup> in dispersion models ($`\mathrm{\Phi }`$ poles) and $`0r^2_E^s0.15`$ fm<sup>2</sup> for $`K`$ loops.
The recent results of the SAMPLE experiment at MIT/Bates and of HAP-PEX at Jefferson Lab came as a big surprise: The $`s`$ quark contribution is much smaller than predicted, and in fact even compatible with zero. The SAMPLE experiment measured essentially $`G_M^s`$, which came out positive though with large error bars $`^\mathrm{?}`$. Extrapolating to $`G_M^s(0)=\mu ^s`$, Hemmert et al. $`^\mathrm{?}`$ obtained $`0.03<\mu _p^s<0.18`$ by use of the slope of $`G_M^s(Q^2)`$ as predicted from HBChPT (note that this theory cannot predict $`\mu ^s`$ itself, because of an unknown low energy constant). The HAPPEX collaboraton obtained a raw asymmetry $`𝒜=5.64\pm 0.75`$ ppm. Since most of this asymmetry was expected on the basis of $`u`$ and $`d`$ quarks, only a small fraction remained as possible $`s`$ quark contribution, leading to the result $`^\mathrm{?}`$
$$G_E^s+0.39G_M^s=0.023\pm 0.034\pm 0.022\pm 0.026$$
(41)
at $`Q^2=0.48`$ (GeV/c)<sup>2</sup>. The error bars in Eq. (41) denote, in order of appearance, the statistical and systematical uncertainties as well as the errors due to our bad knowledge of the neutron from factor $`G_E^n`$ at that momentum transfer. The result is again positive though with large error bars, and taken at face value it rules out most theoretical predictions. A selection of these predictions can be found in Ref. $`^\mathrm{?}`$. Contrary to earlier lattice QCD predictions, a recent lattice calculation finds small negative values $`G_M^s(0)=0.16\pm 0.18`$, which could even shift to more positive values because of systematic errors $`^\mathrm{?}`$. From a comparison of recent data obtained for proton and deuteron targets, it has been suspected that the hadronic radiative corrections to the axial form factor are not yet under control. In view of the importance of this topic, more and new experiments on the strange form factor are underway $`^\mathrm{?}`$.
## 5 COMPTON SCATTERING
The polarizability measures the response of a particle to a quasistatic electromagnetic field. In particular the energy is generally lowered by
$$\mathrm{\Delta }E=\frac{1}{2}\alpha \stackrel{}{E}^2\frac{1}{2}\beta \stackrel{}{H}^2,$$
(42)
where $`\stackrel{}{E}`$ and $`\stackrel{}{H}`$ are the electric and magnetic fields, and $`\alpha `$ and $`\beta `$ the electric and magnetic polarizabilities. In the case of a macroscopic system with N atoms per volume, the polarizabilities are related to the dielectric constant $`\epsilon `$ and the magnetic permeability $`\mu `$ by
$$\epsilon =1N\alpha ,\mu =1N\beta .$$
(43)
The electric polarizability of a metal sphere is essentially given by its volume, it scales with the third power of the radius. In the case of a dielectric sphere an additional factor $`(\epsilon 1)/(\epsilon +2)`$ appears, which reduces the polarizability by orders of magnitude, because $`\epsilon `$ is close to unity. The same is true for the nucleon. If we divide its polarizability by the volume $`V`$, we obtain
$$\frac{\alpha }{V}\frac{10^3\text{fm}^3}{\frac{4}{3}\pi \text{fm}^3}210^4,$$
(44)
i.e. the nucleon is a very rigid object. It is held together by strong interactions, and the applied electromagnetic field cannot easily deform the charge distribution. Of course the nucleon cannot be polarized by putting it between two condensator plates. Instead its polarizability can be measured by Compton scattering: The incoming photon deforms the nucleon, and by measuring the energy and angular distributions of the outgoing photon one can determine the polarizability.
In nonrelativistic quantum mechanics the electric polarizability is given by
$$4\pi \alpha =2\mathrm{}^2\underset{n>0}{}\frac{n\widehat{D}_z0^2}{E_nE_0},$$
(45)
where $`\widehat{D}_z=e\widehat{z}`$ is the dipole operator and $`e^2/4\pi 1/137.`$ Since all excitation energies $`E_nE_0`$ are positive and only the modulus of the transition matrix element $`n\widehat{D}_z0`$ enters, $`\alpha `$ has to be positive in a nonrelativistic model.
Here is a simple prototype problem for a polarizable system $`^\mathrm{?}`$. A nonrelativistic particle with mass $`M`$ and charge $`Q`$ is held by a harmonic oscillator potential with Hooke’s constant $`C=M\omega _0^2`$. If we apply an external electrical field $`\stackrel{}{E}`$, the Hamiltonian is
$$H=\frac{\stackrel{}{p}^2}{2M}+\frac{M\omega _0^2}{2}\stackrel{}{r}^2+Q\stackrel{}{E}\stackrel{}{r},$$
(46)
which can be cast into the form
$$H=\frac{\stackrel{}{p}^2}{2M}+\frac{M\omega _0^2}{2}\left(\stackrel{}{r}+\frac{Q}{M\omega _0^2}\stackrel{}{E}\right)^2\frac{1}{2}\frac{Q^2}{M\omega _0^2}\stackrel{}{E}^2.$$
(47)
The result is
1. a shift in space, $`\mathrm{\Delta }\stackrel{}{r}=\frac{Q}{M\omega _0^2}\stackrel{}{E}`$, leading to an induced dipole moment $`\stackrel{}{d}=Q\mathrm{\Delta }\stackrel{}{r}:=\alpha \stackrel{}{E}`$, and
2. a shift in energy, $`\mathrm{\Delta }E=\frac{1}{2}\alpha \stackrel{}{E}^2`$, with $`\alpha =\frac{Q^2}{M\omega _0^2}`$.
In view of several misrepresentations in the literature, we stress the point that these two definitions of $`\alpha `$, via induced dipole moment or energy shift, should lead to the same value.
A more generic model involves two particles (masses $`M_1`$ and $`M_2`$, charges $`Q_1`$ and $`Q_2`$), held together with a spring constant $`C=\mu \omega _0^2`$, where $`\mu `$ is the reduced mass. An external field $`\stackrel{}{E}`$ induces both an intrinsic dipole moment (expressed in terms of the relative coordinate) and an acceleration of the center of mass. According to classical antenna theory, the scattering amplitude $`f(\omega )`$ is proportional to the acceleration of the induced dipole moments. The final result is $`^\mathrm{?}`$
$`f(\omega )`$ $`=`$ $`{\displaystyle \frac{Q^2}{M}}+{\displaystyle \frac{1}{\mu (\omega _0^2\omega ^2)}}\left({\displaystyle \frac{Q_1M_2Q_2M_1}{M}}\right)^2\omega ^2`$ (48)
$`=`$ $`{\displaystyle \frac{Q^2}{M}}+4\pi \alpha (\omega )\omega ^2.`$
In the limit of $`\omega 0`$, the scattering amplitude reduces to the Thomson term depending only on the total charge $`Q`$ and the total mass $`M`$ of the system. It is the essence of more refined “low energy theorems” (LET) that only such global properties should be visible in that limit. Since the cross section $`d\sigma /d\mathrm{\Omega }f(\omega )^2`$, the internal structure shows up first at $`O(\omega ^2)`$, as interference of the Thomson term with the second term in Eq. (48). In the case of a globally neutral system (e.g. a neutral atom, a neutron or a $`\pi ^0`$), the Thomson term vanishes and the cross section starts at $`O(\omega ^4)`$. This is the familiar case of Rayleigh scattering leading to the blue sky, because most gases absorb in the ultra-violet, $`\omega _0^2\omega ^2`$, with $`\omega `$ a frequency of visible light. If $`\omega `$ increases further, it approaches a singularity in Eq. (48), which is of course avoided by appropriate friction terms, i.e. by a width $`\mathrm{\Gamma }_0`$ of the resonance at $`\omega _0`$.
Compton scattering off the proton is, of course, technically much more complicated than the nonrelativistic model above. The reasons are relativity and the spin degrees of freedom. By use of Lorentz and gauge invariance, crossing symmetry, parity and time reversal invariance, the general Compton amplitude takes the form $`^\mathrm{?}`$
$$T=\epsilon _\mu ^{}\epsilon _\nu \underset{i=1}{\overset{6}{}}𝒪_i^{\mu \nu }\stackrel{~}{A}_i(s,t),$$
(49)
where $`𝒪_i^{\mu \nu }`$ are Lorentz tensors constructed from kinematical variables and $`\gamma `$ matrices, and $`\stackrel{~}{A}_i`$ are Lorentz scalars. In the $`cm`$ frame, these Lorentz structures can be reduced to Pauli matrices combined with unit vectors in the directions of the initial $`(\widehat{k})`$ and final $`(\widehat{k}^{})`$ photons, which yields the result $`^{\mathrm{?},\mathrm{?}}`$
$`T`$ $`=`$ $`A_1(\omega ,t)\stackrel{}{ϵ}^{}\stackrel{}{ϵ}+A_2(\omega ,t)\stackrel{}{ϵ}{}_{}{}^{}{}_{}{}^{}\widehat{k}\stackrel{}{ϵ}\widehat{k}^{}`$ (50)
$`+iA_3(\omega ,t)\stackrel{}{\sigma }(\stackrel{}{ϵ}{}_{}{}^{}{}_{}{}^{}\times \stackrel{}{ϵ})+iA_4(\omega ,t)\stackrel{}{\sigma }(\widehat{k}{}_{}{}^{}\times \widehat{k})\stackrel{}{ϵ}{}_{}{}^{}{}_{}{}^{}\stackrel{}{ϵ}`$
$`+iA_5(\omega ,t)\stackrel{}{\sigma }[(\stackrel{}{ϵ}{}_{}{}^{}{}_{}{}^{}\times \widehat{k})\stackrel{}{ϵ}\widehat{k}{}_{}{}^{}(\stackrel{}{ϵ}\times \widehat{k}^{})\stackrel{}{ϵ}{}_{}{}^{}{}_{}{}^{}\widehat{k}]`$
$`+iA_6(\omega ,t)\stackrel{}{\sigma }[(\stackrel{}{ϵ}{}_{}{}^{}{}_{}{}^{}\times \widehat{k}^{})\stackrel{}{ϵ}\widehat{k}{}_{}{}^{}(\stackrel{}{ϵ}\times \widehat{k})\stackrel{}{ϵ}{}_{}{}^{}{}_{}{}^{}\widehat{k}],`$
with $`\widehat{ϵ}`$ and $`\widehat{ϵ}^{}`$ describing the polarization of the photon in the initial and final states, and $`\stackrel{}{\sigma }`$ the spin of the nucleon.
The low energy theorem predicts the following threshold behavior for the proton amplitudes $`^\mathrm{?}`$:
$`A_1`$ $`=`$ $`{\displaystyle \frac{e^2}{m}}+4\pi (\alpha +\beta \mathrm{cos}\theta )\omega ^2{\displaystyle \frac{e^2}{4m^3}}(1\mathrm{cos}\theta )\omega ^2+\mathrm{},`$
$`A_2`$ $`=`$ $`{\displaystyle \frac{e^2}{m}}\omega 4\pi \beta \omega ^2+\mathrm{},`$
$`A_3`$ $`=`$ $`[(1+2\kappa )(1\mathrm{cos}\theta )\kappa ^2\mathrm{cos}\theta ]{\displaystyle \frac{e^2\omega }{2m^2}}{\displaystyle \frac{(2\kappa +1)e^2}{8m^4}}\mathrm{cos}\theta \omega ^3`$
$`+4\pi [\gamma _1(\gamma _2+2\gamma _4)\mathrm{cos}\theta ]\omega ^3+\mathrm{},`$
$`A_4`$ $`=`$ $`{\displaystyle \frac{(1+\kappa )^2e^2}{2m^2}}\omega +4\pi \gamma _2\omega ^3+\mathrm{},`$
$`A_5`$ $`=`$ $`{\displaystyle \frac{(1+\kappa )^2e^2}{2m^2}}\omega +4\pi \gamma _4\omega ^3+\mathrm{},`$
$`A_6`$ $`=`$ $`{\displaystyle \frac{(1+\kappa )^2e^2}{2m^2}}\omega +4\pi \gamma _4\omega ^3+\mathrm{}.`$ (51)
In the expansion for $`A_1`$ we recover the previously discussed low energy theorem for forward scattering. In addition to $`\alpha `$, however, also the magnetic polarizability $`\beta `$ appears. Since $`\alpha `$ and $`\beta `$ enter differently in $`A_1`$ and $`A_2`$, they can be determined separately by Compton scattering. The amplitudes $`A_1`$ and $`A_2`$ are typical for a scalar (or pseudoscalar) particle, and for this reason we call $`\alpha `$ and $`\beta `$ the scalar polarizabilities. Since the nucleon has a spin, there appear 4 more amplitudes, $`A_3`$ to $`A_6`$, whose leading terms , $`𝒪(\omega )`$, are related to the magnetic moment $`\mu =1+\kappa `$. The subleading terms, $`𝒪(\omega ^3)`$, define 4 new polarizabilities $`\gamma _1`$ to $`\gamma _4`$, the spin or vector polarizabilities of the nucleon. We recall that the differential cross section for small $`\omega `$ is dominated by the Thomson term and that the polarizabilities $`\alpha `$ and $`\beta `$ appear in the cross section at $`𝒪(\omega ^2)`$ via the interference of Thomson and Rayleigh scattering. In addition, however, also the spin-dependent amplitudes contribute at $`𝒪(\omega ^2)`$ for unpolarized Compton scattering, because without polarization the terms with and without the $`\stackrel{}{\sigma }`$ matrices add incoherently in the cross section. For the same reason the spin polarizabilities show up only at $`𝒪(\omega ^4)`$, i.e. are expected to be small and difficult to disentangle from other higher order terms. It is therefore obvious that the 6 polarizabilities cannot be determined from differential cross section measurements only, but that polarization experiments are necessary, in particular the scattering of circularly polarized photons off polarized protons.
In the following we shall again restrict the discussion to forward scattering, i.e. $`\widehat{k}^{}=\widehat{k}`$ or $`\theta =0`$. Due to the transversality condition $`\widehat{ϵ}\widehat{k}=\widehat{ϵ}^{}\widehat{k}^{}`$, only the amplitudes $`A_1`$ and $`A_3`$ contribute in that limit. With the notation $`A_1(\omega ,0)=f(\omega )`$ and $`A_3(\omega ,0)=g(\omega )`$, Eq. (50) can be cast into the form
$$T(\omega ,\theta =0)=\widehat{ϵ}^{}\widehat{ϵ}f(\omega )+i(\widehat{ϵ}^{}\times \widehat{ϵ})\stackrel{}{\sigma }g(\omega ).$$
(52)
Due to the crossing symmetry, the non spin-flip amplitude $`f(\omega )`$ is an even function in $`\omega `$ and the spin-flip amplitude $`g(\omega )`$ is odd. The 2 scattering amplitudes can be determined by scattering circularly polarized photons (spin projection +1) off nucleons polarized in the direction or opposite to the photon momentum (spin projections +1/2 or -1/2), leading to intermediate states with spin projection +3/2 or +1/2 respectively. Denoting the corresponding scattering amplitudes by $`T_{3/2}`$ and $`T_{1/2}`$, we find $`f(\omega )=(T_{1/2}+T_{3/2})/2`$ and $`g(\omega )=(T_{1/2}T_{3/2})/2`$. The optical theorem allows us to express the imaginary parts of $`f`$ and $`g`$ by the sum and difference of the helicity cross sections for physically allowed values of $`\omega `$,
$`\text{Im}f(\omega )`$ $`=`$ $`{\displaystyle \frac{\omega }{4\pi }}{\displaystyle \frac{\sigma _{1/2}+\sigma _{3/2}}{2}}={\displaystyle \frac{\omega }{4\pi }}\sigma _{tot}(\omega )`$
$`\text{Im}g(\omega )`$ $`=`$ $`{\displaystyle \frac{\omega }{4\pi }}{\displaystyle \frac{\sigma _{1/2}\sigma _{3/2}}{2}}={\displaystyle \frac{\omega }{4\pi }}\mathrm{\Delta }\sigma (\omega ).`$ (53)
We further assume that $`f`$ obeys a once-subtracted and $`g`$ an unsubtracted dispersion relation. Finally, we shall restrict the discussion to photon energies below pion threshold $`\omega _0`$, in which case the amplitudes are real and the dispersion relations can be cast into the form
$`4\pi f(\omega )`$ $`=`$ $`4\pi f(0)+{\displaystyle \frac{2\omega ^2}{\pi }}{\displaystyle _{\omega _0}^{\mathrm{}}}{\displaystyle \frac{\sigma _{tot}(\omega ^{})}{\omega ^2\omega ^2}}𝑑\omega ^{},`$ (54)
$`4\pi g(\omega )`$ $`=`$ $`{\displaystyle \frac{2\omega }{\pi }}{\displaystyle _{\omega _0}^{\mathrm{}}}{\displaystyle \frac{\mathrm{\Delta }\sigma (\omega ^{})}{\omega ^{}(\omega ^2\omega ^2)}}𝑑\omega ^{},`$
which involves integrations from the physical threshold for pion production, $`\omega _0`$, to infinity.
Next we make use of the low-energy theorem $`^\mathrm{?}`$, which allows us to express the low-energy behavior of $`f(\omega )`$ and $`g(\omega )`$ by a power series according to Eq. (5),
$`4\pi f(\omega )`$ $`=`$ $`{\displaystyle \frac{e^2}{m}}+4\pi (\alpha +\beta )\omega ^2+[\omega ^4],`$ (55)
$`4\pi g(\omega )`$ $`=`$ $`{\displaystyle \frac{2\pi e^2\kappa ^2}{m^2}}+4\pi \gamma _0\omega ^3+[\omega ^5].`$
If we compare Eqs. (54) and (55), we obtain a series of sum rules, in particular Baldin’s sum rule $`^\mathrm{?}`$
$$\alpha +\beta =\frac{1}{2\pi ^2}_{\omega _0}^{\mathrm{}}\frac{\sigma _{tot}(\omega )}{\omega ^2}𝑑\omega ,$$
(56)
the sum rule of Gerasimov, Drell and Hearn $`^\mathrm{?}`$,
$$\kappa ^2=\frac{2m^2}{\pi e^2}_{\omega _0}^{\mathrm{}}\frac{\sigma _{1/2}(\omega )\sigma _{3/2}(\omega )}{\omega }𝑑\omega ,$$
(57)
and a value for the forward spin polarizability $`^\mathrm{?}`$,
$$\gamma _0=\frac{1}{4\pi ^2}_{\omega _0}^{\mathrm{}}\frac{\sigma _{1/2}(\omega )\sigma _{3/2}(\omega )}{\omega ^3}𝑑\omega .$$
(58)
Both the forward spin polarizability $`\gamma _0`$ and the GDH sum rule depend on the difference of the helicity cross sections,
$$\sigma _{1/2}\sigma _{3/2}|E_{0^+}|^2|M_{1^+}|^2+E_{1^+}^{}M_{1^+}+\mathrm{},$$
(59)
i.e. are dominated by the difference of s-wave pion production (multipole $`E_{0^+}`$) and magnetic excitation of the $`\mathrm{\Delta }`$(1232) resonance (multipole $`M_{1^+}`$).
With the advent of high duty-factor electron accelerators and laser backscattering techniques, new Compton data have been obtained in the 90’s $`^\mathrm{?}`$ and more experiments are expected in the near future. The presently most accurate values for the proton polarizabilities were derived from the work of MacGibbon et al. $`^\mathrm{?}`$ whose experiments were performed with tagged photons at 70 MeV$`\nu 100`$ MeV and untagged ones at the higher energies, and analyzed in collaboration with L’vov $`^\mathrm{?}`$ by means of dispersion relations (in the following denoted by DR) at constant $`t`$. The results were
$`\alpha `$ $`=`$ $`\left(12.1\pm \mathrm{\hspace{0.17em}0.8}\pm \mathrm{\hspace{0.17em}0.5}\right)\times \mathrm{\hspace{0.17em}10}^4\mathrm{f}m^3,`$
$`\beta `$ $`=`$ $`\left(2.1\mathrm{\hspace{0.17em}0.8}\mathrm{\hspace{0.17em}0.5}\right)\times \mathrm{\hspace{0.17em}10}^4\mathrm{f}m^3.`$ (60)
The physics of the $`\mathrm{\Delta }`$(1232) and higher resonances has been the objective of further recent investigations with tagged photons at Mainz $`^\mathrm{?}`$ and with laser-backscattered photons at Brookhaven $`^\mathrm{?}`$. Such data were used to give a first prediction for the so-called backward spin polarizability of the proton $`^\mathrm{?}`$, i.e. the particular combination $`\gamma _\pi =\gamma _1+\gamma _2+2\gamma _4`$ entering the Compton spin-flip amplitude at $`\theta =180^{}`$,
$$\gamma _\pi =\left[27.1\pm \mathrm{\hspace{0.17em}2.2}(\mathrm{s}tat+syst)\genfrac{}{}{0pt}{}{+2.8}{2.4}(\mathrm{m}odel)\right]\times \mathrm{\hspace{0.17em}10}^4\mathrm{f}m^4.$$
(61)
In 1991 Bernard et al. $`^\mathrm{?}`$ evaluated the one-loop contributions to the polarizabilities in the framework of relativistic chiral perturbation theory (ChPT), with the result $`\alpha =10\beta =12.1`$ (here and in the following, the scalar polarizabilities are given in units of $`10^4`$ fm<sup>3</sup> and the spin polarizabilities in units of $`10^4`$ fm<sup>4</sup>). In order to have a systematic chiral power counting, the calculation was then repeated in heavy baryon ChPT, the expansion parameter being an external momentum or the quark mass. To $`O(p^4)`$ the result is $`\alpha =10.5\pm 2.0`$ and $`\beta =3.5\pm 3.6`$, the errors being due to 4 counter terms, which were estimated by resonance saturation $`^\mathrm{?}`$. One of these counter terms describes the paramagnetic contribution of the $`\mathrm{\Delta }`$(1232), which is partly cancelled by large diamagnetic contributions of pion-nucleon loops. In view of the importance of the $`\mathrm{\Delta }`$ resonance, Hemmert et al. proposed to include the $`\mathrm{\Delta }`$ as a dynamical degree of freedom. This added a further expansion parameter, the difference of the $`\mathrm{\Delta }`$ and nucleon masses (“$`ϵ`$ expansion”). A calculation to $`O(ϵ^3)`$ yielded $`\alpha `$ = 12.2 + 0 + 4.2 = 16.4 and $`\beta `$ = 1.2 + 7.2 + 0.7 = 9.1, the 3 separate terms referring to contributions of pion-nucleon loops (identical to the predictions of the $`O(p^3)`$ calculation), $`\mathrm{\Delta }`$-pole terms, and pion-$`\mathrm{\Delta }`$ loops $`^{\mathrm{?},\mathrm{?}}`$. These $`O(ϵ^3)`$ predictions are clearly at variance with the data, in particular $`\alpha +\beta =25.5`$ is nearly twice the rather precise value determined from DR (see below).
The spin polarizabilities have been calculated in both relativistic one-loop ChPT $`^\mathrm{?}`$ and heavy baryon ChPT $`^\mathrm{?}`$. In the latter approach the predictions are $`\gamma _0=4.62.40.2+0=+2.0,`$ (forward spin polarizability) and $`\gamma _\pi =4.6+2.40.243.5=36.7`$ (backward spin polarizability), the 4 separate contributions referring to N$`\pi `$-loops, $`\mathrm{\Delta }`$-poles, $`\mathrm{\Delta }\pi `$-loops, and the triangle anomaly, in that order. It is obvious that the anomaly or $`\pi ^0`$-pole gives by far the most important contribution to $`\gamma _\pi `$, and that it would require surprisingly large higher order contributions to increase $`\gamma _\pi `$ to the value of Ref. $`^\mathrm{?}`$. Similar conclusions were reached in the framework of DR. Using DR at $`t`$ = const, Ref. $`^\mathrm{?}`$ obtained a value of $`\gamma _\pi =34.3`$, while L’vov and Nathan $`^\mathrm{?}`$ worked in the framework of backward DR and predicted $`\gamma _\pi =39.5\pm 2.4`$.
As we have stated before, the most quantitative analysis of the experimental data has been provided by DR. In this way it has been possible to reconstruct the forward non spin-flip amplitude directly from the total photoabsorption cross section by Baldin’s sum rule, which yields a rather precise value for the sum of the scalar polarizabilities
$`\alpha +\beta `$ $`=`$ $`14.2\pm \mathrm{\hspace{0.17em}0.5}(\mathrm{Ref}.{}_{}{}^{54})`$ (62)
$`=`$ $`13.69\pm \mathrm{\hspace{0.17em}0.14}(\mathrm{Ref}.{}_{}{}^{55}).`$
Similarly, the forward spin polarizability can be evaluated by an integral over the difference of the absorption cross sections in states with helicity 3/2 and 1/2,
$`\gamma _0=\gamma _1\gamma _22\gamma _4`$ $`=`$ $`1.34(\mathrm{Ref}.{}_{}{}^{56})`$ (63)
$`=`$ $`0.6(\mathrm{Ref}.{}_{}{}^{52}).`$
The difference can be traced back to the s-wave threshold amplitude $`E_{0^+}(\gamma pn\pi ^+)`$, which used to be $`24.910^3/m_\pi `$ for the SAID $`^\mathrm{?}`$ and is $`28.310^3/m_\pi `$ for the HDT $`^\mathrm{?}`$ multipoles, the latter value agreeing well with the prediction of ChPT, $`28.410^3/m_\pi `$ $`^\mathrm{?}`$. While these predictions relied on pion photoproduction multipoles, the helicity cross sections have now been directly determined by scattering photons with circular polarizations on polarized protons $`^\mathrm{?}`$.
In view of the somewhat inconclusive situation, we are waiting for the new MAMI data for Compton scattering on the proton in and above the $`\mathrm{\Delta }`$-resonance region and over a wide angular range that have been reported preliminarily $`^\mathrm{?}`$. These new data will be most valuable to check the consistency of pion photoproduction and previous Compton scattering results obtained at LEGS, MAMI and other facilities.
Finally, in Fig. 8 we show the potential of double- polarization observables for measuring the spin polarizabilities $`^\mathrm{?}`$. In particular, an experiment with a circularly polarized photon and a polarized proton target should be quite sensitive to the backward spin polarizability $`\gamma _\pi `$, especially at energies between pion threshold and the $`\mathrm{\Delta }`$ resonance. In addition, possible normalization problems can be avoided by measuring appropriate asymmetries. Therefore such polarization experiments hold the promise to disentangle scalar and vector polarizabilities of the nucleon and to quantify the nucleon spin response in an external electromagnetic field.
## 6 PION PHOTOPRODUCTION
The reaction
$$\gamma ^{}(q)+N(p_1)\pi (p)+N(p_2)$$
(64)
is described by a transition matrix element $`\epsilon ^\mu J_\mu `$, with $`\epsilon ^\mu `$ the polarization of the (virtual) photon and $`J_\mu `$ a transition current. This current can be expressed by 6 different Lorentz structures constructed from the independent momenta $`p`$, $`q`$ and $`P=(p_1+p_2)/2`$ and appropriate Dirac matrices. Since the photon couples to the vector current and the pion is pseudoscalar, this transition current has the structure of an axial vector. Written in the $`cm`$ frame, its spacelike ($`\stackrel{}{J}`$) and the timelike ($`\rho `$) components take the form
$`\stackrel{}{J}`$ $`=`$ $`\stackrel{~}{\sigma }F_1+i(\widehat{q}\times \stackrel{}{\sigma })(\stackrel{}{\sigma }\widehat{p})F_2+\stackrel{~}{p}(\stackrel{}{\sigma }\widehat{q})F_3`$ (65)
$`+\stackrel{~}{p}(\stackrel{}{\sigma }\widehat{p})F_4+\widehat{q}(\stackrel{}{\sigma }\widehat{q})F_5+\widehat{q}(\stackrel{}{\sigma }\widehat{p})F_6,`$
$`\rho `$ $`=`$ $`(\stackrel{}{\sigma }\widehat{p})F_7+(\stackrel{}{\sigma }\widehat{q})F_8,`$ (66)
where $`F_1`$ to $`F_8`$ are the CGLN amplitudes $`^\mathrm{?}`$. The structures in front of $`F_1`$ to $`F_6`$ and $`F_7`$ to $`F_8`$ are the axial vectors and pseudoscalars that can be constructed from the $`\stackrel{}{\sigma }`$ matrix and the independent $`cm`$ momenta $`\stackrel{}{p}`$ and $`\stackrel{}{q}`$. We note that $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{p}`$ are the transverse components of $`\stackrel{}{\sigma }`$ and $`\widehat{p}`$, respectively, with regard to $`\widehat{q}`$. With these definitions $`F_1`$ to $`F_4`$ describe the transverse, $`F_5`$ to $`F_6`$ the longitudinal and $`F_7`$ to $`F_8`$ the timelike components of the current. The latter ones are related by current conservation, $`\stackrel{}{q}\stackrel{}{J}\omega \rho =0`$, leading to $`\stackrel{}{q}F_5=\omega F_8`$ and $`\stackrel{}{q}F_6=\omega F_7`$.
The CGLN amplitudes can be decomposed into a series of multipoles $`^\mathrm{?}`$,
$$\{_{l\pm }\}=\{E_{l\pm },M_{l\pm },L_{l\pm }\},$$
(67)
where $`E`$ and $`M`$ denote the transverse electric and magnetic multipoles, and $`L`$ are the longitudinal ones related to scalar (timelike, Coulomb) multipoles $`S`$ by current conservation. These multipoles are complex functions of 2 variables, e.g. $`=(Q^2,W)`$.
The notation of the multipoles is clarified by Fig. 9. The incoming photon carries the multipoles $`EL`$, $`ML`$ and $`SL`$, which are contructed from its spin 1 and the orbital angular momentum. The parity of these multipoles is $`𝒫=(1)^L`$ for $`E`$ and $`S`$, and $`𝒫=(1)^{L+1}`$ for $`M`$. The photon is now coupled to the nucleon with spin $`1/2`$ and $`𝒫=+1`$, which leads to intermediate states with spin $`J=L\pm \frac{1}{2}`$ and the parity of the incoming photon. The outgoing pion has negative intrinsic parity and orbital angular momentum $`l`$, from which we can reconstruct the spin $`J=l\pm \frac{1}{2}`$ and parity $`𝒫=(1)^{l+1}`$ of the intermediate state. This explains the notation of the multipoles, Eq. (67), by the symbols $`E`$, $`M`$ and $`S`$ referring to the type of the photon, and by the index $`l\pm `$ with $`l`$ standing for the pion momentum and the $`\pm `$ sign for the two possibilities to construct the total spin $`J=l\pm \frac{1}{2}`$ in the intermediate states. This notation completely defines the transition, in particular it determines the electromagnetic multipoles and the quantum numbers of the intermediate states.
Let us consider as an example the excitation of the $`\mathrm{\Delta }`$(1232) with the spectroscopic notation $`P_{33}`$. This intermediate state contains a pion in a $`p`$ wave, i.e. $`l=1`$ and $`𝒫=+1`$. The indices “33” refer to isospin $`I=3/2`$ and spin $`J=3/2`$ respectively. The N$`\mathrm{\Delta }`$ transition can therefore take place by $`M1`$ or $`E2`$ photons, for virtual photons also $`S2`$ is allowed. The phase $`\delta _{l\pm }^I`$ of the pion-nucleon final state is $`\delta _{1+}^{3/2}`$, and the photoproduction multipoles are denoted by $`E_{1+}^{3/2}`$, $`M_{1+}^{3/2}`$ and $`L_{1+}^{3/2}`$ (or $`S_{1+}^{3/2}`$), i.e. in the same way as the pion-nucleon phase. As a further example, the threshold production is determined by s-wave pions, i.e. $`l=0`$, $`J=\frac{1}{2}`$, which leads to $`E1`$ or $`S1`$ transitions and multipoles $`E_{0+}`$ or $`S_{0+}`$.
We complete the formalism of pion photoproduction by a discussion of isospin. Since the incoming photon has both isoscalar and isovector components and the produced pion is isovector, the matrix elements take the form $`^\mathrm{?}`$
$$_l^\alpha =\frac{1}{2}[\tau _\alpha ,\tau _0]_l^{()}+\frac{1}{2}\{\tau _\alpha ,\tau _0\}_l^{(+)}+\tau _\alpha _l^0.$$
(68)
The first two amplitudes on the $`rhs`$ can also be combined to
$$_l^{(\frac{3}{2})}=_l^{(+)}_l^{()},_l^{(\frac{1}{2})}=_l^{(+)}+2_l^{()},$$
(69)
where the upper index $`\frac{3}{2}`$ or $`\frac{1}{2}`$ denotes the isospin of the final state. The 4 physical amplitudes are then given in terms of linear combinations of the 3 isospin amplitudes. We note, however, that the isospin symmetry is broken by the mass differences between the nucleons $`(n,p)`$ and pions $`(\pi ^\pm ,\pi ^0)`$ and by explicit Coulomb effects, in particular near threshold.
### 6.1 Threshold pion photoproduction
As has been pointed out before, threshold production is dominated by the multipoles $`E_{0+}`$ ($`s`$-wave pions). For these multipoles there existed a venerable low energy theorem $`^\mathrm{?}`$, which however had to be revised in view of surprising experimental evidence.
Table 3 compares our predictions from dispersion theory to the “classical” low energy theorem (LET), ChPT and experiment. Note that ChPT $`^\mathrm{?}`$ contains the lowest order loop corrections, while “LET” is based on tree graphs only. Due to the coupling between the channels, the real part of $`E_{0+}(\gamma p\pi ^0p)`$ obtains large contributions from the imaginary parts of the higher multipoles via the dispersion integrals. Altogether these contributions nearly cancel the large contribution of the Born terms, which correspond to the result of pseudoscalar coupling, leading to a total threshold value $`^\mathrm{?}`$
$`\text{Re}E_{0+}^{\text{thr}}(p\pi ^0)`$ $`=`$ $`7.63+4.150.41+2.32+0.29+0.07=1.22,`$
$`\text{Re}E_{0+}^{\text{thr}}(n\pi ^0)`$ $`=`$ $`5.23+4.150.41+3.680.930.05=1.19,`$ (70)
where the individual contributions on the $`rhs`$ are, in that order, the Born term, $`M_{1+},E_{1+},E_{0+},M_1`$ and higher multipoles.
As we see from Table 3, the discrepancy between the “classical” LET and the experiment is very substantial in the case of $`\pi ^0`$ production on the proton. The reason for this was first explained in the framework of ChPT by pion-loop corrections. An expansion in the mass ratio $`\mu =m_\pi /M1/7`$ leads to the result $`^\mathrm{?}`$
$$E_{0+}(\pi ^0p)=\frac{eg_{\pi N}}{8\pi m_\pi }\left\{\mu \mu ^2\frac{3+\kappa _p}{2}\mu ^2\frac{M^2}{16f_\pi ^2}+\mathrm{}\right\},$$
(71)
where $`g_{\pi N}`$ is the pion-nucleon coupling constant and $`f_\pi 93`$ MeV the pion decay constant. We observe that the leading term is proportional to $`\mu `$, which suppresses this process relative to charged pion production. The leading terms of these expansions can be understood, to some degree, by simply relating the dipole moments in the respective pion-nucleon states. In particular the expansion for $`\gamma n\pi ^0n`$ starts at $`O(\mu ^2)`$, because both particles in the final state are neutral. The third term on the $`rhs`$ of Eq. (71) is the loop correction. Though formally of higher order in $`\mu `$, its numerical value is larger than the leading term!
While the threshold cross section receives its forward-backward asymmetry essentially from the combination $`\text{Re}\{E_{0^+}^{}(M_{1^+}+3E_{1^+}M_1^{})\}`$, the photon asymmetry $`\mathrm{\Sigma }`$ is dominated by $`\text{Re}\{M_{1^+}^{}(E_{1^+}+M_1^{})\}`$ and the target asymmetry $`T`$ by $`\text{Im}\{E_{0^+}^{}(E_{1^+}M_{1^+})\}`$. Since $`E_{1^+}`$ is small, the value of $`\mathrm{\Sigma }`$ is surprisingly sensitive to the multipole $`M_1^{}`$ resonating at the Roper resonance $`\text{N}^{}`$(1440). The observable T, on the other side, measures the phase of pion-nucleon s-wave scattering at threshold relative to the phase of the $`\mathrm{\Delta }`$ (1232) multipole.
Finally, the energy dependence of $`E_{0+}(\pi ^0p)`$ near threshold is shown in Fig. 10. The discrepancy between the “classical” LET and the experimental data is clearly seen, and one also observes a “Wigner cusp” at the $`\gamma p\pi ^+n`$ threshold. In particular, the imaginary amplitude rises sharply due to the strong coupling to this channel. Since charged pion production is much more likely to happen, neutral pions will often be produced by rescattering $`\gamma p\pi ^+n\pi ^0p`$.
### 6.2 Pion production in the resonance region
The search for a deformation of the “elementary” particles is a longstanding issue. Such a deformation is evidence for a strong tensor force between the constituents, originating in the case of the nucleon from the residual force of gluon exchange. Depending on one’s favourite model, such effects can be described by d-state admixture in the quark wave function $`^\mathrm{?}`$, tensor correlations between the pion cloud and the quark bag $`^{\mathrm{?},\mathrm{?}}`$, or by exchange currents accompanying the exchange of mesons between the quarks $`^\mathrm{?}`$. Unfortunately, it would require a target with a spin of at least 3/2 (e.g. $`\mathrm{\Delta }`$ matter) to observe a static deformation. An alternative is to measure the transition quadrupole moment between the nucleon and the $`\mathrm{\Delta }`$, i.e. the amplitude $`E_{1+}`$, which is sensitive to model parameters responsible for a possible deformation of the hadrons.
The experimental quantity of interest is the ratio $`R_{EM}=E_{1+}/M_{1+}`$ in the region of the $`\mathrm{\Delta }`$. The two amplitudes $`E_{1+}`$ and $`M_{1+}`$ are related to the helicity amplitudes, which may be determined by scattering an incident photon with circular polarization off a target nucleon with its spin oriented in the direction or opposite to the photon momentum $`\stackrel{}{q}`$,
$`A_1`$ $`=`$ $`A_{1/2}={\displaystyle \frac{1}{\sqrt{2q}}}N^{}(J,M={\displaystyle \frac{1}{2}})J_+N(J_i={\displaystyle \frac{1}{2}},M_i={\displaystyle \frac{1}{2}})`$
$`A_3`$ $`=`$ $`A_{3/2}={\displaystyle \frac{1}{\sqrt{2q}}}N^{}(J,M={\displaystyle \frac{3}{2}})J_+N(J_i={\displaystyle \frac{1}{2}},M_i=+{\displaystyle \frac{1}{2}}),`$ (72)
where $`J_+`$ is the hadronic current corresponding to the absorption of a photon with positive helicity on the nucleon $`N`$ with spin $`J_i=\frac{1}{2}`$ and spin projection $`M_i`$, leading to a resonance state $`N^{}`$ with spin $`J\frac{1}{2}`$ and spin projection $`M`$. It is obvious that all resonances can generally contribute to $`A_1`$, while only resonances with $`J\frac{3}{2}`$ will contribute to $`A_3`$. The helicity-conserving process $`A_1`$ can also occur on an individual (massless) quark, whereas $`A_3`$ is forbidden in that approximation. Hence perturbative QCD predicts that $`A_3`$ should vanish for high momentum transfer, i.e. for electroproduction and $`Q^2=\stackrel{}{q}^{2}\omega ^2\mathrm{}`$.
We shall now compare the prediction of Kälbermann and Eisenberg $`^\mathrm{?}`$ with our analysis of the modern pion photoproduction data $`^{\mathrm{?},\mathrm{?}}`$. The helicity amplitudes for the $`N\mathrm{\Delta }`$ transition will be given in the usual units of 10<sup>-3</sup> GeV$`^{\frac{1}{2}}`$, and the prediction of Ref. $`^\mathrm{?}`$ is obtained for a bag radius of 1 fm, which was the preferred value in the 1980’s. The result is
$$\begin{array}{cc}\mathrm{Ref}.{}_{}{}^{74}:A_1=130\hfill & A_3=250\hfill \\ \mathrm{Ref}.{}_{}{}^{58}:A_1=131\pm 1\hfill & A_3=252\pm 1,\hfill \end{array}$$
the agreement being truly astounding though somewhat accidental, because the theoretical value depends on the bag radius. However, the result is relatively stable, even a drastic decrease of the bag radius to 0.6 fm will change the helicity amplitudes by only 10%. This success of the chiral bag model is even more outstanding when compared with the results of the quark model without pionic degress of freedom. From a selection of ten quark model calculations published over the past 20 years we find $`113A_182`$ and $`195A_358`$, values far off the experimental data.
The helicity amplitudes are related to the electric and magnetic multipoles,
$`M_{1^+}`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{3}}}(\sqrt{3}A_{1/2}+3A_{3/2}),`$
$`E_{1^+}`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{3}}}(\sqrt{3}A_{1/2}A_{3/2}).`$ (73)
Since $`A_{\frac{3}{2}}\sqrt{3}A_{\frac{1}{2}}`$ according to Eq. (6.2), the model predicts that $`E_{1^+}`$ (electric quadrupole excitation $`E2`$) is very much smaller than $`M_{1^+}`$ (magnetic dipole excitation $`M1`$). A few years after the pioneering work of Ref. $`^\mathrm{?}`$, we obtained, for a bag radius of 0.6 fm, the ratio $`R=E_{1^+}/M_{1^+}=2.8\%`$ $`^\mathrm{?}`$. This result differed by a factor of two from the then accepted experimental value $`R_{1988}1.3\%`$. However, it is quite close to the recent MAMI data of Beck et al. $`^\mathrm{?}`$, $`R_{1997}=(2.5\pm 0.2\pm 0.2)\%`$, and to our global analysis of the data $`^\mathrm{?}`$, $`R_{1998}=(2.5\pm 0.1)\%`$.
As may be seen from Fig. 11 , the ratio $`R=R_{EM}`$ changes rapidly with the energy $`W_{cm}`$ of the pion-nucleon system. The reason for this energy dependence is the nonresonant background, which is particularly large in the case of the small $`E_{1^+}`$ multipole. The historically first prediction of that number is due to Chew et al. $`^\mathrm{?}`$ in 1957 who found $`R0`$ from a dispersion theoretical analysis. Such value was later explained by Becchi and Morpurgo $`^\mathrm{?}`$ in the framework of the constituent quark model. In the following years the quark models were refined by introducing tensor correlations, with the result of finite, small and usually negative values for $`R`$. Such correlations have been motivated in different ways, by hyperfine interactions between the quarks $`^\mathrm{?}`$, pion-loop effects $`^\mathrm{?}`$ and, more recently, exchange currents $`^\mathrm{?}`$. In analogy with heavy even-even nuclei having “intrinsic” deformation, finite values of $`E2`$ are often referred to as ”bag deformation” or ”deformation of the nucleon”, although a quadrupole moment cannot be observed for an object with spin $`J<1`$. Ideally one could probe the static quadrupole moment of the $`\mathrm{\Delta }`$ by experiments like $`\pi N\mathrm{\Delta }\mathrm{\Delta }\gamma \pi N\gamma `$, however a closer look shows that this is hardly a realistic possibility. In conclusion it is precisely the $`N\mathrm{\Delta }`$ transition quadrupole moment that provides us with information on tensor correlations in the nucleon, which can be translated, e.g., into a $`d`$-state admixture in the quark wave function. In such a model the $`\mathrm{\Delta }`$ would have an oblate deformation, much smaller than a frisbee and much larger than the earth, in absolute numbers quite comparable to the deuteron, which however has a prolate deformation.
Meson electroproduction allows us to study the dependence of the multipoles on momentum transfer, $`M_{l\pm }=M_{l\pm }(Q^2)`$, i.e. to probe the spatial distribution of the transition strength. In addition, the virtual photon carries a longitudinal field introducing a further multipole, $`S_{l\pm }`$. The $`Q^2`$ dependence of the $`N\mathrm{\Delta }`$ multipoles is displayed in Fig. 12. In the top figure, we show the results for the magnetic multipole divided by the standard dipole form factor. The data are compared to the predictions of our unitary isobar model (UIM) $`^\mathrm{?}`$. This model contains the usual Born terms, vector meson exchange in the t-channel and nucleon resonances in the s-channel, unitarized partial wave by partial wave with the appropriate pion-nucleon phases and inelasticities. The center piece of Fig. 12 shows the ratio $`R=R(Q^2)`$ compared to mostly older and strongly fluctuating data. More recent data from Jefferson Lab $`^\mathrm{?}`$ indicate, however, that even at $`Q^22.8`$ and 4 (GeV/c)<sup>2</sup> this ratio remains negative and of the order of a few per cent. This is surprising, because perturbative QCD predicts that the helicity amplitude $`A_3`$ should vanish for $`Q^2\mathrm{}`$ and, hence, the ratio $`R`$ should approach +100% (see Eq. (73)). Finally, the bottom figure shows the corresponding longitudinal-transverse ratio $`S_{1^+}/M_{1^+}`$. Recent experimental data at ELSA, MAMI and MIT $`^\mathrm{?}`$ at $`Q^20.5`$ (GeV/c$`)^2`$ yield ratios of about -7 %, slightly below our prediction, while the preliminary data from the Jefferson Lab $`^\mathrm{?}`$ at larger $`Q^2`$ seem to indicate considerably lower values between -10 % and -20 %. From perturbative QCD one expects that this ratio should vanish for $`Q^2\mathrm{}`$.
The modern precision experiments will be continued to the higher resonances. Concerning the $`N^{}`$(1440) or Roper resonance, both data and predictions are still in a deploratory state, and it will require double-polarization experiments to find out about the nature of that resonance. One possibility to tackle the problem will be pion production by linearly polarized photons on longitudinally polarized protons. Such an experiment measures the polarization observable $`GImM_1^{}ReM_{1^+}`$, i.e. an interference of the $`\mathrm{\Delta }`$ resonance with the absorptive part of the Roper multipole $`M_1^{}`$.
The existing information on some of the higher resonances is shown in Fig. 13. For our discussion in the next chapter it is important to note:
1. Because of its quantum numbers $`J^P=\frac{1}{2}^{}`$, the resonance $`S_{11}`$(1535) is only excited via the $`A_{1/2}`$ amplitude. As function of $`Q^2`$, this amplitude drops much slower than any other resonance of the nucleon. With a resonance position very close to $`\eta `$ production threshold, the $`S_{11}`$(1535) has an $`\eta `$ branching ratio of about 50 %, while this ratio is of the order of 1 % or less for all other resonances.
2. The resonances $`D_{13}`$(1520) and $`F_{15}`$(1680) carry most of the electric dipole and quadrupole strengths, respectively. For real photons $`(Q^2=0)`$ their helicity amplitudes $`A_{1/2}^p`$ are nearly zero, but already at $`Q^20.5`$ (GeV/c)<sup>2</sup> $`A_{1/2}^p`$ and $`A_{3/2}^p`$ are of equal importance, and in accordance with pQCD, $`A_{3/2}^p`$ decreases rapidly for $`Q^2\mathrm{}`$.
## 7 SUM RULES
As has been stated in Eq. (57), the GDH sum rule connects the integral
$$I=_{\nu _0}^{\mathrm{}}\frac{\sigma _{1/2}(\nu )\sigma _{3/2}(\nu )}{\nu }𝑑\nu $$
(74)
with the anomalous magnetic moment.
On the basis of the pion-nucleon multipoles and certain assumptions for the higher channels, various authors have estimated this integral. As shown in Table 4, the absolute value of the proton integral $`I_p`$ has been consistently overpredicted, while the neutron integral $`I_n`$ comes out too small. This has the consequence that not even the sign of the isovector combination $`I_pI_n`$ agrees with the sum rule value. This apparent discrepancy has led to speculations that the GDH integral should not converge for various reasons, e.g. due to a generalized current algebra, because of fixed axial vector poles or influences of the Higgs particle. None of these arguments is too convincing at present. In fact one should realize that the GDH integrand is an oscillating function of photon energy, with multipole contributions of alternating sign. Therefore, little details matter and a stable result requires very exact data. Comparing again the results obtained with the SAID $`^\mathrm{?}`$ and HDT $`^\mathrm{?}`$ multipoles, the generally accepted threshold value of $`E_{0^+}`$ reduces the “discrepancy” with the sum rule value by about $`25\%`$ (see Table 4 and Ref. $`^\mathrm{?}`$).
The first direct measurement of the helicity cross sections was recently performed at MAMI in the energy region 200 MeV $`<\nu <800`$ MeV $`^\mathrm{?}`$. The experiment will be extended to the higher energies at ELSA. Some preliminary results are shown in Fig. 14, which contains only 5 % of the data taken in the 1998 run. The figure shows the importance of charged pion production near threshold (multipole $`E_{0+}`$), and the dominance of the multipole $`M_{1+}^{3/2}`$ in the $`\mathrm{\Delta }`$ resonance region. At yet higher energies the data lie above the prediction for one-pion production, which indicates considerable two-pion contributions. These data establish that the forward spin polarizability should be $`\gamma _00.810^4`$ fm<sup>4</sup>. Furthermore the preliminary data saturate the GDH sum rule at $`\nu 800`$ MeV if one accepts our predictions $`^\mathrm{?}`$ for the energy range between threshold and 200 MeV. However, more data are urgently required at energies both below 200 MeV and above 800 MeV.
In view of the difficulty to obtain even the proper sign for the proton-neutron difference from the older data (see Table 4), it is of considerable interest to measure the GDH for the neutron. However, such investigations are difficult due to nuclear binding effects. While it is generally assumed that <sup>2</sup>H and <sup>3</sup>He are good neutron targets, the sum rule requires to integrate over all regions of phase space and not only the region of quasifree kinematics. In fact there exists a GDH sum rule for systems of any spin, and hence every nucleus should have a well-defined value for the GDH integral. With the definition of Table 4, one finds the small value $`I(^2`$H$`)=0.65\mu `$b due to the fact that the deutron lies very close to the Schmidt line. However, a loosely bound system of neutron and proton would be expected to have $`I_p+I_n=438\mu `$b, which differs by 3 orders of magnitude from the deuteron value! Obviously the large contributions from pion production have to be canceled by binding effects in the deuteron. As has been shown by Arenhövel and collaborators $`^\mathrm{?}`$, such contribution is mainly due to the transition from the <sup>3</sup>S<sub>1</sub> ground state of <sup>2</sup>H to the <sup>1</sup>S<sub>0</sub> resonance at 68 keV. Weighted with the inverse power of excitation energy, the absorption cross section for this low-lying resonance cancels the huge cross sections due to pion production. We note that the opposite sign of the two contributions is due to the fact that the spins of the 3 quarks become aligned by the transition $`N\mathrm{\Delta }`$, while the nucleon spins are parallel in the deuteron $`(^3`$S<sub>1</sub>) but antiparallel in the <sup>1</sup>S resonance. However, in addition to this low lying resonance, there are also sizeable sum rule contributions by break-up reactions $`\gamma +dn+p`$ in the range below and above pion threshold. Such effects are usually triggered by meson exchange currents (the virtual pions below or the real pions above threshold are reabsorbed by the other nucleon) or isobaric currents (a $`\mathrm{\Delta }`$ is produced but decays by final state interaction without emitting a pion). In addition there are also contributions of coherent $`\pi ^0`$ production, i.e. $`\gamma +d\pi ^0+d`$. It is general to all these processes that they cannot occur on a free nucleon, though they are certainly driven by pion production and nucleon resonances. This leads to the serious question: Which part of the GDH integral is nuclear structure and hence should be subtracted, and how should one divide the rest into the contributions of protons and neutrons? The problem is not restricted to the deuteron but quite general. For example the “neutron target” <sup>3</sup>He has the same sum rule as the nucleon except that one has to replace $`e,m`$ and $`\kappa `$ of the nucleon by the charge $`(Q=2e)`$, mass $`(M3m_N)`$ and anomalous magnetic moment of <sup>3</sup>He. The result is $`I(^3`$He$`)=496\mu `$b, while we naively expect $`I_n=233\mu `$b, because the spins of the two protons are antiparallel and hence should not contribute to the helicity asymmetry.
These considerations can be generalized to virtual photons by electron scattering. While the coincidence cross section for the reaction $`\stackrel{}{e}+\stackrel{}{p}e^{}+N+\pi `$ contains 18 different response functions $`^\mathrm{?}`$, only 4 responses remain after integration over the angles of pion emission, which is exactly the result of Eq. (17). The 4 partial cross sections can in principle be separated by a super-Rosenbluth plot if one varies the polarizations. These are the transverse polarizations $`\epsilon `$ of the virtual photon, the polarization $`P_e`$ of the electron ($`\pm 1`$ for the relativistic case), and the nucleon’s polarization in the scattering plane of the electron, with components $`P_z`$ in the direction and $`P_x`$ perpendicular to the virtual photon momentum.
The multipole content of one-pion production to the partial cross sections is $`^{\mathrm{?},\mathrm{?}}`$
$`\sigma _T^{(1\pi )}`$ $`=`$ $`4\pi {\displaystyle \frac{|𝐤_\pi ^{cm}|}{k^{cm}}}{\displaystyle \underset{l}{}}{\displaystyle \frac{1}{2}}(l+1)^2`$
$`[(l+2)(|E_{l+}|^2+|M_{l+1,}|^2)+l(|M_{l+}|^2+|E_{l+1,}|^2)]`$
$`=`$ $`4\pi {\displaystyle \frac{|𝐤_\pi ^{cm}|}{k^{cm}}}\{|E_{0+}|^2+2|M_{1+}|^2+6|E_{1+}|^2+|M_1|^2+2|E_2|^2\pm \mathrm{}\},`$
$`\sigma _L^{(1\pi )}`$ $`=`$ $`4\pi {\displaystyle \frac{|𝐤_\pi ^{cm}|}{k^{cm}}}\left({\displaystyle \frac{Q}{\omega ^{cm}}}\right)^2{\displaystyle \underset{l}{}}{\displaystyle \frac{1}{2}}(l+1)^3[|L_{1+}|^2+|L_{l+1,}|^2]`$
$`=`$ $`4\pi {\displaystyle \frac{|𝐤_\pi ^{cm}|}{k^{cm}}}\left({\displaystyle \frac{Q}{\omega ^{cm}}}\right)^2\left\{|L_{0+}|^2+8|L_{1+}|^2+|L_1|^2+8|L_2|^2\pm \mathrm{}\right\},`$
$`\sigma _{TT^{}}^{(1\pi )}`$ $`=`$ $`4\pi {\displaystyle \frac{|𝐤_\pi ^{cm}|}{k^{cm}}}{\displaystyle \underset{l}{}}{\displaystyle \frac{1}{2}}(l+1)[(l+2)(|E_{l+}|^2+|M_{l+1,}|^2)`$
$`+l(|M_{l+}|^2+|E_{l+1,}|^2)2l(l+2)(E_{l+}^{}M_{l+}E_{l+1,}^{}M_{l+1,})]`$
$`=`$ $`4\pi {\displaystyle \frac{|𝐤_\pi ^{cm}|}{k^{cm}}}\{|E_{0+}|^2+|M_{1+}|^26E_{1+}^{}M_{1+}3|E_{1+}|^2+|E_2|^2\pm \mathrm{}\}`$
$`\sigma _{LT^{}}^{(1\pi )}`$ $`=`$ $`4\pi {\displaystyle \frac{|𝐤_\pi ^{cm}|}{k^{cm}}}\left({\displaystyle \frac{Q}{\omega ^{cm}}}\right){\displaystyle \underset{l}{}}{\displaystyle \frac{1}{2}}(l+1)^2`$
$`[L_{l+}^{}((l+2)E_{l+}+lM_{l+})+L_{l+1,}^{}(lE_{1+,}+(l+2)M_{l+1,})]`$
$`=`$ $`4\pi {\displaystyle \frac{|𝐤_\pi ^{cm}|}{k^{cm}}}\left({\displaystyle \frac{Q}{\omega ^{cm}}}\right)\{L_{0+}^{}E_{0+}2L_{l+}^{}(M_{1+}+3E_{1+})`$
$`+L_1^{}M_1+L_2^{}E_2\pm \mathrm{}\}.`$
Since the partial wave decomposition is defined in the hadronic $`cm`$ frame, the appropriate $`cm`$ values of the kinematical observables have to be used in these equations, in particular the $`cm`$ momentum and the $`cm`$ energy of the virtual photon, $`k^{cm}=\frac{m}{W}k`$ and $`\omega ^{cm}=\frac{1}{W}\sqrt{m^2\nu ^2Q^2(W^2m^2)}`$ respectively. We note that $`\omega ^{cm}`$ has a zero if $`W=\sqrt{m^2+Q^2}`$, which is compensated by a corresponding zero in the longitudinal multipole. This situation can be avoided by using the “scalar” multipoles (rather to be called “Coulomb” or “time-like” multipoles!),
$$S_{l\pm }=\frac{k^{cm}}{\omega ^{cm}}L_{l\pm }.$$
(79)
While $`\sigma _T`$ and $`\sigma _L`$ are the sum of squares of transverse $`(E_{l\pm },M_{l\pm })`$ and longitudinal $`(L_{l\pm })`$ multipoles respectively, the interference structure functions $`\sigma _{TT}^{}=(\sigma _{3/2}\sigma _{1/2})/2`$ and $`\sigma _{LT}^{}`$ contain multipole contributions of alternating sign. The multipoles involved are now functions of energy and momentum transfer, $`_{l\pm }=_{l\pm }(\nu ,Q^2)`$.
The 4 cross sections are related to the familiar structure functions of deep inelastic lepton scattering,
$$\{\sigma _T,\sigma _L;\sigma _{LT}^{},\sigma _{TT}^{}\}\{F_1,F_2;G_1,G_2\}.$$
(80)
In the Bjorken scaling region, the 2 arguments $`\nu `$ and $`Q^2`$ of these functions can be replaced by the scaling variable $`x=Q^2/2m\nu `$, which leads to the definition of quark distribution functions. For the spin structure functions $`G_1`$ and $`G_2`$ we find
$`g_1(\nu ,Q^2)`$ $`=`$ $`{\displaystyle \frac{\nu }{m}}G_1(\nu ,Q^2)g_1(x)={\displaystyle \frac{1}{2}}{\displaystyle e_i^2(f_i^{}f_i^{})}`$
$`g_2(\nu ,Q^2)`$ $`=`$ $`{\displaystyle \frac{\nu ^2}{m^2}}G_2(\nu ,Q^2)g_2(x)={\displaystyle \frac{1}{2}}{\displaystyle e_i^2(f_i^{}f_i^{})},`$ (81)
where the arrows indicate the different directions of the quark spins. With these definitions we can express a set of generalized sum rules in terms of both the quark spin functions of Eq. (7) and the cross sections of Eq. (79), e.g.
$`I_1(Q^2)`$ $`=`$ $`{\displaystyle _{\nu _0}^{\mathrm{}}}{\displaystyle \frac{d\nu }{\nu }}G_1(\nu ,Q^2)={\displaystyle \frac{2m}{Q^2}}{\displaystyle _0^{x_0}}𝑑xg_1(x,Q^2){\displaystyle \frac{2m^2}{Q^2}}\mathrm{\Gamma }`$ (82)
$`=`$ $`{\displaystyle \frac{m^2}{2\pi e^2}}{\displaystyle _{\nu _0}^{\mathrm{}}}𝑑\nu (1x)\left(\sigma _{1/2}\sigma _{3/2}+2{\displaystyle \frac{Q}{\nu }}\sigma _{LT}^{}\right),`$
$`I_2(Q^2)`$ $`=`$ $`{\displaystyle \frac{1}{m}}{\displaystyle _{\nu _0}^{\mathrm{}}}𝑑\nu G_2(\nu ,Q^2)={\displaystyle \frac{2m^2}{Q^2}}{\displaystyle _0^{x_0}}𝑑xg_2(x,Q^2)`$ (83)
$`=`$ $`{\displaystyle \frac{m^2}{2\pi e^2}}{\displaystyle _{\nu _0}^{\mathrm{}}}𝑑\nu (1x)\left(\sigma _{1/2}+\sigma _{3/2}+2{\displaystyle \frac{\nu }{Q}}\sigma _{LT}^{}\right),`$
with $`\nu _0`$ and $`x_0`$ the lowest threshold for inelastic reactions.
Eq. (82) is a possible generalization of the GDH sum rule, because $`I_1(0)=\kappa ^2/4`$. However, a large variety of generalized GDH sum rules can be obtained by adding different fractions of the interference term $`\frac{Q}{\nu }\sigma _{LT}^{}`$, which vanishes both in the real photon limit, $`Q^20`$, and in the asymptotic region, $`Q^2\mathrm{}`$. The most obvious choice would be to simply drop this term in Eq. (82). As it stands, however, the definition of $`I_1`$ is the natural definition of an integral over the spin structure function $`g_1`$. In particular it has the asymptotic behaviour indicated in Eq. (82), with $`\mathrm{\Gamma }`$ a constant. The fact that the experimental value of $`\mathrm{\Gamma }`$ differed from earlier predictions $`^\mathrm{?}`$ led to the “spin crisis” and taught us that less than half of the nucleon’s spin is carried by the quarks $`^\mathrm{?}`$.
The integral Eq. (83) for the second spin structure functions $`G_2`$ shows distinct differences in comparison with Eq. (82). First, the helicity cross sections $`\sigma _{1/2}`$ and $`\sigma _{3/2}`$ appear with different sign. This has the consequence that in the sum $`I_{1+2}=I_1+I_2`$ only the longitudinal-transverse contribution $`\sigma _{LT}^{}`$ remains. Second, the latter contribution now appears as $`\frac{\nu }{Q}\sigma _{LT}^{}`$, which is finite in the real-photon limit, $`Q^20`$. While the generalized GDH integral of Eq. (82) ist not a sum rule, i.e. not related to another observable except for the real photon point, the Burkhardt-Cottingham (BC) sum rule predicts that $`I_2(Q^2)`$ can be expressed by the magnetic $`(G_M)`$ and electric $`(G_E)`$ Sachs form factors at each momentum transfer $`^\mathrm{?}`$,
$$I_2(Q^2)=\frac{1}{4}G_M(Q^2)\frac{G_M(Q^2)G_E(Q^2)}{1+Q^2/4m^2}.$$
(84)
According to Eq. (84) the integral $`I_2`$ approaches the value $`\kappa \mu /4`$ for real photons $`(Q^2=0)`$ and drops with $`Q^{10}`$ for $`Q^2\mathrm{}`$. As a result the sum $`I_{1+2}(0)`$ should take the value $`\kappa (\mu 1)/4`$, i.e. $`\kappa ^2/4`$ and 0 for proton and neutron respectively. However, there are strong indications that the BC integrand gets large contributions at higher energies, which in fact will affect its convergence. At least for the proton, however, the “sum rule” seems to work quite well if we restrict ourselves to the resonance region.
In the case of the proton, the GDH sum rule predicts $`\mathrm{\Gamma }_1<0`$ for small $`Q^2`$, while all experiments for $`Q^2>1`$ (GeV/c)<sup>2</sup> yield positive values. Clearly, the value of the sum rule has to change rapidly at low $`Q^2`$, with some zero-crossing at $`Q_0^2<1`$ (GeV/c)<sup>2</sup>. This evolution of the sum rule was first parametrized by Anselmino et al. $`^\mathrm{?}`$ in terms of vector meson dominance. Burkert, Ioffe and others $`^\mathrm{?}`$ refined this model considerably by treating the contributions of the resonances explicitly. Soffer and Teryaev $`^\mathrm{?}`$ suggested that the rapid fluctuation of $`I_1`$ should be analyzed in conjunction with $`I_2`$, because $`I_1+I_2`$ is known for both $`Q^2=0`$ and $`Q^2\mathrm{}`$. Though this sum is related to the practically unknown longitudinal-transverse interference cross section $`\sigma _{LT}^{}`$ and therefore not yet determined directly, it can be extrapolated smoothly between the two limiting values of $`Q^2`$. The rapid fluctuation of $`I_1`$ then follows by subtraction of the BC value of $`I_2`$. We also refer the reader to a recent evaluation of the $`Q^2`$-dependence of the GDH sum rule in a constituent quark model $`^\mathrm{?}`$, and to a discussion of the constraints provided by chiral perturbation theory at low $`Q^2`$ and twist-expansions at high $`Q^2`$ (see Ref. $`^\mathrm{?}`$).
In Fig. 15 we give our predictions $`^\mathrm{?}`$ for the integrals $`I_1(Q^2)`$ and $`I_2(Q^2)`$ in the resonance region, i.e. integrated up to $`W_{max}=2GeV`$. As can be seen, our model is able to generate the dramatic change in the helicity structure quite well. While this effect is basically due to the single-pion component predicted by the UIM, the eta and multipion channels are quite essential to shift the zero-crossing of $`I_1`$ from $`Q^2=0.75(GeV/c)^2`$ to 0.52 $`(GeV/c)^2`$ and 0.45 $`(GeV/c)^2`$, respectively. This improves the agreement with the SLAC data $`^\mathrm{?}`$. However, some differences remain. Due to a lack of data in the $`\mathrm{\Delta }`$ region, the SLAC data are likely to underestimate the $`\mathrm{\Delta }`$ contribution and thus to overestimate the $`I_1`$ integral or the corresponding first moment $`\mathrm{\Gamma }_1`$. A few more data points in the $`\mathrm{\Delta }`$ region would be very useful in order to clarify the situation, and we are looking forward to the results of the current experiments at Jefferson Lab $`^\mathrm{?}`$.
Concerning the integral $`I_2`$, our results are in good agreement with the predictions of the BC sum rule (see Fig. 15, lower part). The remaining differences are of the order of $`10\%`$ and should be attributed to contributions beyond $`W_{\mathrm{max}}=2`$ GeV and the scarce experimental data for $`\sigma _{LT}^{}`$.
We recall at this point that the results mentioned above refer to the proton. Unfortunately, we find some serious problems for the neutron, for which our model predicts both $`I_1(0)`$ and $`I_2(0)`$ larger than expected from the sum rules. This has the consequence that our prediction for $`I_{1+2}(0)`$ has a relatively large positive value while it should vanish by sum rule arguments. The reason for this striking discrepancy could well be due to the discussed problems with “neutron targets”. On the other hand it could also be an indication of sizeable contributions at the higher energies, which could possibly cancel for the proton but add in the case of the neutron. In this context it is interesting to note that a recent parametrization of deep inelastic scattering predicts sizeable high-energy contributions with different signs for proton and neutron $`^\mathrm{?}`$.
A more general argument is that the convergence of sum rules cannot be given for granted, and thus the good agreement of our model with the BC sum rule could be accidental and due to a particular model prediction for the essentially unknown longitudinal-transverse interference term. As can be seen from Fig. 15, the contribution of $`\sigma _{LT^{}}`$ is quite substantial for $`I_2`$ even at the real photon point due to the factor $`\nu /Q`$ in Eq. (83). This contribution, however, is constrained by the positivity relation $`|\sigma _{LT}^{}|\sqrt{\sigma _L\sigma _T}`$. The dash-dotted line shows the integral for the upper limit of this inequality and a similar effect would occur for the lower limit. The surprisingly large upper limit can be understood in terms of multipoles. In a realistic description of the integrated cross section $`\sigma _{LT^{}}`$, the large $`M_{1+}`$ multipole can only interfer with the small $`L_{1+}`$ multipole. The upper and lower limits of the positivity relation overestimate the structure function considerably due to an unphysical “interference” between $`s`$ and $`p`$ waves.
## 8 SUMMARY
The new generation of electron accelerators with high energy, intensity and duty-factor has made it possible to perform new classes of coincidence experiments involving polarization observables. These investigations have already provided new data with unprecedented precision, and they will continue to do so for the years to come. Some of the interesting topics and challenging questions are
* a full separation of the electric and magnetic form factors of neutron and proton by double-polarization experiments,
* the search for strange quarks in the nucleon by parity-violating electron scattering,
* new and more precise information on the scalar and vector polarizabilities of the nucleon by a combined analysis of Compton scattering and photoproduction as well as extensions to generalized polarizabilities via virtual Compton scattering,
* the threshold amplitudes for the production of Goldstone bosons and tests of chiral field theories,
* the quadrupole strength for $`\mathrm{\Delta }_{1232}`$ excitation as a measure of tensor correlations among the constituents,
* photo- and electroexcitation of the higher resonances, e.g. the $`N_{1440}^{}`$ (Why does the Roper occur at such a low excitation energy? Where is its Coulomb monopole strength?), the $`N_{1535}^{}`$ (Is it really a resonance or a threshold effect of $`\eta `$ production?), and the helicity structure of the main dipole $`(N_{1520}^{})`$ and quadrupole ($`N_{1680}^{})`$ resonances for both proton and neutron,
* investigations of individual decay channels including energy and angular distributions in order to find out how much of the excitation strength is actually due to resonances as opposed to background and threshold effects, and more generally the question how to extract the “intrinsic” quark structure from the experimental data, which necessarily contain the hadronization in terms of mesons,
* ongoing experiments to determine the helicity structure of photo- and electroproduction in the resonance region by use of double-polarization observables, which in turn are related to deep inelastic scattering and the quark spin structure by means of sum rules and related integrals over the excitation spectrum.
Our present understanding of nonperturbative QCD is still in a somewhat deplorable phenomenological state. The ongoing experimental activities will change that situation within short by providing new and detailed information on low-energy QCD in general and the nucleon’s structure in particular. This rich phenomenology will without doubt challenge the theoretical approaches and, as is strongly to be hoped, eventually pave the way for a more quantitative understanding of nonperturbative QCD.
## References
TESTS AND PROBLEMS
1. KINEMATICS
2. Prove that $`s+t+u=m_1^2+m_2^2+m_3^2+m_4^2`$, for the reaction $`p_1+p_2p_3+p_4`$.
3. Calculate the threshold energy in the lab frame for the reactions
* $`p(\gamma ,\gamma ^{})p^{}`$
* $`p(\gamma ,\pi )p^{}`$
* replace the incident real photon by a virtual one with
$`m_\gamma ^{}^2=Q^2<0`$
4. Which energy should an accelerator have to electroproduce a $`K^+`$ at $`Q^2=1`$ (GeV/c)<sup>2</sup>? How much energy would you like to have before you schedule such an experiment?
5. In case of the reaction $`p_1+p_2p_3+p_4`$ with 4 scalar particles, how many Lorentz scalars and vectors can be constructed? How many are independent?
6. Find the kinematical limits for Compton scattering in the s-channel. Use $`\nu =(su)/4M`$ and $`t`$ as orthogonal coordinates, and relate them in the $`cm`$ frame for forward and backward scattering. What about the other parts of the hyperbola of Fig. 3?
7. FORM FACTORS
8. Evaluate the vector current $`J_u=\overline{u}_p^{}(F_1\gamma _\mu +\frac{i\sigma _{\mu \nu }q^\nu }{2m}F_2)u_p`$ in the Breit frame, $`\stackrel{}{p}=\frac{1}{2}\stackrel{}{q}`$ and $`\stackrel{}{p}^{}=+\frac{1}{2}\stackrel{}{q}`$, and identify the Sachs form factors $`G_E`$ and $`G_M`$ of Eq. (23).
9. The electric Sachs form factor of the proton can be approximated by the dipole form $`G_E^p=(1+(Q/.84`$ GeV)$`{}_{}{}^{2})^2=G_D(Q^2)`$. Calculate the charge distribution by a Fourier transform according to Eq. (25).
10. STRANGENESS
11. Derive the structure of the Lorentz tensor $`W_{\mu \nu }=J_\mu J_\nu ^{}`$, constructed from a general vector current $`J_\mu `$ for an unpolarized fermion. Use independent 4-momenta $`q_\mu =p_{2\mu }p_{1\mu }`$ and $`P_\mu =\frac{1}{2}(p_{1\mu }+p_{2\mu })`$, and impose current conservation.
12. In the case of $`\stackrel{}{e}+pe^{}+p^{}`$, there also appears an antisymmetric tensor $`\eta _{\mu \nu }^{(𝒜)}=\frac{i}{2m^2}ϵ_{\mu \nu \alpha \beta }q^\alpha P^\beta `$ with a $`\pm `$ sign in front depending on the helicity of the electron. Why does this term not contribute to the cross section derived from photon exchange? What would be required to see this term?
13. A simple model of the proton says that part of the time it appears as a neutron surrounded by a $`\pi ^+`$ cloud. If this system has orbital momentum $`l=l_z=1`$, with which probability should it occur in order to describe the anomalous magnetic moment $`\kappa _p`$ of the proton?
14. Estimate the contribution of $`\mathrm{\Lambda }^{}K^+`$ configurations to the strangeness form factors of Eq. (32) following the procedure of No 10. Which sign does the strangeness magnetic moment $`\mu ^s`$ carry according to the model? What about the strangeness radius $`<r^2>_s`$?
15. COMPTON SCATTERING
16. Derive an upper limit for the electric polarizability of proton and neutron in a two-body model as in No 10. Use the definition of Eq. (45) and note that the excitation spectrum lies at $`E_nE_0>m_\pi `$. Express the result in terms of the radius $`<r^2>_E^{n,p}`$. How general is the result?
17. A generic model of a polarizable system is the following (see Ref. $`^\mathrm{?}`$): Two objects with masses $`M_{1,2}`$ and charges $`Q_{1,2}`$ are bound by a spring (oscillator frequency $`\omega _0^2=C/\mu `$, $`C=`$ Hooke’s constant, $`\mu =`$ reduced mass). An electric $`\stackrel{}{E}=\stackrel{}{E}_0e^{i\omega t}`$ induces a dipole moment $`\stackrel{}{D}=\alpha (\omega )\stackrel{}{E}_0`$.
* Determine $`\alpha (\omega )`$ and consider the cases of (I) equal particles, (II) $`M_2\mathrm{}`$, $`Q_20`$.
* If $`Q=Q_1+Q_20`$, the system will be accelerated even in the limit $`\omega 0`$. Calculate $`\stackrel{}{D}^{}`$ for the $`cm`$ coordinate.
* Calculate $`\stackrel{}{D}^{}`$ for the relative coordinate.
* Classical antenna theory says that the cross section is
$$\frac{d\sigma }{d\mathrm{\Omega }}=|f(\omega )|^2\left(\frac{|\stackrel{}{D}^{}|}{|\stackrel{}{E}|}\right)^2.$$
Compare this result with Eq. (48) and discuss the scattering amplitude $`f`$ and the cross section for $`\omega =0`$, $`\omega \omega _0`$, $`\omega \omega _0`$ and $`\omega \omega _0`$. Which kind of scattering occurs for $`Q=0`$?
18. Estimate the polarizability for the following systems, approximated by 2-body configurations, using the result of No 13 (see Ref. $`^\mathrm{?}`$).
* H atom = $`p+e^{}`$ , $`\mathrm{}\omega _010`$ eV
* deuteron = $`p+n`$ , $`\mathrm{}\omega _04.5`$ eV
* $`{}_{}{}^{208}Pb=82p+126n`$ , $`\mathrm{}\omega _014`$ eV
* $`p=2u+1d`$ , $`\mathrm{}\omega _0500`$ eV
* $`n=1u+2d`$ , $`\mathrm{}\omega _0500`$ eV
19. PION PHOTOPRODUCTION
20. Threshold pion production is given by the s-wave multipoles $`E_{0+}`$. Estimate these multipoles for the 4 physical channels $`\gamma +𝒩\pi +𝒩`$ by evaluating the squares of the electric dipole moments of the $`\pi 𝒩`$ configuration. Compare these results with Table 3.
21. Which multipoles connect the $`𝒩`$ with the following resonances:
$`P_{33}(1232),J^p=\frac{3}{2}^+`$; $`P_{11}(1440),J^p=\frac{1}{2}^+;D_{13}(1520),J^p=\frac{3}{2}^{};S_{11}(1535)`$, $`J^p=\frac{1}{2}^{};F_{15}(1680),J^p=\frac{5}{2}^+`$.
See Eq. (67) and the text following that equation.
22. SUM RULES
23. The integrand of the GDH sum rule has the multipole decomposition of Eq. (76). Draw a figure of the GDH integrand for the 4 physical channels as function of $`\omega `$, using the result of No. 16 and the information given in Section 6.
24. A possible generalization $`I_{GDH}`$ of the GDH integral is obtained from Eq. (83) by dropping the term in $`\sigma _{LT}^{}`$. Discuss the sign of $`I_{GDH}`$ for $`Q^20`$ and $`Q^2\mathrm{}`$. Give qualitative arguments why the sign change takes place already at relatively small $`Q^2`$.
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# References
In this paper we recalculate the two-loop beta-function coefficients in the two-dimensional Sine-Gordon (SG) model in a two-parameter perturbative expansion around the asymptotically free (AF) point. The study of the SG model in the vicinity of this point is especially important since this region is used in the description of the Kosterlitz-Thouless (KT) phase transition in the two-dimensional $`O(2)`$ nonlinear $`\sigma `$-model, better known as the XY model<sup>1</sup><sup>1</sup>1 For a review of the Sine-Gordon description of the Kosterlitz-Thouless theory, see . This was the motivation of the authors of , who have undertaken a systematic study of perturbation theory in a two-parameter expansion around the AF point. They calculated the renormalization group beta-functions up to the two-loop coefficients. The beta-function coefficients were also calculated in by a completely different technique based on string theory. The results found in differ from those of at the two-loop level. The question of two-loop beta-function coefficients were considered also in for a class of generalized Sine-Gordon models. The results, when specialised to the case of the ordinary SG model, agree with those of , but disagree with those of . In the short distance expansion of some Sine-Gordon correlation functions were calculated using conformal perturbation theory. This allowed the extraction of the one- and two-loop beta-function coefficients around the AF point. The resulting two-loop beta-functions differ from all the previous results.
In view of the role the Sine-Gordon model is playing in the description of the KT phase transition it is very important to know the correct two-loop beta-function coefficients. The purpose of the present paper is to show that, in fact, the two-loop results of Amit et al. are the correct ones. We show this first by comparing the SG beta-function to known results in the chiral Gross-Neveu model , which is known to be equivalent to the SG model at its AF point. We also check the beta-function coefficients by considering the renormalization of $`2n`$-point functions of exponentials of the SG field.
Following we consider the Euklidean Lagrangian
$$=\frac{1}{2}_\mu \varphi _\mu \varphi +\frac{m_0^2}{2}\varphi ^2+\frac{\alpha _0}{\beta _0^2a^2}\left[1\mathrm{cos}(\beta _0\varphi )\right],$$
(1)
where $`m_0`$ is an IR regulator mass and $`a`$ is the UV cutoff (of dimension length). UV regularized correlation functions are calculated by using
$$G_0(x)=\frac{1}{2\pi }K_0\left(m_0\sqrt{x^2+a^2}\right)$$
(2)
where $`K_0`$ is the modified Bessel function, as the $`\varphi `$ propagator. Our strategy is slightly different from , who really considered the renormalization of the massive SG model (1) of mass $`m_0`$. We treat $`m_0`$ as an IR regulator mass and consider IR stable physical quantities for which we can take the limit $`m_00`$ already at the UV regularized level (before UV renormalization).
The model (1) is renormalizable in a two-parameter perturbative expansion around the point corresponding to the couplings $`\alpha _0=0`$, $`\beta _0^2=8\pi `$. Writing
$$\beta _0^2=8\pi (1+\delta _0)$$
(3)
the two bare expansion parameters are $`\alpha _0`$ and $`\delta _0`$ and physical quantities can be made UV finite by the renormalizations
$`\alpha _0`$ $`=`$ $`Z_\alpha \alpha ;Z_\alpha =1+g_1\delta \mathrm{}+\alpha ^2(\overline{g}_2\mathrm{}^2+g_2\mathrm{})+\delta ^2(\overline{g}_3\mathrm{}^2+g_3\mathrm{})+\mathrm{},`$ (4)
$`1+\delta _0`$ $`=`$ $`Z_\varphi ^1(1+\delta );Z_\varphi =1+f_1\alpha ^2\mathrm{}+\alpha ^2\delta (\overline{f}_2\mathrm{}^2+f_2\mathrm{})+\mathrm{},`$ (5)
where $`\alpha `$ and $`\delta `$ are the renormalized couplings and $`\mathrm{}=\mathrm{ln}\mu a`$ with $`\mu `$ an arbitrary renormalization point. The dots stand for terms higher order in perturbation theory and the numerical coefficients $`g_1,f_1`$ etc. can be calculated by renormalizing correlation functions. The results of Amit et al. are
$$f_1=\frac{1}{32},g_1=2,f_2=\frac{3}{32},g_2=\frac{5}{64},g_3=0,$$
(6)
those of and are
$$f_1=\frac{1}{32},g_1=2,f_2=\frac{1}{32},g_2=\frac{1}{32},g_3=0$$
(7)
and finally found
$$f_1=\frac{1}{32},g_1=2,f_2=\frac{1}{32},g_2=\frac{1}{16},g_3=0.$$
(8)
We see that the one-loop coefficients are the same but not all two-loop coefficients agree. The subject of this paper is to recalculate these numbers.
The renormalization group (RG) beta-functions can be calculated by solving the equations
$$𝒟\alpha =𝒟\delta =0,$$
(9)
where, as usual, the RG operator is defined by
$$𝒟=a\frac{}{a}+\beta _\alpha (\alpha _0,\delta _0)\frac{}{\alpha _0}+\beta _\delta (\alpha _0,\delta _0)\frac{}{\delta _0}.$$
(10)
One finds
$`\beta _\alpha `$ $`=`$ $`g_1\alpha _0\delta _0g_2\alpha _0^3g_3\alpha _0\delta _0^2+\mathrm{},`$ (11)
$`\beta _\delta `$ $`=`$ $`f_1\alpha _0^2+(f_1+f_2)\alpha _0^2\delta _0+\mathrm{}`$ (12)
It is well-known that, in the case of several couplings, the higher beta-function coefficients are not all scheme independent. Indeed, considering the perturbative redefinitions
$$\stackrel{~}{\alpha }_0=\alpha _0+c_1\alpha _0\delta _0+\mathrm{}\stackrel{~}{\delta }_0=\delta _0+c_2\alpha _0^2+\mathrm{}$$
(13)
one finds that in addition to the one loop coefficients $`f_1`$ and $`g_1`$ only the following two two-loop coefficient combinations are invariant.
$$g_3,J=2g_2f_2.$$
(14)
The RG analysis with two couplings can be made similar to the case of a single coupling by changing the variables from $`\alpha _0`$ and $`\delta _0`$ to the pair $`Q`$ and $`\delta _0`$, where $`Q`$ is a RG invariant (solution of the $`𝒟Q=0`$ equation), given in perturbation theory by
$$Q=f_1\alpha _0^2+g_1\delta _0^2+2g_2\alpha _0^2\delta _0+F_2\delta _0^3+\mathrm{},$$
(15)
where $`F_2=\frac{2}{3}g_3\frac{2}{3}g_1+\frac{2g_1}{3f_1}J`$. Now $`Q`$, being a RG invariant, can almost be treated as if it were a numerical constant and $`\delta _0`$ as the ‘true’ coupling. The beta-function in these variables is
$$\beta (\delta _0,Q)=Q+2\delta _0^2+AQ\delta _0+B\delta _0^3+\mathrm{},$$
(16)
where $`A=1J/f_1`$ and $`B=2(Ag_3)/3`$.
It is well-known that the SG model can also be formulated in terms of two fermion fields, interacting with a chirally symmetric current-current interaction . A special case of the two-fermion model corresponds to the $`SU(2)`$-symmetric chiral Gross-Neveu model. This correspondence is evident in the bootstrap aproach, since the SG S-matrix in the limit $`\beta _0\sqrt{8\pi }`$ becomes the $`SU(2)`$ chiral Gross-Neveu S-matrix. This asymptotically free model has to correspond to one of the possible RG trajectories in the two-parameter SG language. It is easy to see that it has to be the $`Q=0`$ trajectory, since this is the only trajectory going through the origin ($`\delta _0=\alpha _0=0`$) of the parameter space. More precisely, the chiral Gross-Neveu model must correspond to the negative half of the $`Q=0`$ trajectory, which is an UV asymptotically free trajectory. Making the identification
$$\delta _0=\frac{1}{\pi }g^2,$$
(17)
where $`g`$ is the coupling of the $`SU(2)`$ Gross-Neveu model, the Gross-Neveu beta-function becomes
$$\beta (g)=\frac{1}{\pi }g^3+\frac{B}{2\pi ^2}g^5+\mathrm{}$$
(18)
Using the results of Amit et al., (6), $`B=2`$, using the results of and , (7), $`B=4/3`$ and finally $`B=8/3`$ if we trust , (8). Comparing (18) to the results of the beta-function calculations performed directly in the fermion language we see that the correct Gross-Neveu beta-function is reproduced if $`B=2`$. Thus the two-loop results of Amit et al. are correct after all! This was the observation<sup>2</sup><sup>2</sup>2We thank P. Forgács who made this observation first and called our attention to it. that served as our motivation for the present study. The correctness of the two-loop Gross-Neveu beta-function coefficient has been checked by studying the system in the presence of an external field . Using this method the value of this coefficient can be read off from the bootstrap S-matrix and the results are in agreement with .
We now turn to the explicit calculation of the renormalization parameters (4,5). The first quantity we consider is the two-point function of the $`U(1)`$ current$`J_\mu =i\frac{\beta _0}{2\pi }ϵ_{\mu \nu }_\nu \varphi `$,
$$J_\mu (x)J_\nu (y)=\frac{d^2p}{(2\pi )^2}\left(\frac{p_\mu p_\nu }{p^2}\delta _{\mu \nu }\right)e^{ip(xy)}I(p).$$
(19)
The advantage of considering this physical quantity is that it is IR stable. Putting $`m_0=0`$ we find
$$I(p)=\frac{2}{\pi }\left\{1+\delta _0+\frac{\alpha _0^2}{32}\left(\mathrm{ln}pa+K+\frac{1}{2}\right)+\frac{\alpha _0^2\delta _0}{16}(\mathrm{ln}pa+K)^2+\mathrm{}\right\},$$
(20)
where $`K=\mathrm{\Gamma }^{}(1)1\mathrm{ln}2`$. Since the current is conserved there is no operator renormalization required here and (20) must become finite after the substitutions (4,5). From this requirement we get
$$f_1=\frac{1}{32},g_1=2,f_2=\frac{3}{32}.$$
(21)
To determine the remaining two-loop coefficients $`g_2`$ and $`g_3`$ we have to calculate $`Z_\alpha `$, the renormalization constant corresponding to $`\alpha _0`$. For this purpose we need a quantity with a perturbative series starting at $`𝒪(\alpha _0)`$. We have chosen the $`2n`$-point correlation function
$$X=𝒜(x_1)\mathrm{}𝒜(x_{2n}),$$
(22)
where
$$𝒜(x)=\left(\frac{1}{a}\right)^{\frac{1}{2n^2}}e^{\frac{i\beta _0}{2n}\varphi (x)}.$$
(23)
Although, in contrast to the Noether current, the operator (23) needs to be renormalized, for large enough $`n`$ the dimension of (23) is so small that there is no operator mixing and the operator renormalization constant can simply be determined from the correlation function
$$Y=𝒜(x_1)\mathrm{}𝒜(x_n)𝒜^{}(y_1)\mathrm{}𝒜^{}(y_n).$$
(24)
A second order calculation gives
$$\begin{array}{cc}\hfill Y=& M^{\left(\frac{1}{n^2}\right)}\{1+\frac{\delta _0}{n^2}L+\frac{\delta _0^2}{2n^4}L^2+\frac{\alpha _0^2}{64n^3}L^2\hfill \\ & +L(\frac{\alpha _0^2}{64n^2}\frac{\alpha _0^2}{128}[W\left(\frac{1}{n}\right)+W(\frac{1}{n})])+\mathrm{}\},\hfill \end{array}$$
(25)
where
$$M=\frac{_{i<j}|x_ix_j|_{k<l}|y_ky_l|}{_{i,k}|x_iy_k|},$$
(26)
$`L=\mathrm{ln}Ma^n`$ and the dots stand for finite $`𝒪(\alpha _0^2)`$ terms as well as higher order terms. $`W(\mu )`$ is defined by
$$W(\mu )=1+_0^1𝑑zz^\mu F(\mu ,\mu ;1;z)+_0^1\frac{dz}{z^2}\left[F(\mu ,\mu ;1;z)1\mu ^2z\right],$$
(27)
where $`F(\alpha ,\beta ;\gamma ;z)`$ is the standard hypergeometric function. (25) can be made finite by the renormalization $`Y_R=Z_{2n}Y`$, where
$$Z_{2n}=1\frac{1}{n}\mathrm{}\delta +\frac{1}{2n^2}\mathrm{}^2\delta ^2+\frac{1}{64n}\mathrm{}^2\alpha ^2+k_1\mathrm{}\alpha ^2+\mathrm{},$$
(28)
with
$$k_1=\frac{1}{64n}+\frac{n}{128}\left[W\left(\frac{1}{n}\right)+W\left(\frac{1}{n}\right)\right].$$
(29)
For the $`2n`$-point function $`X`$ a second order calculation gives
$$\begin{array}{cc}\hfill X=& \frac{\alpha _0}{16\pi }N^{\left(\frac{1}{n^2}\right)}F\{1+\delta _0\mathrm{\Psi }+\frac{1}{2}\delta _0^2\mathrm{\Psi }^2\hfill \\ & +\frac{n\alpha _0^2}{128n+64}\mathrm{\Psi }[\mathrm{\Psi }+4+\frac{1}{n}nW\left(\frac{1}{n}\right)]+\mathrm{}\},\hfill \end{array}$$
(30)
where
$$N=\underset{i<j}{}|x_ix_j|,F=d^2z\frac{1}{_i|zx_i|^{\frac{2}{n}}}$$
(31)
and
$$\mathrm{\Psi }=1+\frac{1}{n^2}\mathrm{ln}\left(Na^{n(2n1)}\right)\frac{2}{nF}\underset{j}{}d^2z\frac{\mathrm{ln}\left|\frac{zx_j}{a}\right|}{_i|zx_i|^{\frac{2}{n}}}.$$
(32)
In (30) the dots represent finite terms of $`𝒪(\alpha _0^2)`$ and $`𝒪(\delta _0^2)`$ as well as higher terms. Renormalizing $`X`$ by requiring $`X_R=Z_{2n}X`$ to be finite after coupling constant renormalization gives
$$g_3=0\mathrm{and}g_2=\frac{1}{16}+\frac{n}{128}\left[W\left(\frac{1}{n}\right)W\left(\frac{1}{n}\right)\right].$$
(33)
At first sight $`g_2`$ seems to be $`n`$-dependent which would mean that the $`2n`$-point function (22) cannot really be made finite with wave function plus coupling constant renormalization. Luckily, however, one can see that using the identity
$$W(\mu )W(\mu )=2\mu (|\mu |<1)$$
(34)
satisfied by the hypergeometric function, $`g_2`$ is equal to the $`n`$-independent constant $`5/64`$. Moreover, (33) together with (21) reproduce (6), the results of . The nontrivial cancellation of the $`n`$-dependence makes us more confident that these are the correct two-loop coefficients.
Acknowledgements
We would like to thank P. Forgács for calling our attention to the fact that the results of reproduce the correct two-loop coefficients for the chiral Gross-Neveu model. We thank R. Konik for a correspondence.
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# Complementarity of the Maldacena and Randall-Sundrum Pictures*footnote **footnote *Research supported in part by NSF Grant PHY-9722090.
## Abstract
We revive an old result, that one-loop corrections to the graviton propagator induce $`1/r^3`$ corrections to the Newtonian gravitational potential, and compute the coefficient due to closed loops of the U(N) $`𝒩=4`$ super-Yang-Mills theory that arises in Maldacena’s AdS/CFT correspondence. We find exact agreement with the coefficient appearing in the Randall-Sundrum brane-world proposal. This provides more evidence for the complementarity of the two pictures.
preprint: UM-TH-00-12 hep-th/0003237
It is an old, and seemingly forgotten, result that one-loop corrections to the graviton propagator induce $`1/r^3`$ corrections to the gravitational potential :
$$V(r)=\frac{Gm_1m_2}{r}\left(1+\frac{\alpha G}{r^2}\right),$$
(1)
where $`G`$ is the four-dimensional Newton’s constant, $`\mathrm{}=c=1`$ and $`\alpha `$ is a purely numerical coefficient given, in the case of spins $`s1`$, by $`45\pi \alpha =12N_1+3N_{1/2}+N_0`$, where $`N_s`$ are the numbers of particle species of spin $`s`$ going around the loop . However, the importance of this result has recently become apparent in attempts to relate two topical but, at first sight, different developments in quantum gravity. These are Maldacena’s AdS/CFT correspondence and the Randall-Sundrum brane-world mechanism .
The AdS/CFT correspondence in general relates the gravitational dynamics of a $`(d+1)`$-dimensional anti-de Sitter spacetime, AdS<sub>d+1</sub>, to a $`d`$-dimensional conformal field theory, CFT<sub>d</sub>. In the case of $`d=4`$, Maldacena’s conjecture, based on the decoupling limit of D3-branes in Type $`IIB`$ string theory compactified on $`S^5`$, then relates the dynamics of AdS<sub>5</sub> to an $`𝒩=4`$ superconformal $`U(N)`$ Yang-Mills theory on its four-dimensional boundary . Other compactifications are also possible, leading to different SCFT’s on the boundary. We note that, by choosing Poincaré coordinates on AdS<sub>5</sub>, the metric may be written as
$$ds^2=e^{2y/L}(dx^\mu )^2+dy^2,$$
(2)
in which case the superconformal Yang-Mills theory is taken to reside at the boundary $`y\mathrm{}`$.
The Randall-Sundrum mechanism, on the other hand, was originally motivated, not via the decoupling of gravity from D3-branes, but rather as a possible mechanism for evading Kaluza-Klein compactification by localizing gravity in the presence of an uncompactified extra dimension. This was accomplished by inserting a positive tension 3-brane (representing our spacetime) into AdS<sub>5</sub>. In terms of the Poincaré patch of AdS<sub>5</sub> given above, this corresponds to removing the region $`y<0`$, and either joining on a second partial copy of AdS<sub>5</sub>, or leaving the brane at the end of a single patch of AdS<sub>5</sub>. In either case the resulting Randall-Sundrum metric is given by
$$ds^2=e^{2|y|/L}(dx^\mu )^2+dy^2,$$
(3)
where $`y(\mathrm{},\mathrm{})`$ or $`y[0,\mathrm{})`$ for a ‘two-sided’ or ‘one-sided’ Randall-Sundrum brane respectively.
The similarity of these two scenarios led to the notion that they are in fact closely tied together. To make this connection clear, consider the one-sided Randall-Sundrum brane. By introducing a boundary in AdS<sub>5</sub> at $`y=0`$, this model is conjectured to be dual to a cutoff CFT coupled to gravity, with $`y=0`$, the location of the Randall-Sundrum brane, providing the UV cutoff. This extended version of the Maldacena conjecture then reduces to the standard AdS/CFT duality as the boundary is pushed off to $`y\mathrm{}`$, whereupon the cutoff is removed and gravity becomes completely decoupled. Note in particular that this connection involves a single CFT at the boundary of a single patch of AdS<sub>5</sub>. For the case of a brane sitting between two patches of AdS<sub>5</sub>, one would instead require two copies of the CFT, one for each of the patches.
It has been suggested that a crucial test of this Randall-Sundrum version of the Maldacena conjecture would be to compare the $`1/r^3`$ corrections to Newton’s law in both pictures. From the above, we see that the contribution of a single CFT, with $`(N_1,N_{1/2},N_0)=(N^2,4N^2,6N^2)`$, is
$$V(r)=\frac{Gm_1m_2}{r}\left(1+\frac{2N^2G}{3\pi r^2}\right).$$
(4)
Using the AdS/CFT relation $`N^2=\pi L^3/2G_5`$ and the one-sided brane-world relation $`G=2G_5/L`$ , where $`G_5`$ is the five-dimensional Newton’s constant and $`L`$ is the radius of AdS<sub>5</sub>, this becomes
$$V(r)=\frac{Gm_1m_2}{r}\left(1+\frac{2L^2}{3r^2}\right).$$
(5)
The coefficient of the $`1/r^3`$ term is $`2/3`$ of the Randall-Sundrum result quoted in , but in fact agrees with the more thorough analysis of . We shall confirm below that a more careful analysis of the Randall-Sundrum picture using the results of yields exactly the same answer as the above AdS/CFT calculation, thus providing strong evidence for the conjectured duality of the two pictures.
First we derive (4) in more detail by computing the lowest order quantum corrections to solutions of Einstein’s equations. Working with linearized gravity, we begin by writing the metric as
$$g_{\mu \nu }=\eta _{\mu \nu }+h_{\mu \nu },$$
(6)
so that
$$\sqrt{g}g^{\mu \nu }\stackrel{~}{g}^{\mu \nu }=\eta ^{\mu \nu }\stackrel{~}{h}^{\mu \nu }+\mathrm{},$$
(7)
where
$$\stackrel{~}{h}_{\mu \nu }=h_{\mu \nu }\frac{1}{2}\eta _{\mu \nu }h^\alpha {}_{\alpha }{}^{}.$$
(8)
In harmonic gauge, $`_\mu \stackrel{~}{g}^{\mu \nu }=0`$ (i.e. $`_\mu \stackrel{~}{h}^{\mu \nu }=0`$), the classical linearized Einstein equation reads
$$\text{ }\text{ }\stackrel{~}{h}_{\mu \nu }^c(x)=16\pi GT_{\mu \nu }^{}(x),$$
(9)
where the superscript $`c`$ denotes the classical contribution. Fourier transforming to momentum space results in
$$\stackrel{~}{h}_{\mu \nu }^c(p)=16\pi G\mathrm{\Delta }_4(p)T_{\mu \nu }^{}(p),$$
(10)
where $`\mathrm{\Delta }_4(p)=1/p^2`$ is the four-dimensional massless scalar propagator.
Incorporating one-loop corrections, the quantum corrected metric becomes
$$\stackrel{~}{h}_{\mu \nu }^{}=\stackrel{~}{h}_{\mu \nu }^c+\stackrel{~}{h}_{\mu \nu }^q,$$
(11)
where the quantum correction $`\stackrel{~}{h}_q^{\mu \nu }`$ is given in momentum space by
$$\stackrel{~}{h}_q^{\mu \nu }(p)=D^{\mu \nu \alpha \beta }(p)\mathrm{\Pi }_{\alpha \beta \gamma \delta }(p)\stackrel{~}{h}_c^{\gamma \delta }(p).$$
(12)
$`D^{\mu \nu \alpha \beta }`$ is the graviton propagator,
$$D^{\mu \nu \alpha \beta }(p)=\frac{1}{2}\mathrm{\Delta }_4(p)(\eta ^{\mu \alpha }\eta ^{\nu \beta }+\eta ^{\mu \beta }\eta ^{\nu \alpha }\eta ^{\mu \nu }\eta ^{\alpha \beta }+\mathrm{}),$$
(13)
and $`\mathrm{\Pi }_{\alpha \beta \gamma \delta }`$ is the one-loop graviton self-energy, which by symmetry and Lorentz invariance must be of the general form
$`\mathrm{\Pi }_{\alpha \beta \gamma \delta }(p)`$ $`=`$ $`p^4[\mathrm{\Pi }_1(p^2)\eta _{\alpha \beta }\eta _{\gamma \delta }+\mathrm{\Pi }_2(p^2)(\eta _{\alpha \gamma }\eta _{\beta \delta }+\eta _{\alpha \delta }\eta _{\beta \gamma })+\mathrm{\Pi }_3(p^2)(\eta _{\alpha \beta }\widehat{p}_\gamma \widehat{p}_\delta +\eta _{\gamma \delta }\widehat{p}_\alpha \widehat{p}_\beta )`$ (15)
$`+\mathrm{\Pi }_4(p^2)(\eta _{\alpha \gamma }\widehat{p}_\beta \widehat{p}_\delta +\eta _{\alpha \delta }\widehat{p}_\beta \widehat{p}_\gamma +\eta _{\beta \gamma }\widehat{p}_\alpha \widehat{p}_\delta +\eta _{\beta \delta }\widehat{p}_\alpha \widehat{p}_\gamma )+\mathrm{\Pi }_5(p^2)\widehat{p}_\alpha \widehat{p}_\beta \widehat{p}_\gamma \widehat{p}_\delta ].`$
The ellipses in (13) refer to gauge dependent terms in the propagator which make no contribution if coupled to conserved sources. Combining (12), (13) and (15), one thus obtains the quantum corrected metric in the form
$$h_{\mu \nu }^q(p)=p^2\left[2\mathrm{\Pi }_2(p)\delta _\mu ^\alpha \delta _\nu ^\beta +\mathrm{\Pi }_1(p)\eta _{\mu \nu }^{}\eta _{}^{\alpha \beta }+(\mathrm{\Pi }_3(p)+\mathrm{})\widehat{p}_\mu ^{}\widehat{p}_\nu ^{}\eta _{}^{\alpha \beta }\right]\stackrel{~}{h}_{\alpha \beta }^c,$$
(16)
where non-physical gauge-dependent terms have again been dropped. Finally, combining both classical and one-loop quantum results at the linearized level yields
$`h_{\mu \nu }(p)`$ $`=`$ $`16\pi G\mathrm{\Delta }_4(p)[T_{\mu \nu }(p)\frac{1}{2}\eta _{\mu \nu }T^\alpha {}_{\alpha }{}^{}(p)]`$ (18)
$`16\pi G[2\mathrm{\Pi }_2(p)T_{\mu \nu }(p)+\mathrm{\Pi }_1(p)\eta _{\mu \nu }T^\alpha {}_{\alpha }{}^{}(p)].`$
Note that we have ignored the gauge-dependent term in $`h_{\mu \nu }`$ proportional to $`\widehat{p}_\mu \widehat{p}_\nu `$. It makes no contribution when $`h_{\mu \nu }`$ is attached to a conserved source $`T_{\mu \nu }`$ satisfying $`p^\mu T_{\mu \nu }=p^\nu T_{\mu \nu }=0`$.
The actual form of the one-loop $`\mathrm{\Pi }_i`$’s depend on the theory at hand. However for any massless theory in four-dimensions, after cancelling the infinities with the appropriate counterterms, the finite remainder must necessarily have the form
$$\mathrm{\Pi }_i(p)=32\pi G\left(a_i\mathrm{ln}\frac{p^2}{\mu ^2}+b_i\right),$$
(19)
where $`a_i`$ and $`b_i`$, $`(i=1,2,3,4,5)`$, are numerical coefficients and $`\mu `$ is an arbitrary subtraction constant having the dimensions of mass. In order to make connection with the Newtonian potential, we Fourier transform (18) back to coordinate space. For the static potential we obtain the expected $`1/r`$ behavior at the classical level, while the quantum term generates the claimed $`1/r^3`$ correction. In addition, the constant parts in (19) give rise to a regulator-dependent $`\delta ^3(𝐫)`$ contact interaction. However we have no real expectation that this one-loop perturbative result remains valid when continued down to zero size. Moreover, possible $`r^3\mathrm{ln}\mu r`$ terms come only from the $`\widehat{p}_\mu \widehat{p}_\nu `$ terms in (16) and hence drop out. For a point source, $`T_{00}(x)=m\delta ^3(𝐫)`$, we obtain to this order
$`g_{00}`$ $`=`$ $`\left(1{\displaystyle \frac{2Gm}{r}}{\displaystyle \frac{2\alpha G^2m}{r^3}}\right),`$ (20)
$`g_{ij}`$ $`=`$ $`\left(1+{\displaystyle \frac{2Gm}{r}}+{\displaystyle \frac{2\beta G^2m}{r^3}}\right)\delta _{ij},`$ (21)
where, in agreement with , $`\alpha =432\pi (a_1+2a_2)`$ and $`\beta =432\pi a_1`$. This yields the potential given in (1). Explicit calculations of the self-energy (19) for spin 1 , spin 1/2 (two-component fermions) and (real conformally coupled) spin 0 yield<sup>§</sup><sup>§</sup>§Note that a symmetry factor of $`1/2`$ was omitted in Ref. ; this was subsequently corrected in Ref.
$`a_i(s=1)=4a_i(s=1/2)`$ $`=12a_i(s=0)`$ (23)
$`={\displaystyle \frac{1}{120(4\pi )^2}}(2,3,2,3,4).`$
Note that all spins contribute with the same sign as they must by general positivity arguments on the self-energy . Thus
$$\alpha =2\beta =\frac{1}{45\pi }(12N_1+3N_{1/2}+N_0)=\frac{2N^2}{3\pi },$$
(24)
as quoted in the introductory paragraph above.
This $`\alpha `$ coefficient also determines that part of the Weyl anomaly involving the square of the Weyl tensor :
$$g_{\mu \nu }T^{\mu \nu }=b\left(F+\frac{2}{3}\text{ }\text{ }R\right)+b^{}G,$$
(25)
where
$`F`$ $`=`$ $`C_{\mu \nu \rho \sigma }C^{\mu \nu \rho \sigma }=R_{\mu \nu \rho \sigma }R^{\mu \nu \rho \sigma }2R_{\mu \nu }R^{\mu \nu }+R^2,`$ (26)
$`G`$ $`=`$ $`R_{\mu \nu \rho \sigma }R^{\mu \nu \rho \sigma }=R_{\mu \nu \rho \sigma }R^{\mu \nu \rho \sigma }4R_{\mu \nu }R^{\mu \nu }+{\displaystyle \frac{1}{3}}R^2,`$ (27)
and where $`b`$ and $`b^{}`$ are constants
$`b`$ $`=`$ $`{\displaystyle \frac{1}{120(4\pi )^2}}[12N_1+3N_{1/2}+N_0],`$ (28)
$`b^{}`$ $`=`$ $`{\displaystyle \frac{1}{720(4\pi )^2}}[124N_1+11N_{1/2}+2N_0].`$ (29)
Note that for the $`𝒩=4`$ SCFT, the coefficient of the (Riemann)<sup>2</sup> term, $`b+b^{}`$, vanishes . The same result is obtained if one calculates the holographic Weyl anomaly using the AdS/CFT correspondence . Thus $`b=3\alpha /128\pi =c/(4\pi )^2`$, where the $`c`$ is the central charge given in the normalization of . For the central charge, one obtains $`c=\pi L^3/8G_5`$ , so that
$$G\alpha =\frac{GL^3}{3G_5}=\frac{2L^2}{3},$$
(30)
where the second equality makes use of the brane-world relation $`G=2G_5/L`$. Although we have focused on the $`𝒩=4`$ SCFT to relate the coefficient appearing in Newton’s law to the central charge, the result (30) is universal, being independent of which particular CFT appears in the AdS/CFT correspondence, which is just as well since the Randall-Sundrum coefficient does not depend on the details of the fields propagating on the brane.
We now turn to this brane-world, where the five-dimensional action has the form
$$S=d^5x\sqrt{g_{(5)}}[M^3R_{(5)}\mathrm{\Lambda }]+d^4x\sqrt{g_{(4)}}_{\mathrm{brane}}.$$
(31)
Here $`M`$ is the five-dimensional Planck mass, $`M^3=1/(16\pi G_5)`$, and $`\mathrm{\Lambda }`$ is the cosmological constant in the bulk. Small fluctuations of the metric on the brane may be represented by
$$ds^2=e^{2|y|/L}[\eta _{\mu \nu }+h_{\mu \nu }(x,y)]dx^\mu dx^\nu +dy^2,$$
(32)
where $`L`$ is the ‘radius’ of AdS,
$$R_{MNPQ}^{(5)}=\frac{1}{L^2}(g_{MP}^{(5)}g_{NQ}^{(5)}g_{MQ}^{(5)}g_{NP}^{(5)}),$$
(33)
and is related to $`\mathrm{\Lambda }`$ by $`\mathrm{\Lambda }=12M^3/L^2`$. The brane-world geometry has been chosen such that $`x^\mu `$ are coordinates along the 3-brane, while $`y`$ is the coordinate perpendicular to the brane (which sits at $`y=0`$).
Both brane and bulk quantities are contained in the linearized metric $`h_{\mu \nu }(x,y)`$. However, for comparison with the CFT on the brane, we are only concerned with the former. Hence we consider a matter source on the brane, and examine $`h_{\mu \nu }(x)h_{\mu \nu }(x,y=0)`$. For this case, the results of indicate
$`h_{\mu \nu }(p)`$ $`=`$ $`{\displaystyle \frac{2}{LM^3}}\mathrm{\Delta }_4(p)[T_{\mu \nu }(p)\frac{1}{2}\eta _{\mu \nu }T^\alpha {}_{\alpha }{}^{}(p)]`$ (35)
$`{\displaystyle \frac{1}{M^3}}\mathrm{\Delta }_{KK}(p)[T_{\mu \nu }(p)\frac{1}{3}\eta _{\mu \nu }T^\alpha {}_{\alpha }{}^{}(p)].`$
This expression has a clear physical meaning; $`\mathrm{\Delta }_4(p)`$, the four-dimensional massless propagator, corresponds to the zero-mode graviton localized on the brane, while
$$\mathrm{\Delta }_{KK}(p)=\frac{1}{p}\frac{K_0(pL)}{K_1(pL)}$$
(36)
is the propagator for the continuum Kaluza-Klein graviton modes. Comparing the first term of (35) to (18), we obtain the relation between four- and five-dimensional Newton’s constants, $`G=2G_5/L=1/(8\pi LM^3)`$ given above. Note that in the above we have taken the brane to be at the end of a single patch of AdS<sub>5</sub>, as was done in . This corresponds to the case at hand, since the AdS/CFT relations we have employed above pertain to a single copy of AdS<sub>5</sub>.
The continuum graviton modes give rise to corrections to the Newtonian potential. At large distances, corresponding to $`pL1`$, a small argument expansion for Bessel functions yields
$$\mathrm{\Delta }_{KK}(p)=\frac{L}{2}\left(\mathrm{ln}\frac{p^2L^2}{4}+2\gamma \right)+𝒪(p^2),$$
(37)
and, just as in (19), is the source of the $`1/r^3`$ correction to the Newtonian potential. For a static gravitational source of mass $`m`$ on the brane, $`T_{00}(p)=2\pi \delta (p_0)m`$, evaluating the Fourier transform for $`rL`$ yields the linearized metric
$`h_{00}`$ $`=`$ $`{\displaystyle \frac{2Gm}{r}}\left(1+{\displaystyle \frac{2L^2}{3r^2}}+\mathrm{}\right),`$ (38)
$`h_{ij}`$ $`=`$ $`{\displaystyle \frac{2Gm}{r}}\left(1+{\displaystyle \frac{L^2}{3r^2}}+\mathrm{}\right)\delta _{ij},`$ (39)
from which one may read off the Newtonian potential (5).
Moreover, all the metric components in (39) agree with those of (21) and not merely the $`g_{00}`$ component. In momentum space, this may be traced to the behavior of $`h_{\mu \nu }`$ in the two pictures, namely (18) and (35). In (35) the factor of $`1/3`$ in the non-leading term, as compared with factor $`1/2`$ in the leading term, is attributable to the fact that the Kaluza-Klein gravitons are massive. Whereas in (18), it is because the CFT requires loop corrections with $`\mathrm{\Pi }_2(p)=\frac{3}{2}\mathrm{\Pi }_1(p)`$, which is in fact satisfied, as far as the $`\mathrm{ln}p^2`$ term is concerned, since $`a_2=\frac{3}{2}a_1`$.
We have thus demonstrated that the $`1/r^3`$ corrections to Newton’s law are identical between the Maldacena and Randall-Sundrum pictures. This was examined in the context of a single CFT corresponding to a one-sided brane-world scenario. Had we chosen instead to take the brane-world to be sitting between two patches AdS<sub>5</sub> (one on either side), as was the case considered in , we would have obtained a factor of two in the relation between Newton’s constants, with a corresponding factor in the propagator, (35). While this would ensure the correct four-dimensional behavior of gravity, given in (39), the two-sided brane-world relation $`G=G_5/L`$ will modify the comparison with the one-loop CFT result, (30). To compensate for this mismatch, one may assume that the two-sided brane-world is dual to two copies of the CFT coupled to gravity, as is implicit in . This leads to the natural picture that a one-sided brane corresponds to a single CFT while a two-sided brane corresponds to two CFTs.
An intriguing feature of this comparison of the gravitational potential in both pictures is a highlighting of the classical/quantum nature of this duality, as seen in the relation
$$\mathrm{\Pi }_2(p)+𝒪(G^2)=\frac{L}{4}\mathrm{\Delta }_{KK}(p).$$
(40)
The propagator for the continuum graviton modes in the Randall-Sundrum picture thus incorporates all quantum effects of matter on the brane. It may be worthwhile to examine this relation at the two-loop or higher level. Nevertheless, this agreement at one-loop lends strong support to the conjectured duality between the two pictures.
We would like to acknowledge the referees for providing valuable comments. JTL wishes to thank I. Giannakis and H.C. Ren for fruitful discussions on linearized gravity in the brane-world.
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# On the Particle Definition in the presence of Black Holes
## 1 Introduction
At present we know two very fundamental and effectual theories in order to describe nature, quantum theory and general relativity. But a satisfactory unification of these distinct theories is still missing. One possibility to achieve some progress towards this aim is expected to be provided by the consideration of quantised fields in given classical space-times. Within this treatment the metric plays the role of an external background field. Various investigations have been devoted to this topic during the last decades, to mention only some of the most important initial papers in chronological order: In 1972 Fulling noticed the non-uniqueness of the particle interpretation in curved space-times (which may be regarded as the basis for several effects). Two years later Hawking found out that black holes are not completely black but possess a thermal behaviour caused by the conversion of the initial vacuum fluctuations. A short time after this striking discovery Davies showed that also a uniformly accelerated mirror – treated as a mirror at rest in the Rindler metric – creates a thermal spectrum. In 1976 Unruh recognised the fact that even a uniformly accelerated observer in the Minkowski vacuum feels environed by a thermal bath. Many examinations have been accomplished since these basic papers, see e.g. and references therein. Now black holes are very interesting touchstones in order to test candidates for the theory that unifies general relativistic and quantum aspects. The representations of black holes within the underlying theory are expected to reproduce their main properties.
The main intention of this article is to provide a canonical approach to quantum field theory in specific curved space-times and a related particle definition, together with an investigation of the consequences of this formalism.
This paper is organised as follows: In Section 2 we set up the equations which describe the quantum field under consideration together with the assumptions necessary for the particle definition. Based on methods of functional analysis we propose in Sec. 3 a canonical approach to the particle definition via the diagonalisation of the Hamiltonian. This procedure is applied to some flat space-times in Section 4 in order to elucidate the underlying mechanism. Sec. 5 considers the space-time of a black hole and Section 5.1 presents the particle definition for the exterior domain. In Sec. 5.2 we examine the quantised field inside the black hole and deduce the unstable behaviour of its time-evolution. Possible consequences of this instability are discussed in Section 5.3 and its implications for the sonic analogues of black holes are pointed out in Section 5.4. We shall close with a summary, some discussions and an outline.
Throughout this article natural units with $`G=\mathrm{}=c=k_\mathrm{B}=1`$ will be used. Lowercase Greek indices such as $`\mu ,\nu `$ vary from 0 (time) to 3 (space) and describe space-time components. Lowercase Roman indices $`i,j`$ range from 1 to 3 and label only the space. (For both we employ the Einstein sum convention.) Uppercase Roman indices $`I,J`$ may assume all natural numbers while uppercase Greek indices like $`\mathrm{\Gamma },\mathrm{\Lambda }`$ may be more general, e.g. continuous.
## 2 Equations of motion
We consider a minimally coupled, massless and neutral (i.e. real) scalar field $`\mathrm{\Phi }`$ whose propagation in the space-time $`(M,g^{\mu \nu })`$ is described by the action
$`𝒜={\displaystyle \underset{M}{}}d^4x{\displaystyle \frac{\sqrt{g}}{2}}g^{\mu \nu }(_\mu \mathrm{\Phi })(_\nu \mathrm{\Phi }),`$ (1)
with $`g=det(g_{\mu \nu })`$. Possible potential terms like a mass term $`m^2\mathrm{\Phi }^2`$ or a conformal coupling term $`\mathrm{\Phi }^2/6`$ (where $`=_\mu ^\mu `$ denotes the Ricci scalar) do not alter the main conclusions, see Secs. 3.2 and 5.2. The same holds true for a charged and thus complex field $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^{}`$. For reasons of simplicity and considering the scalar field $`\mathrm{\Phi }`$ as a model for the photon field we restrict ourselves to the most simple action in Eq. (1).
Provided that the spatial surface terms arising from the integration by parts vanish the variation of the action $`\delta 𝒜=0`$ leads to the Klein-Fock-Gordon equation
$`\mathrm{}\mathrm{\Phi }={\displaystyle \frac{1}{\sqrt{g}}}_\mu \left(\sqrt{g}g^{\mu \nu }_\nu \mathrm{\Phi }\right)=0.`$ (2)
The corresponding inner product is defined by
$`\left(\psi |\varphi \right)\stackrel{\mathrm{def}}{=}i{\displaystyle \underset{\mathrm{\Sigma }}{}}𝑑\mathrm{\Sigma }^\mu \psi ^{}\underset{\mu }{\overset{}{}}\varphi `$ (3)
with $`\psi \underset{\mu }{\overset{}{}}\varphi =\psi _\mu \varphi \varphi _\mu \psi `$. In this definition the surface element $`d\mathrm{\Sigma }_\mu `$ already contains factors like $`\sqrt{g_\mathrm{\Sigma }}`$ with the result $`d\mathrm{\Sigma }_\mu dx^\mu =d^4x\sqrt{g}`$. For functions which fulfil the Klein-Fock-Gordon equation $`\mathrm{}\psi =\mathrm{}\varphi =0`$ the inner product is independent of the special surface $`\mathrm{\Sigma }`$, cf. for instance . To see that, one has to use Gauss’ law
$`{\displaystyle \underset{M_\mathrm{\Sigma }}{}}𝑑\mathrm{\Sigma }_\mu A^\mu ={\displaystyle \underset{M_\mathrm{\Sigma }}{}}d^4x\sqrt{g}_\mu A^\mu `$ (4)
and again to require vanishing spatial surface terms.
### 2.1 Preconditions
Now we have to specify the assumptions which are necessary for an appropriate particle definition. At first we demand a strongly causal space-time $`M`$. This condition ensures the essential physical principle of distinguishing cause and effect and forbids (for instance) the occurrence of closed time-like curves, cf. .
As another requirement we impose a static metric of the space-time $`M`$
$`ds^2=g_{\mu \nu }dx^\mu dx^\nu =g_{00}(𝒓)dt^2+g_{ij}(𝒓)dx^idx^j.`$ (5)
These two assumptions allow to factorise the space-time $`M=G`$ into time $`t`$ and space $`𝒓G`$ with an open domain $`G^3`$. The Killing vector corresponding to the time translation symmetry permits the definition of a conserved energy. This fact is substantial for a physical reasonable particle definition, see also Sec. 2.2 below. On the other hand, the selection of a particular Killing vector refers to a class of associated observers whose time evolution is generated by this vector field. In general different Killing vectors generate distinct particle definitions applying for the different observers, see Sec. 4.3 below.
The third precondition we need is called non-degenerated signature. This fixes the signature of the metric inside the domain $`G`$
$`𝒓G:g_{00}(𝒓)>0;\left(g_{ij}(𝒓)\right)_{ij}<0.`$ (6)
In the latter inequality $`\left(g_{ij}(𝒓)\right)_{ij}`$ has to be understood as a matrix (and not as the single components), i.e. $`𝒑^3:𝒑^2>0p_ig^{ij}p_j<0`$. Both quantities ($`g_{00}`$ and $`g_{ij}`$) are continuous and regular inside $`G`$ but may diverge or vanish by approaching the boundary $`G`$. This might be the case for a horizon situated at $`G`$. (A $`g^{00}`$-component of the metric that vanishes over a finite volume would create primary constraints, see e.g. .)
As demonstrated above the possibility of performing the integration by parts is a really important issue. Accordingly, our last assumption is a physical complete region $`G`$. This simply enforces the vanishing of the surface terms. The occurrence of such boundary contributions always indicates the interaction with a system behind the surface. Such a region would not be physical complete. It should be noted here that the validity of the integration by parts also includes periodicity in angular coordinates, such as $`f(\phi )=f(\phi +2\pi )`$. For the above specified space-time $`M=G`$ the spatial surface terms read
$`{\displaystyle \underset{G}{}}𝑑S_i\mathrm{\Phi }g^{ij}_j\mathrm{\Phi }=0.`$ (7)
There are several ways to achieve the equation above. For Dirichlet boundary conditions one would demand $`\mathrm{\Phi }=0`$ at $`G`$ and for Neumann type $`dS_ig^{ij}_j\mathrm{\Phi }=0`$ at $`G`$. But there is also a third possibility for the disappearance of the surface terms, namely if the components of the metric themselves which are orthogonal to the surface $`dS_i`$ approach zero at $`G`$, i.e. $`dS_ig^{ij}(G)=0`$. As stated above this might be the case for a horizon situated at the boundary $`G`$.
Strictly speaking, there exist various definitions of a horizon, such as the event, apparent, Cauchy, particle, and putative horizon, cf. and . The definition of the particle horizon refers to a special observer at a given time-like word-line whereas the other horizons can be defined in an observer-independent way. Hence, the vanishing of the spatial surface terms (without constraints on the fields) implies a priori only a particle horizon at $`G`$. However, with some additional requirements on the space-time, for instance spherical symmetry and asymptotic flatness, the (particle) horizon at $`G`$ meets the other definitions as well.
After having established the properties of the space-time $`M`$, we arrive at the conclusion that it is globally hyperbolic (i.e. strongly causal and complete, cf. ) and the spatial domain $`G`$ represents for every time $`t`$ a Cauchy surface.
### 2.2 Energy
For a time-independent metric the Noether theorem demands the existence of a conserved energy. The energy-momentum tensor for the scalar field reads
$`T_{\mu \nu }={\displaystyle \frac{2}{\sqrt{g}}}{\displaystyle \frac{\delta 𝒜}{\delta g^{\mu \nu }}}=(_\mu \mathrm{\Phi })(_\nu \mathrm{\Phi }){\displaystyle \frac{g_{\mu \nu }}{2}}(_\gamma \mathrm{\Phi })(^\gamma \mathrm{\Phi }).`$ (8)
By virtue of the Klein-Fock-Gordon equation $`\mathrm{}\mathrm{\Phi }=0`$ the covariant divergence of the energy-momentum tensor vanishes
$`_\mu T_\nu ^\mu ={\displaystyle \frac{1}{\sqrt{g}}}_\mu \left(\sqrt{g}T_\nu ^\mu \right){\displaystyle \frac{1}{2}}T^{\alpha \beta }_\nu g_{\alpha \beta }=0.`$ (9)
In general this covariant equation does not lead to any conserved quantities due to the exchange of energy and momentum between the gravitational and the scalar field (second term). But for a stationary metric ($`_0g_{\alpha \beta }=0`$) it is possible to construct a conserved energy flux $`j^\mu `$ utilising the ($`\nu =0`$)-components
$`_\mu j^\mu =_\mu \left(\sqrt{g}T_0^\mu \right)=0.`$ (10)
This local conservation law allows for the introduction of a conserved energy as a global quantity via
$`E\stackrel{\mathrm{def}}{=}{\displaystyle \underset{G}{}}d^3x\sqrt{g}T_0^0={\displaystyle \underset{G}{}}d^3𝒓T_0^0.`$ (11)
For a Minkowski space-time where $`T^{00}=T_0^0=T_{00}`$ holds this definition coincides (of course) with the usual energy. Another argument for the above defined energy for being the correct choice is the following: Starting from the action $`𝒜`$ we may define the Lagrange function $`L`$ such that
$`𝒜\stackrel{\mathrm{def}}{=}{\displaystyle 𝑑tL}`$ (12)
holds. The Hamilton function $`H`$ as the Legendre transform of this Lagrange function exactly coincides with the energy of the field $`H=E`$.
## 3 Particle definition
To provide a canonical definition of particles one has to indicate which properties the particles should exhibit. For a free (linear) field we expect the particles to evolve independently and to carry a certain energy. As shown in Section 2.2, for a static metric the energy $`E`$ of the field $`\mathrm{\Phi }`$ and its Hamilton function $`H`$ coincide. Consequently, both requirements can be satisfied by the diagonalisation of $`E=H`$ or, equivalently, the Lagrange function $`L`$. Having defined the particles via diagonalisation of $`H=E`$, the corresponding vacuum $`|0`$ coincides with the ground state of the Hamiltonian and the energy. Of course, the procedure described above does not represent the only one possibility to accomplish the particle definition. Another approach is based on the ”one-particle structure” of classical solutions of the field equation, see e.g. and the remarks in Sec. 3.6.
According to the definition in Section 2.2 the Lagrange function governing the dynamics of the field reads
$`L={\displaystyle \frac{1}{2}}{\displaystyle \underset{G}{}}d^3𝒓g^{00}(𝒓)\dot{\mathrm{\Phi }}^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{G}{}}d^3𝒓g^{ij}(𝒓)(_i\mathrm{\Phi })(_j\mathrm{\Phi }),`$ (13)
where $`d^3𝒓`$ denotes the spatial integration with the volume element $`d^3𝒓=\sqrt{g}d^3x`$.
To diagonalise this expression one has to deal with an elliptic partial differential operator which requires some functional analysis. All of the used theorems can be found in and are not cited explicitly in the following.
### 3.1 Hilbert space theory
To work with mathematically well-defined quantities we have to set up some definitions. $`C_0^{\mathrm{}}(G)`$ denotes the set of all infinitely differentiable functions $`u:G`$ of compact support inside the open domain $`G`$. For two functions of this kind $`u,vC_0^{\mathrm{}}(G)`$ we define a scalar product via
$`\left\{u\right|v\}_1\stackrel{\mathrm{def}}{=}{\displaystyle \underset{G}{}}d^3𝒓g^{00}(𝒓)u^{}(𝒓)v(𝒓).`$ (14)
The assumption of a non-degenerated signature in Section 2.1 is essential for this definition. Without a positive $`g^{00}`$ the above expression would be a pseudo-scalar product instead of a scalar product with $`\left\{u\right|u\}=0u=0`$. The latter property is necessary for investigations concerning convergence. As every scalar product induces a norm $`u^2=\left\{u\right|u\}`$ it is now possible to define a Hilbert space as the completion of all $`C_0^{\mathrm{}}(G)`$ functions with respect to this norm
$`L_2(G,g^{00})\stackrel{\mathrm{def}}{=}\overline{C_0^{\mathrm{}}(G)}^{\{|\}_1}.`$ (15)
Because every $`C_0^{\mathrm{}}(G)`$-function can be $`L_2(G,g^{00})`$-approximated by linear combinations of step functions, this Hilbert space is separable.
The same procedure may be performed for vector-valued functions $`𝒖:G^3`$. Again we may define a scalar product for two smooth functions of compact support $`𝒖,𝒗[C_0^{\mathrm{}}(G)]^3`$ due to the non-degenerated signature
$`\left\{𝒖\right|𝒗\}_3\stackrel{\mathrm{def}}{=}{\displaystyle \underset{G}{}}d^3𝒓g^{ij}(𝒓)u_i^{}(𝒓)v_j(𝒓),`$ (16)
and in analogy the corresponding Hilbert space reads
$`L_2^3(G,g^{ij})\stackrel{\mathrm{def}}{=}\overline{[C_0^{\mathrm{}}(G)]^3}^{\{|\}_3}.`$ (17)
The advantage of the scalar products defined in such a way becomes evident if we use the linear partial differential operator
$`𝒟:C_0^{\mathrm{}}(G)L_2(G,g^{00})`$ $``$ $`L_2^3(G,g^{ij})`$
$`\varphi (𝒓)`$ $``$ $`\left(_i\varphi (𝒓)\right)_i`$ (18)
to cast the Lagrange function into the simple form
$`L={\displaystyle \frac{1}{2}}\left\{\dot{\mathrm{\Phi }}\right|\dot{\mathrm{\Phi }}\}_1{\displaystyle \frac{1}{2}}\left\{𝒟\mathrm{\Phi }\right|𝒟\mathrm{\Phi }\}_3.`$ (19)
Nevertheless, this is still not a representation which is suitable for diagonalisation. For that purpose we have to perform the spatial integration by parts (see Section 2). In terms of functional analysis this means the construction of the adjoint operator. The domain of definition $`\mathrm{Def}(𝒟)=C_0^{\mathrm{}}(G)`$ of the $`𝒟`$-operator is dense in $`L_2(G,g^{00})`$. As a consequence, its adjoint $`𝒟^{}`$ exists as a linear operator $`𝒟^{}:\mathrm{Def}(𝒟^{})L_2^3(G,g^{ij})L_2(G,g^{00})`$. For $`[C_0^{\mathrm{}}(G)]^3`$-functions the spatial integration by parts is always possible. Accordingly, the domain of definition of the adjoint $`𝒟^{}`$ contains these functions $`[C_0^{\mathrm{}}(G)]^3\mathrm{Def}(𝒟^{})`$ and is thereby also dense in $`L_2^3(G,g^{ij})`$. Therefore the twice adjoint $`𝒟^{}`$ exists as a linear operator as well $`𝒟^{}:\mathrm{Def}(𝒟^{})L_2(G,g^{00})L_2^3(G,g^{ij})`$. Of course, these operators describe physical reality only if one ensures the possibility of the spatial integration by parts via physical reasons as done in Section 2.1.
### 3.2 $`𝒦`$-operator
Now we are in the position to cast the Lagrange function into a form which can be utilised for the diagonalisation of the system. With the definition of the elliptic partial differential operator $`𝒦\stackrel{\mathrm{def}}{=}𝒟^{}𝒟^{}`$ (see also ) we arrive at
$`L`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{\dot{\mathrm{\Phi }}\right|\dot{\mathrm{\Phi }}\}_1{\displaystyle \frac{1}{2}}\left\{\mathrm{\Phi }\right|𝒟^{}𝒟^{}|\mathrm{\Phi }\}_1`$ (20)
$`\stackrel{\mathrm{def}}{=}`$ $`{\displaystyle \frac{1}{2}}\left\{\dot{\mathrm{\Phi }}\right|\dot{\mathrm{\Phi }}\}_1{\displaystyle \frac{1}{2}}\left\{\mathrm{\Phi }\right|𝒦|\mathrm{\Phi }\}_1.`$
Every linear operator of the form $`𝒦=𝒟^{}𝒟^{}`$ is non-negative and self-adjoint and thus can be diagonalised. Let us study the domain of definition $`\mathrm{Def}(𝒦)`$ of this operator. The twice adjoint operator $`𝒟^{}`$ is the closure of the original operator $`𝒟`$, i.e. $`\overline{𝒟}=𝒟^{}`$. Its domain of definition is the completion of all $`C_0^{\mathrm{}}(G)`$-functions
$`\mathrm{Def}(\overline{𝒟})=\overline{C_0^{\mathrm{}}(G)}^{\{|\}_𝒟}`$ (21)
with respect to the graph scalar product which is defined via
$`\left\{u\right|v\}_𝒟\stackrel{\mathrm{def}}{=}\left\{u\right|v\}_1+\left\{𝒟u\right|𝒟v\}_3.`$ (22)
One observes that the operator $`\mathrm{𝟏}+𝒦`$ is exactly the Friedrich extension (which is self-adjoint, see ) of the original operator $`\mathrm{𝟏}+𝒦_{C_0^{\mathrm{}}(G)}`$ mediated via the graph scalar product. As a result, if the domain $`G`$ has boundaries $`G`$ with Dirichlet boundary conditions, these boundary conditions are already incorporated into the domain of definition of the operators $`\overline{𝒟}`$ and $`𝒦`$, i.e.
$`\varphi \mathrm{Def}(𝒦)\mathrm{Def}(\overline{𝒟})\varphi (G)=0.`$ (23)
To incorporate Neumann boundary conditions one has to start with an operator like $`(u^i)^i_iu^i`$ and to proceed in the same way.
As mentioned in Sec. 2, additional potential terms do not alter the main conclusions. If we assume the scalar curvature to be a bounded $`m^2_{\mathrm{Max}}`$ and smooth $`C^{\mathrm{}}(G)`$ function we may introduce a new operator via
$`:L_2(G,g^{00})`$ $``$ $`L_2(G,g^{00})`$
$`\varphi (𝒓)`$ $``$ $`\left(m^2+\right)\varphi (𝒓).`$ (24)
Obviously this operator is bounded, non-negative, and self-adjoint. In addition, since $`\mathrm{Def}(𝒟^{}\overline{𝒟})\mathrm{Def}()=L_2(G,g^{00})`$, we may define a modified $`𝒦`$-operator via
$`𝒦\stackrel{\mathrm{def}}{=}𝒟^{}\overline{𝒟}+_{\mathrm{Def}(𝒟^{}\overline{𝒟})},`$ (25)
which is still self-adjoint and non-negative.
### 3.3 Spectral theory
As mentioned above, every self-adjoint operator can be diagonalised. One way to reveal this statement in a more explicit form is the following theorem: For every self-adjoint operator $`𝒦`$ there exists a spectral family $``$ of orthogonal projections with
$`𝒦={\displaystyle \lambda 𝑑(\lambda )}.`$ (26)
$`d(\lambda )`$ contributes only for values $`\lambda `$ being in the spectrum $`\sigma (𝒦)`$ of the $`𝒦`$-operator $`\lambda \sigma (𝒦)`$. The spectrum $`\sigma (𝒦)`$ of an operator $`𝒦`$ contains all complex numbers $`z`$ for which the resolvent $`(z)\stackrel{\mathrm{def}}{=}(z𝒦)^1`$ does not exist, i.e. $`(z𝒦)^1`$ is not a well and densely defined and bounded operator. For a self-adjoint and non-negative operator $`𝒦`$ the spectrum is purely real and non-negative $`\sigma (𝒦)_+`$. It splits up into two parts, the point spectrum $`\sigma _\mathrm{p}`$ and the continuous spectrum $`\sigma _\mathrm{c}`$. The point spectrum is the set of all proper eigenvalues $`\lambda `$ corresponding to proper eigenfunctions
$`\sigma _\mathrm{p}=\{\lambda :|f_\lambda \}:𝒦\left|f_\lambda \right\}=\left|f_\lambda \right\}\lambda \}.`$ (27)
The continuous spectrum contains all numbers $`\lambda `$ where $`(\lambda 𝒦)^1`$ formally exists, but is not bounded
$`\sigma _\mathrm{c}=\{\lambda \backslash \sigma _\mathrm{p}:(\lambda 𝒦)^1=\mathrm{}\}.`$ (28)
The discrete spectrum $`\sigma _\mathrm{d}`$ is that part of the point spectrum $`\sigma _\mathrm{p}`$ which incorporates all isolated points $`\lambda `$ of $`\sigma _\mathrm{p}`$ with a finite number of corresponding eigenfunctions $`\left|f_\lambda \right\}`$. The continuous spectrum $`\sigma _\mathrm{c}`$ may also be divided into two parts, the absolute continuous spectrum $`\sigma _{\mathrm{ac}}`$, where $`d(\lambda )/d\lambda `$ exists as a weakly integrable operator, and the remaining singular continuous spectrum $`\sigma _{\mathrm{sc}}`$.
To provide some physical insight into these abstract quantities we shall investigate the spectrum for a few examples. The discrete spectrum $`\sigma _\mathrm{d}`$ describes localised states, such as bound states or states of a field confined in a finite volume. The point spectrum $`\sigma _\mathrm{p}`$ may contain more points with additional characteristics. E.g., if the operator $`𝒦`$ governs the dynamics of the Maxwell field $`A_\mu `$ there is an infinite set of eigenfunctions at the point $`\lambda =0`$. These functions correspond to the gauge invariance of this theory and do not change physical quantities. The absolute continuous spectrum $`\sigma _{\mathrm{ac}}`$ represents usually the scattering states, but the singular continuous spectrum $`\sigma _{\mathrm{sc}}`$ may be related to more strange phenomena, like quasi-bound states, scattering states in average, fractal measure $`d\mu _{\mathrm{}}`$ (cf. Sec. 3.4 below), chaotic behaviour, etc.
Fortunately, for smooth and regular coefficients $`g_{\mu \nu }`$ with an appropriate asymptotic behaviour the spectrum of the $`𝒦`$-operator is either purely discrete $`\sigma (𝒦)=\sigma _\mathrm{d}`$ (for a finite volume) or absolute continuous $`\sigma (𝒦)=\sigma _{\mathrm{ac}}`$ (for an infinite volume, see e.g. ).
### 3.4 Spectral theorem
For our main intention, the diagonalisation of the Lagrange function, it is suitable to make use of the following theorem: For every self-adjoint operator $`𝒦`$ acting on a separable Hilbert space there exists a unitary transformation $`𝒰`$ which diagonalises it : $`𝒰𝒦𝒰^{}=`$. $``$ denotes the multiplication by argument: $`(f)(\lambda )=\lambda f(\lambda )`$. Because $`𝒦`$ is $``$-real, i.e. $`(𝒦\varphi )^{}=𝒦(\varphi ^{})`$, we may construct a quasi-unitary transformation
$`𝒱=\left(\begin{array}{c}\mathrm{}(𝒰)\hfill \\ \mathrm{}(𝒰)\hfill \end{array}\right):L_2(G,g^{00}){\displaystyle \underset{\mathrm{}}{}}L_2(\sigma (𝒦),\mu _{\mathrm{}})\stackrel{\mathrm{def}}{=}L_2(\sigma ,𝒱),`$ (31)
which is $``$-real $`(𝒱\varphi )^{}=𝒱(\varphi ^{})`$ and does also diagonalise the operator $`𝒱𝒦𝒱^{}=`$.
Accordingly, the Hilbert space $`L_2(\sigma ,𝒱)`$ is restricted to real numbers and the associated scalar product reads
$`\left\{u\right|v\}_\sigma =\left\{𝒱^{}u\right|𝒱^{}v\}_1={\displaystyle \underset{\mathrm{}}{}}{\displaystyle \underset{\sigma (𝒦)}{}}𝑑\mu _{\mathrm{}}(\lambda )u_{\mathrm{}}(\lambda )v_{\mathrm{}}(\lambda )\stackrel{\mathrm{def}}{=}\begin{array}{c}\\ {\displaystyle }\\ \mathrm{\Gamma }\end{array}u_\mathrm{\Gamma }v_\mathrm{\Gamma }.`$ (35)
Because $`L_2(\sigma ,𝒱)`$ is a real Hilbert space over $``$, the usual complex conjugation of the first argument in the scalar product disappears.
For a discrete spectrum $`\sigma (𝒦)=\sigma _\mathrm{d}`$ the measure $`d\mu _{\mathrm{}}(\lambda )`$ denotes simply a sum and for an absolute continuous spectrum $`\sigma (𝒦)=\sigma _{\mathrm{ac}}`$ an elementary integral possibly together with a $`\mathrm{}`$-summation, cf. Eq. (35). For example, the $`\mathrm{}`$-sum may describe the angular quantum numbers for the Laplacian in spherical coordinates $`\mathrm{}=\mathrm{},m`$. Both, the $`\mathrm{}`$-summation and the integration with the measure $`d\mu _{\mathrm{}}(\lambda )`$ are now abbreviated by the index $`\mathrm{\Gamma }`$.
Performing the transformation of the fields $`\left|Q\right\}_\sigma =𝒱\left|\mathrm{\Phi }\right\}_1`$ the Lagrange function can be diagonalised
$`L={\displaystyle \frac{1}{2}}\left\{\dot{Q}\right|\dot{Q}\}_\sigma {\displaystyle \frac{1}{2}}\left\{Q\right||Q\}_\sigma ={\displaystyle \frac{1}{2}}\begin{array}{c}\\ {\displaystyle }\\ \mathrm{\Gamma }\end{array}\left(\dot{Q}_\mathrm{\Gamma }^2\omega _\mathrm{\Gamma }^2Q_\mathrm{\Gamma }^2\right),`$ (39)
with $`\omega _\mathrm{\Gamma }^2\stackrel{\mathrm{def}}{=}\lambda _\mathrm{\Gamma }\sigma (𝒦)_+`$ which will be called eigenfrequencies.
One should note that the $`\left|Q\right\}_\sigma `$ still depend on time $`\left|Q(t)\right\}_\sigma `$, only the spatial dependence is transformed by $`𝒱`$. Owing to the reality of the transformation $`𝒱`$ the amplitudes $`Q_\mathrm{\Gamma }(t)`$ are real as well.
### 3.5 Canonical quantisation
Starting with the diagonal Lagrange function in Eq. (39) we are able to perform the canonical quantisation procedure by imposing the usual equal time commutation relations
$`[\left\{u\right|\widehat{Q}(t)\}_\sigma ,\left\{\widehat{P}(t)\right|v\}_\sigma ]`$ $`=`$ $`i\left\{u\right|v\}_\sigma ,`$
$`[\left\{u\right|\widehat{Q}(t)\}_\sigma ,\left\{\widehat{Q}(t)\right|v\}_\sigma ]`$ $`=`$ $`[\left\{u\right|\widehat{P}(t)\}_\sigma ,\left\{\widehat{P}(t)\right|v\}_\sigma ]=0,`$ (40)
which hold for all $`\left|u\right\}_\sigma `$ and $`\left|v\right\}_\sigma `$. In this representation the canonical conjugated momenta are simply determined by $`\left|P\right\}_\sigma =|dQ/dt\}_\sigma `$.
Due to the isometry of the transformation $`𝒱`$ these commutation relations are completely equivalent to the corresponding relations for the field $`\widehat{\mathrm{\Phi }}`$. For a static metric the inner product is related to the scalar product via
$`\left(\psi |\varphi \right)=i\left\{\psi \right|\dot{\varphi }\}_1i\left\{\dot{\psi }\right|\varphi \}_1.`$ (41)
As a consequence, the relations above are indeed identical to the commutators of the fields
$`[\left(\psi |\widehat{\mathrm{\Phi }}\right),\left(\widehat{\mathrm{\Phi }}|\varphi \right)]=\left(\psi |\varphi \right).`$ (42)
The Hamiltonian splits up into an infinite set of commuting parts describing harmonic oscillators that are appropriate for a particle definition
$`\widehat{H}={\displaystyle \frac{1}{2}}\left\{\widehat{P}\right|\widehat{P}\}_\sigma +{\displaystyle \frac{1}{2}}\left\{\widehat{Q}\right||\widehat{Q}\}_\sigma ={\displaystyle \frac{1}{2}}\begin{array}{c}\\ {\displaystyle }\\ \mathrm{\Gamma }\end{array}\left(\widehat{P}_\mathrm{\Gamma }^2+\omega _\mathrm{\Gamma }^2\widehat{Q}_\mathrm{\Gamma }^2\right).`$ (46)
In terms of the creators $`\left|\widehat{A}^{}\right\}`$ and annihilators
$`\left|\widehat{A}\right\}={\displaystyle \frac{1}{\sqrt{2}}}\left(^{1/4}\right|\widehat{Q}(t=0)\}+i^{1/4}\left|\widehat{P}(t=0)\right\})`$ (47)
the Hamiltonian can be cast into the form
$`\widehat{H}={\displaystyle \frac{1}{2}}\left\{\widehat{A}^{}\right|^{1/2}|\widehat{A}\}_\sigma +{\displaystyle \frac{1}{2}}\left\{\widehat{A}\right|^{1/2}|\widehat{A}^{}\}_\sigma =\begin{array}{c}\\ {\displaystyle }\\ \mathrm{\Gamma }\end{array}{\displaystyle \frac{\omega _\mathrm{\Gamma }}{2}}\left(\widehat{A}_\mathrm{\Gamma }^{}\widehat{A}_\mathrm{\Gamma }+\widehat{A}_\mathrm{\Gamma }\widehat{A}_\mathrm{\Gamma }^{}\right).`$ (51)
For a discrete spectrum $`\sigma (𝒦)=\sigma _\mathrm{d}`$ this already defines the physical particles because we have now creation and annihilation operators $`\widehat{A}_\mathrm{\Gamma }^{}`$ and $`\widehat{A}_\mathrm{\Gamma }`$ that diagonalise the Hamiltonian (which is also the energy operator).
For a continuous spectrum $`\sigma (𝒦)=\sigma _\mathrm{c}`$ the quantities $`\widehat{P}_\mathrm{\Gamma }`$, $`\widehat{Q}_\mathrm{\Gamma }`$, $`\widehat{A}_\mathrm{\Gamma }`$ and $`\widehat{A}_\mathrm{\Gamma }^{}`$ are not well-defined operators but operator-valued distributions because
$`[\widehat{Q}_\mathrm{\Gamma }(t),\widehat{P}_\mathrm{\Lambda }(t)]=i\delta (\mathrm{\Gamma },\mathrm{\Lambda }),`$ (52)
which is a Dirac $`\delta `$-distribution for continuous indices $`\mathrm{\Gamma },\mathrm{\Lambda }`$. But a product of two distributions acting on the same linear space (e.g. the Schwartz/Sobolev space $`𝔖_1`$ for the $`\delta `$-distribution), i.e. with the same index $`\mathrm{\Gamma }`$, is not well-defined. This reflects the infinite-volume divergence in quantum field theory. Consequently, the Hamiltonian in Eqs. (46) and (51) is not well-defined. It may only be viewed as a formal expression until an appropriate regularisation method has been applied.
### 3.6 Vacuum definition
In order to get rid of the singularities discussed above and to obtain well-defined operators $`\widehat{a}_I`$ we introduce a complete orthonormal and real basis $`\left|b_I\right\}_\sigma `$ with $`I`$ of the separable Hilbert space $`L_2(\sigma ,𝒱)`$ and define
$`\widehat{a}_I\stackrel{\mathrm{def}}{=}\left\{\widehat{A}\right|b_I\}_\sigma .`$ (53)
For a discrete spectrum $`\sigma (𝒦)=\sigma _\mathrm{d}`$ we may choose $`b_I(\mathrm{\Gamma })=\delta _{\mathrm{\Gamma }I}`$ which leads us back to the operators $`\widehat{A}_\mathrm{\Gamma }`$. For a continuous spectrum $`\sigma (𝒦)=\sigma _\mathrm{c}`$ this coincidence does not hold. Due to $`\left\{b_I\right|b_J\}_\sigma =\delta _{IJ}`$ with a Kronecker-$`\delta _{IJ}`$ the $`\widehat{a}_I`$ are well-defined operators with $`[\widehat{a}_I,\widehat{a}_J^{}]=\delta _{IJ}`$ instead of operator-valued distributions with $`[\widehat{A}_\mathrm{\Gamma },\widehat{A}_\mathrm{\Lambda }^{}]=\delta (\mathrm{\Gamma },\mathrm{\Lambda })`$. Unfortunately, for a continuous spectrum $`\sigma (𝒦)=\sigma _\mathrm{c}`$ the operators $`\widehat{a}_I`$ are now well-defined, but do not exactly diagonalise the Hamiltonian. But – as we shall see later in Section 3.7 – one may choose an appropriate basis $`\left|b_I\right\}_\sigma `$ for which the operators $`\widehat{a}_I`$ approximately diagonalise the Hamiltonian.
The corresponding number operators take the usual form $`\widehat{n}_I\stackrel{\mathrm{def}}{=}\widehat{a}_I^{}\widehat{a}_I`$. The Fock space $`𝔉`$ which contains all pure states $`|\mathrm{\Psi }`$ of the quantum field $`\widehat{\mathrm{\Phi }}`$ is now defined as the completion of the linear hull of the proper eigenvectors of these commuting operators $`\widehat{n}_I`$ for all indices $`I`$
$`𝔉\stackrel{\mathrm{def}}{=}\overline{\stackrel{}{\mathrm{lin}}\left\{\right|\mathrm{\Psi }:_I:\widehat{n}_I|\mathrm{\Psi }=|\mathrm{\Psi }n_I\}}.`$ (54)
As a consequence, the spectrum of the operators $`\widehat{n}_I`$ in this Fock space is a pure point spectrum $`\sigma (\widehat{n}_I)=\sigma _\mathrm{p}`$. With the same arguments as already used for the quantisation of the harmonic oscillator the commutation relations $`[\widehat{a}_I,\widehat{a}_J^{}]=\delta _{IJ}`$ imply $`\sigma (\widehat{n}_I)=`$. It should be mentioned here that this definition of the Fock space $`𝔉`$ is slightly different to the frequently employed approach based on the one-particle Hilbert space $`L_2^{}(G,g^{00})`$ (see e.g. )
$`𝔉=\left(\right)_{\mathrm{symm}}\left(\right)_{\mathrm{symm}}\mathrm{}.`$ (55)
Nevertheless, these distinct definitions are related if one divides the Fock space in Eq. (54) into orthogonal subspaces labelled by different values of the total number of particles $`n_{\mathrm{total}}\stackrel{\mathrm{def}}{=}_In_I`$.
Accordingly, the vacuum is defined as the eigenvector of all commuting operators $`\widehat{n}_I`$ with eigenvalue zero
$`_I:\widehat{n}_I|0=0\mathrm{i}.\mathrm{e}.\widehat{a}_I|0=0.`$ (56)
This definition is independent of the special choice of the basis $`\left|b_I\right\}_\sigma `$. To prove this statement, we use the completeness of the basis $`\left|b_I\right\}_\sigma `$ to obtain the following result
$`\left|\zeta \right\}_\sigma :\{\widehat{A}|\zeta \}_\sigma |0=0.`$ (57)
If we regularise the formal expression for the Hamiltonian in Eq. (51) via insertion of a complete basis (in principle together with a convergence factor, $`\mathrm{exp}(I\epsilon )`$ for instance)
$`\widehat{H}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{I}{}}\left(\left\{^{1/4}\widehat{A}^{}\right|b_I\}_\sigma \left\{b_I\right|^{1/4}\widehat{A}\}_\sigma +\left\{^{1/4}\widehat{A}\right|b_I\}_\sigma \left\{b_I\right|^{1/4}\widehat{A}^{}\}_\sigma \right)`$ (58)
$`=`$ $`{\displaystyle \underset{I}{}}\left(\left\{\widehat{A}^{}\right|^{1/4}b_I\}_\sigma \left\{^{1/4}b_I\right|\widehat{A}\}_\sigma +{\displaystyle \frac{1}{2}}\left\{^{1/4}b_I\right|^{1/4}b_I\}_\sigma \right),`$
it appears as a divergent sum of some non-negative operators of the structure $`\widehat{X}_I^{}\widehat{X}_I`$ and remaining $``$-numbers. The above defined vacuum is the ground state of all operators $`\widehat{X}_I^{}\widehat{X}_I`$ and in this regard also the ground state of the Hamiltonian. Hence, the divergent amount of $``$-number terms represents the zero-point energy. The infinite summation over the index $`I`$ corresponds to the sum over arbitrary high frequencies and – for a continuous spectrum – the summation of an infinite number of basis elements for a given frequency interval. The first infinity, the infinite energy divergence, is always present in quantum field theory and the latter, the infinite volume divergence, only for non-discrete spectra.
In the Minkowski space-time the above defined vacuum coincides (of course) with the usual Minkowski vacuum $`|0=|0_\mathrm{M}`$. In the Schwarzschild space-time this state – which is the ground state of the Hamiltonian – is called the Boulware state $`|0=|\mathrm{\Psi }_\mathrm{B}`$.
The particle definition presented above can be reproduced utilising the well-known approach based on the inner product: The basis elements $`\left|b_I\right\}_\sigma `$ of the Hilbert space $`L_2(\sigma ,𝒱)`$ are normalised and therefore correspond to functions $`e_i(𝒓)`$ via $`\left|e_I\right\}_1=𝒱^{}\left|b_I\right\}_\sigma `$ which are also normalised $`\left\{e_I\right|e_J\}_1=\delta _{IJ}`$ and build up a basis of the Hilbert space $`L_2(G,g^{00})`$. As a consequence, the operators $`\widehat{a}_I`$ correspond to localised wave-packets $`\left|F_I(t)\right\}_1`$ which are defined as follows
$`\left|F_I(t)\right\}_1=\left(4𝒦\right)^{1/4}\mathrm{exp}(i𝒦^{1/2}t)\left|e_I\right\}_1.`$ (59)
These quantities are solutions of the Klein-Fock-Gordon equation $`_t^2\left|F_I\right\}_1=𝒦\left|F_I\right\}_1`$ and normalised with respect to the inner product
$`(F_I|F_J)=(F_I^{}|F_J^{})=\delta _{IJ},(F_I^{}|F_J)=(F_I|F_J^{})=0.`$ (60)
Comparison with the particle definition via the inner product verifies indeed the identification
$`\widehat{a}_I=\left(F_I|\widehat{\mathrm{\Phi }}\right).`$ (61)
The functions $`F_I`$ and $`F_I^{}`$ form a complete set of solutions of the Klein-Fock-Gordon equation. Hence, the field $`\widehat{\mathrm{\Phi }}`$ may be expanded via
$`\widehat{\mathrm{\Phi }}={\displaystyle \underset{I}{}}\widehat{a}_IF_I+\widehat{a}_I^{}F_I^{},`$ (62)
which demonstrates again the equivalence of the approaches.
### 3.7 Eigenfunctions
For a point spectrum $`\sigma _\mathrm{p}`$ there exist proper eigenfunctions $`f_\mathrm{\Gamma }L_2(G,g^{00})`$ with $`𝒦\left|f_\mathrm{\Gamma }\right\}_1=\omega _\mathrm{\Gamma }^2\left|f_\mathrm{\Gamma }\right\}_1`$, but for a continuous spectrum $`\sigma _\mathrm{c}`$ this is of course not the case. Nevertheless, it is in many cases possible to find an analogue. If pointwise defined functions $`f_\mathrm{\Gamma }(𝒓)`$ (or – more generally – locally integrable functions $`f_\mathrm{\Gamma }L_1^{\mathrm{local}}`$) exist such that
$`\left\{\zeta \right|𝒱\varphi \}_\sigma =\begin{array}{c}\\ {\displaystyle }\\ \mathrm{\Gamma }\end{array}\zeta _\mathrm{\Gamma }{\displaystyle \underset{G}{}}d^3𝒓f_\mathrm{\Gamma }(𝒓)\varphi (𝒓)`$ (66)
holds for all $`\varphi C_0^{\mathrm{}}(G)`$ and $`\zeta C_0^{\mathrm{}}(\sigma )`$ the functions $`f_\mathrm{\Gamma }(𝒓)`$ are called (generalised) eigenfunctions of the $`𝒦`$-operator. In contrast to the proper ($`\sigma _\mathrm{p}`$) eigenfunctions with $`f_\mathrm{\Gamma }\mathrm{Def}(𝒦)L_2(G,g^{00})`$ the generalised ($`\sigma _\mathrm{c}`$) eigenfunctions do not belong to the Hilbert space $`f_\mathrm{\Gamma }L_2(G,g^{00})`$ and (of course) also not to the domain of definition of the $`𝒦`$-operator $`f_\mathrm{\Gamma }\mathrm{Def}(𝒦)`$. However, due to $`𝒱𝒦=𝒱`$ also the generalised eigenfunctions fulfil the pointwise/local (generalised) eigenvalue equation $`𝒦_{\mathrm{local}}f_\mathrm{\Gamma }(𝒓)=\omega _\mathrm{\Gamma }^2f_\mathrm{\Gamma }(𝒓)`$. This is a very important relation for the calculation of these eigenfunctions. If the (generalised) eigenfunctions exist, the transformation of the fields $`\left|\mathrm{\Phi }\right\}_1=𝒱^{}\left|Q\right\}_\sigma `$ can be described by the pointwise/local identity
$`\widehat{\mathrm{\Phi }}(𝒓,t)=\begin{array}{c}\\ {\displaystyle }\\ \mathrm{\Gamma }\end{array}\widehat{Q}_\mathrm{\Gamma }(t)f_\mathrm{\Gamma }(𝒓).`$ (70)
Even though the generalised eigenfunctions are not in $`L_2(G,g^{00})`$, they may be thought as a (singular) limiting case of $`L_2(G,g^{00})`$-functions: In the following considerations we assume $`d\mu (\lambda )=d\lambda `$ and $`\sigma =`$ for reasons of simplicity. The $`L_2(\sigma ,𝒱)`$-basis functions $`b_I(\lambda )`$ can be squeezed and translated $`b_{I(\mathrm{\Gamma })}^\epsilon (\lambda )=b_I(\lambda /\epsilon \lambda _\mathrm{\Gamma }/\epsilon )/\sqrt{\epsilon }`$ and are still a basis of $`L_2(\sigma ,𝒱)`$. Evaluating the (singular) limiting case of these squeezed basis functions
$`\underset{\epsilon 0}{lim}{\displaystyle \frac{b_{I(\mathrm{\Gamma })}^\epsilon (\lambda )}{\sqrt{\epsilon }}}=\underset{\epsilon 0}{lim}{\displaystyle \frac{b_I(\lambda \epsilon ^1\lambda _\mathrm{\Gamma }\epsilon ^1)}{\epsilon }}=𝒩_{I(\mathrm{\Gamma })}\delta (\lambda \lambda _\mathrm{\Gamma }),`$ (71)
where $`𝒩_{I(\mathrm{\Gamma })}`$ denotes some normalisation factor, we observe that every generalised eigenfunction $`f_\mathrm{\Gamma }(𝒓)`$ can be locally approximated by appropriately chosen wave packets $`\left|e_{I(\mathrm{\Gamma })}^\epsilon \right\}_1=𝒱^{}\left|b_{I(\mathrm{\Gamma })}^\epsilon \right\}_\sigma `$
$`\underset{\epsilon 0}{lim}{\displaystyle \frac{e_{I(\mathrm{\Gamma })}^\epsilon (𝒓)}{\sqrt{\epsilon }}}=𝒩_{I(\mathrm{\Gamma })}f_\mathrm{\Gamma }(𝒓).`$ (72)
Accordingly, also the operator-valued distributions $`\widehat{A}_\mathrm{\Gamma }`$ may be considered as a singular limiting case of the regular operators $`\widehat{a}_{I(\mathrm{\Gamma })}^\epsilon `$
$`\underset{\epsilon 0}{lim}{\displaystyle \frac{\widehat{a}_{I(\mathrm{\Gamma })}^\epsilon }{\sqrt{\epsilon }}}=𝒩_{I(\mathrm{\Gamma })}\widehat{A}_\mathrm{\Gamma }.`$ (73)
The divergent factor $`1/\sqrt{\epsilon }`$ indicates the singular character of the generalised eigenfunctions (e.g. plane waves) in contrast to the regular basis elements (wave packets). Of course, in realistic experiments one never deals with plane waves, but wave packets. On the other hand, the calculations with plane waves are usually much simpler. Hence, in the following we shall perform our evaluations with eigenfunctions always bearing in mind their character as a singular limiting case of regular objects.
### 3.8 Continuum normalisation
To investigate the physical consequences caused by the singular behaviour of the product of two distributions $`\widehat{N}_\mathrm{\Gamma }=\widehat{A}_\mathrm{\Gamma }^{}\widehat{A}_\mathrm{\Gamma }`$ – expressed by the factor $`1/\epsilon `$ – we consider a quantum field confined in a finite volume $`V`$ and study the limiting case $`V\mathrm{}`$. This limit may be interpreted as the transition from a discrete spectrum $`\sigma (𝒦)=\sigma _\mathrm{d}`$ to a continuous one $`\sigma (𝒦)=\sigma _\mathrm{c}`$. For a 3-dimensional cubic volume $`V`$ the indices $`\mathrm{\Gamma }`$ correspond, for example, to discrete wave-numbers $`𝒌`$. In the continuum limit $`V\mathrm{}`$ the $`𝒌`$-sum transforms into an integral over $`d^3𝒌`$ via
$`{\displaystyle \underset{𝒌}{}}𝒩_VV{\displaystyle d^3𝒌}.`$ (74)
$`𝒩_V`$ denotes a normalisation factor which depends on the imposed boundary conditions (Dirichlet, Neumann, periodic, etc.) and the shape of the domain $`G`$. The Kronecker-$`\delta _{𝒌,𝒌^{}}`$ converts into a Dirac-$`\delta ^3(𝒌𝒌^{})`$ in an analogue way $`𝒩_VV\delta _{𝒌,𝒌^{}}\delta ^3(𝒌𝒌^{})`$. Ergo, the singularity $`\delta ^3(𝒌𝒌)`$ displays the infinite-volume divergence $`\delta (𝒌,𝒌)=\delta ^3(𝒌𝒌)=𝒩_VV`$. Recalling the formal expression for the Hamiltonian in Eq. (51)
$`\widehat{H}=\begin{array}{c}\\ {\displaystyle }\\ \mathrm{\Gamma }\end{array}\omega _\mathrm{\Gamma }\left(\widehat{A}_\mathrm{\Gamma }^{}\widehat{A}_\mathrm{\Gamma }+{\displaystyle \frac{1}{2}}\delta (\mathrm{\Gamma },\mathrm{\Gamma })\right)=\begin{array}{c}\\ {\displaystyle }\\ \mathrm{\Gamma }\end{array}\left(\omega _\mathrm{\Gamma }\widehat{N}_\mathrm{\Gamma }+{\displaystyle \frac{\omega _\mathrm{\Gamma }}{2}}\delta (\mathrm{\Gamma },\mathrm{\Gamma })\right),`$ (81)
we observe that – in addition to the mode summation/integration – its indefinite character exactly exhibits this divergence. Indeed, if one examines the continuum limit of the Hamiltonian
$`\widehat{H}^\mathrm{d}={\displaystyle \underset{𝒌}{}}\left(\widehat{N}_𝒌^\mathrm{d}+{\displaystyle \frac{1}{2}}\right)|𝒌|\widehat{H}^\mathrm{c}={\displaystyle d^3𝒌\left(\widehat{N}_𝒌^\mathrm{c}+\frac{1}{2}\delta ^3(𝒌𝒌)\right)|𝒌|},`$ (82)
the limiting number ”operator” $`\widehat{N}_𝒌^\mathrm{c}`$ can be identified via $`𝒩_VV\widehat{N}_𝒌^\mathrm{d}\widehat{N}_𝒌^\mathrm{c}`$. The singular character of this formal expression may be exemplified with the following consideration: If the state of the quantum field corresponds to thermal equilibrium at a temperature $`T>0`$, the expectation value of the number operator $`\widehat{N}_𝒌^\mathrm{d}`$ equals the Bose-Einstein distribution for arbitrary large but finite volumes $`V`$. For an infinite volume the expectation value of the quantity $`\widehat{N}_𝒌^\mathrm{c}`$ diverges owing to the factor $`𝒩_VV`$. The expectation values of the regular operators $`\widehat{n}_I`$ are (of course) still finite and behave as the Bose-Einstein distribution evaluated at some averaged frequency $`\omega _I`$.
### 3.9 Bogoliubov coefficients
So far we have considered static space-times and developed an appropriate particle definition. If we now drop the restriction to stationary metrics and take dynamical space-times into account, the question concerning particle creation arises. A variation of the metric $`g_{\mu \nu }`$ induces a change of the $`𝒦`$-operator and – possibly – the corresponding Hilbert space $`L_2(G,g^{00})`$. A function, which belongs initially to $`L_2(G^{\mathrm{in}},g_{\mathrm{in}}^{00})`$ may be later (e.g. if a horizon has formed) not in $`L_2(G^{\mathrm{out}},g_{\mathrm{out}}^{00})`$ but a distribution with respect to the $`L_2(G^{\mathrm{out}},g_{\mathrm{out}}^{00})`$-scalar product. As a consequence, it is not clear whether the Bogoliubov coefficient (see e.g. ) describing the particle creation
$`\beta _{IJ}=\left(F_I^{}|F_J\right)`$ (83)
exists for all $`F_IL_2(G^{\mathrm{in}},g_{\mathrm{in}}^{00})`$ and $`F_JL_2(G^{\mathrm{out}},g_{\mathrm{out}}^{00})`$ or not. However, $`C_0^{\mathrm{}}`$-functions belong to the domain of definition of all (tempered) distributions. Thus for $`F_I^{\mathrm{in}}C_0^{\mathrm{}}(G^{\mathrm{in}})`$ and $`F_J^{\mathrm{out}}C_0^{\mathrm{}}(G^{\mathrm{out}})`$ the Bogoliubov coefficients always exist provided the metric can be cast into an analytic form. Similar to the previous Sections all other quantities (e.g. generalised eigensolutions) have to be approximated with $`C_0^{\mathrm{}}`$-functions.
## 4 Flat space-time examples
In the previous Section we have derived a canonical definition of particles for curved space-times which fulfil certain conditions. In the following we are going to apply this approach to the most simple example of a flat space-time in order to achieve a deeper insight into the physical consequences of the used mathematical theorems.
For the unbounded 1+1 dimensional Minkowski space-time with $`ds^2=dt^2dx^2`$ the $`𝒦`$-operator reads
$`𝒦={\displaystyle \frac{^2}{x^2}},`$ (84)
together with the domain $`G=(\mathrm{}<x<\mathrm{})`$. The infinite volume of this domain and the regularity of the metric cause a purely absolute continuous spectrum $`\sigma (𝒦)=\sigma _{\mathrm{ac}}=_+`$. The unitary transformation $`𝒰`$ (see Section 3.4) is simply the one-dimensional Fourier transformation $`𝒰=`$. The quasi-unitary transformation $`𝒱`$ takes the real and imaginary parts separately leading to the generalised eigenfunctions $`\mathrm{sin}\omega x`$ and $`\mathrm{cos}\omega x`$. Hence, the expansion of the field $`\widehat{\mathrm{\Phi }}`$ takes the following form
$`\widehat{\mathrm{\Phi }}(x,t)`$ $`=`$ $`\begin{array}{c}\\ {\displaystyle }\\ \omega \end{array}\widehat{Q}_{\omega ,\mathrm{c}}(t)\mathrm{cos}(\omega x)+\widehat{Q}_{\omega ,\mathrm{s}}(t)\mathrm{sin}(\omega x)`$ (88)
$`=`$ $`𝒩_\mu {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d\omega }{\sqrt{2\omega }}}(\widehat{A}_{\omega ,\mathrm{c}}^{}e^{i\omega t}\mathrm{cos}(\omega x)+\widehat{A}_{\omega ,\mathrm{s}}^{}e^{i\omega t}\mathrm{sin}(\omega x)+\mathrm{h}.\mathrm{c}.),`$ (89)
with a normalisation factor $`𝒩_\mu `$ depending on the explicit form of the measure $`d\mu (\omega ^2)`$, e.g. $`d\mu (\omega ^2)=d\omega `$ or $`d\mu (\omega ^2)=d\omega /2\pi `$ etc. The spectrum of the $`𝒦`$-operator discussed above is twice degenerated, i.e. there are two independent generalised eigenfunctions ($`\mathrm{sin}\omega x`$ and $`\mathrm{cos}\omega x`$) for every point $`\lambda =\omega ^2`$ of the spectrum $`\sigma `$. This degeneracy of the spectrum allows for the definition of particles with a definite direction of propagation: With a simple linear transformation we may rearrange the expansion of the field
$`\widehat{\mathrm{\Phi }}(x,t)=𝒩_\mu {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d\omega }{\sqrt{4\omega }}}(\widehat{A}_{\omega ,+}^{}e^{i\omega t+i\omega x}+\widehat{A}_{\omega ,}^{}e^{i\omega ti\omega x}+\mathrm{h}.\mathrm{c}.).`$ (90)
The new introduced quantities $`\widehat{A}_{\omega ,+}`$, $`\widehat{A}_{\omega ,}`$, $`\widehat{A}_{\omega ,+}^{}`$ and $`\widehat{A}_{\omega ,}^{}`$ obey the same commutation relations as the original ones $`\widehat{A}_{\omega ,\mathrm{c}}`$, $`\widehat{A}_{\omega ,\mathrm{s}}`$, $`\widehat{A}_{\omega ,\mathrm{c}}^{}`$ and $`\widehat{A}_{\omega ,\mathrm{s}}^{}`$. Thus they also describe particles. In contrast to the original particles which correspond to standing waves with different phases ($`\mathrm{sin}\omega x`$ and $`\mathrm{cos}\omega x`$) the new particles describe left-moving and right-moving waves according to $`\mathrm{exp}(\pm i\omega x)`$. These complex functions are suitable for a definition of particles $`\widehat{A}_\mathrm{\Gamma }`$ but do not correspond to Hermitian amplitudes $`\widehat{Q}_\mathrm{\Gamma }`$.
As a second example we study the situation of a bounded domain $`G=(0<x<\mathrm{})`$ in a 1+1 dimensional Minkowski space-time with a Dirichlet boundary condition (a mirror) at $`x=0`$. Even though $`𝒦`$ seems to have the same form as in Eq. (84) it denotes a different operator as a result of the boundary condition. With the same arguments the spectrum is purely absolute continuous $`\sigma =\sigma _{\mathrm{ac}}=_+`$. In contrast to the previous example this spectrum is not degenerated. Every point $`\lambda =\omega ^2`$ of $`\sigma `$ corresponds to exactly one generalised eigenfunction, i.e. $`\mathrm{sin}\omega x`$. As a consequence, the definition of particles with a certain direction (left-moving or right-moving) is not possible. This result is physical reasonable if one takes conservation laws into account. Every left-moving component will be reflected by the mirror at $`x=0`$ after some period of time and turns its direction into right-moving and vice versa.
A finite domain $`G=(0<x<L)`$ in a 1+1 dimensional Minkowski space-time with Dirichlet boundary conditions at $`x=0`$ and $`x=L`$ of course possesses a purely discrete spectrum $`\sigma =\sigma _\mathrm{d}`$ with proper eigenfunctions $`\mathrm{sin}(\pi x/L)`$. The insertion of mirrors represented by Dirichlet boundary conditions usually lowers the ”density” of the spectrum $`\sigma `$, i.e. the number of eigenfunctions.
### 4.1 Ingoing and outgoing particles
Now we shall extend our investigations to the 3+1 dimensional Minkowski space-time described by different coordinate systems. Using spherical coordinates $`r,\vartheta ,\phi `$ it will turn out that the definition of ingoing or outgoing particles is not possible within the canonical approach. This is a consequence of the spectral properties of the operator
$`𝒦=^2={\displaystyle \frac{^2}{𝒓^2}},`$ (91)
together with the domain $`G=^3`$. Expressed by Cartesian coordinates $`𝒓=(x,y,z)^T`$ the generalised eigenfunctions take the simple form $`\mathrm{sin}(𝒌𝒓)`$ and $`\mathrm{cos}(𝒌𝒓)`$ with $`|𝒌|=\omega `$. As it is well-known these functions form a complete basis of $`L_2(^3)`$.
Employing spherical coordinates $`r,\vartheta ,\phi `$ the Cartesian eigenfunctions can be expanded with the aid of the equality
$`\mathrm{exp}(i𝒌𝒓)={\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}i^{\mathrm{}}(2\mathrm{}+1)j_{\mathrm{}}(\omega r)P_{\mathrm{}}(\mathrm{cos}\theta ),`$ (92)
with $`𝒌𝒓=\omega r\mathrm{cos}\theta `$ and the Legendre polynomials $`P_{\mathrm{}}`$. By inspection we recognise the following fact: For a given angular behaviour the spherical Bessel functions $`j_{\mathrm{}}(\omega r)`$ are already complete to describe the radial dependence. The Neumann $`n_{\mathrm{}}(\omega r)`$ or Hankel functions $`h_{\mathrm{}}^\pm (\omega r)`$ are not required and would be ”over-complete”. Therefore they do not describe additional degrees of freedom and do not enter the particle definition. This result can also be derived by considering the spectrum of the $`𝒦`$-operator. Due to the singular behaviour of the functions $`n_{\mathrm{}}(\omega r)`$ and $`h_{\mathrm{}}^\pm (\omega r)`$ at $`r=0`$ they are not eigenfunctions.
To acquire real eigenfunctions we have to introduce redefined spherical harmonics $`𝒴_\mathrm{}m(\vartheta ,\phi )\stackrel{\mathrm{def}}{=}𝒩_{\mathrm{}}P_{\mathrm{}}^m(\mathrm{cos}\vartheta )\mathrm{cos}(m\phi )`$ for $`m0`$ and $`𝒴_\mathrm{}m(\vartheta ,\phi )\stackrel{\mathrm{def}}{=}𝒩_{\mathrm{}}P_{\mathrm{}}^m(\mathrm{cos}\vartheta )\mathrm{sin}(m\phi )`$ for $`m<0`$. $`P_{\mathrm{}}^m`$ denote the associated Legendre polynomials and $`𝒩_{\mathrm{}}`$ are normalisation factors. Accordingly, the complete set of real and orthogonal eigenfunctions reads
$`f_{\omega \mathrm{}m}(r,\vartheta ,\phi )=𝒩_\omega \mathrm{}j_{\mathrm{}}(\omega r)𝒴_\mathrm{}m(\vartheta ,\phi ).`$ (93)
Again we observe the occurrence of exactly one eigenfunction $`f_{\omega \mathrm{}m}`$ per eigenfrequency $`\omega `$ for a fixed angular dependence $`\mathrm{},m`$. The regularity at $`r=0`$ plays the role of an effective boundary condition and forbids the existence of additional eigenfunctions such as $`n_{\mathrm{}}(\omega r)`$ or $`h_{\mathrm{}}^\pm (\omega r)`$. As a consequence, within the canonical approach it is not possible to define radial ingoing or outgoing particles in the Minkowski space-time. Functions like $`\mathrm{exp}(\pm i\omega r)/r`$ are not eigenfunctions of the Laplacian and therefore not solutions of the wave equation
$`\mathrm{}{\displaystyle \frac{\mathrm{exp}(i\omega t\pm i\omega r)}{r}}=4\pi e^{i\omega t}\delta ^3(𝒓)0.`$ (94)
Expanding the field $`\widehat{\mathrm{\Phi }}`$ into functions that do not satisfy the equation of motion $`\mathrm{}\widehat{\mathrm{\Phi }}=0`$ would abandon the independence of the distinct particles. Functions like $`\mathrm{exp}(\pm i\omega r)/r`$ correspond to the resolvents $`(\omega ^2\pm i\epsilon )`$ of the operator $`𝒦`$ (remember $`\sigma (𝒦)`$). Particles are defined with respect to the eigenfunctions which are representations of the spectral family $``$ of the operator $`𝒦`$. Into this spectral family $``$ the resolvents themselves do not enter, but linear combinations of them: $`(\lambda i\epsilon )(\lambda +i\epsilon )(\lambda )`$ which again leads to $`\mathrm{sin}(\omega r)/r=\omega j_{\mathrm{}=0}(\omega r)`$.
The impossibility of defining radial ingoing and outgoing particles is not restricted to the Minkowski space-time, this holds also for arbitrary spherically symmetric metrics
$`ds^2=g_{00}(r)dt^2+g_{11}(r)dr^2+r^2d\mathrm{\Omega }^2,`$ (95)
provided that the coefficients of the metric $`g_{00}`$ and $`g_{11}`$ are smooth and analytic functions. Such functions can be Taylor expanded
$`g_{00}(r)=g_{00}(0)+g_{00}^{\prime \prime }(0){\displaystyle \frac{r^2}{2}}+𝒪(r^3),`$ (96)
where $`g_{00}^{}(0)`$ and $`g_{11}^{}(0)`$ have to vanish for smoothness. After the separation of the angular variables with $`𝒴_\mathrm{}m(\vartheta ,\phi )`$ the radial dependence of the eigenfunctions is governed by a second-order ordinary differential equation in $`r`$. Provided its coefficients are smooth and regular the solutions of such an equation are uniquely determined by the first two (non-vanishing) terms of their Laurent expansion. For the evaluation of these initial data only the terms $`g_{00}(0)`$, $`g_{00}^{}(0)=0`$, $`g_{11}(0)`$ and $`g_{11}^{}(0)=0`$ are of relevance. For that purpose the radial part of the metric can be approximated by $`ds^2=g_{00}(0)dt^2+g_{11}(0)dr^2`$. Ergo the behaviour of the corresponding eigenfunctions is (up to a simple scale transformation with $`g_{00}(0)`$ and $`g_{11}(0)`$ respectively) asymptotically ($`r0`$) the same as in the Minkowski space-time. Consequently, also in these more general spherically symmetric metrics there exists exactly one eigenfunction for given $`\omega ,\mathrm{},m`$ which forbids the definition of radial ingoing and outgoing particles.
In view of conservation law arguments the nonexistence of ingoing and outgoing particles in regular space-times appears very plausible: Every ingoing component will bounce off at the origin after some period of time and eventually turn into outgoing and vice versa.
### 4.2 Rindler metric
As stated in Section 2.1 the particle interpretation crucially depends on the selection of a particular time-like Killing vector. In the following we shall consider an example where this dependence will become more evident. In the previous treatments we focused on the Killing vector mediating the Minkowski time translation symmetry. Of course this Killing field corresponds to usual observers at rest. But there exist further time-like Killing vectors in the Minkowski space-time – associated with special Lorentz boosts – which result in a deviating particle interpretation.
Starting with the 1+1 dimensional Minkowski metric $`ds^2=dt^2dx^2`$ and performing the coordinate transformation
$`t`$ $`=`$ $`\rho \mathrm{sinh}\kappa \tau ,`$
$`x`$ $`=`$ $`\rho \mathrm{cosh}\kappa \tau ,`$ (97)
one arrives at the Rindler metric $`ds^2=\kappa ^2\rho ^2d\tau ^2d\rho ^2`$. The quantity $`\kappa `$ is called the surface gravity, see e.g. . For fixed $`\rho `$ the transformation describes an accelerated motion. With respect to the new time coordinate $`\tau `$ the Rindler metric is static and thus allows for a particle definition according to Section 3. These particles may be interpreted as those seen by an accelerated observer. The corresponding $`𝒦`$-operator can be cast into the form
$`𝒦=\kappa ^2\rho {\displaystyle \frac{}{\rho }}\rho {\displaystyle \frac{}{\rho }},`$ (98)
with $`G=(0<\rho <\mathrm{})`$. The surface term (see Sec. 2) at $`\rho =0`$ vanishes without imposing any condition on the field $`\mathrm{\Phi }`$ due to $`\sqrt{g}=\kappa \rho =0`$ at $`\rho =0`$. Indeed, the Rindler metric possesses a horizon there. Since the occurrence of this horizon depends on the choice of the coordinates and thereby on the observer, it is a particle horizon (with respect to all world-lines $`\rho =\mathrm{const}`$) but not an event (or apparent, etc.) horizon, see Sec. 2.1 and .
For further investigations it is convenient to introduce the tortoise coordinate $`\rho _{}=\mathrm{ln}(\kappa \rho )/\kappa `$. In terms of this coordinate the metric reads $`ds^2=e^{2\kappa \rho _{}}\left(d\tau ^2d\rho _{}^2\right)`$ resulting in the operator
$`𝒦={\displaystyle \frac{^2}{\rho _{}^2}},`$ (99)
with $`G=(\mathrm{}<\rho _{}<\mathrm{})`$. As a consequence, the spectrum is twice degenerated and the corresponding eigenfunctions read $`\mathrm{sin}(\omega \rho _{})`$ and $`\mathrm{cos}(\omega \rho _{})`$ or $`\mathrm{exp}(\pm i\omega \rho _{})`$, respectively. Returning to the coordinate $`\rho `$ the eigenfunctions behave as $`\mathrm{exp}(\pm i\omega \mathrm{ln}(\kappa \rho )/\kappa )=(\kappa \rho )^{\pm i\omega /\kappa }`$. Even though the domain is bounded $`G=(0<\rho <\mathrm{})`$ there are two eigenfunctions per eigenvalue which allows for the definition of left-moving and right-moving particles. This indicates the absence of real boundary conditions on the field $`\mathrm{\Phi }`$ at the horizon $`\rho =0`$. In this regard the horizon is the opposite of a mirror. Even for a finite domain $`G=(0<\rho <L)`$ the spectrum of the $`𝒦`$-operator is still continuous due to the horizon: $`G=(\mathrm{}<\rho _{}<L_{})`$.
### 4.3 Unruh effect
After having performed a particle definition for the Minkowski and the Rindler observer, the question about the relationship of these two approaches arises. Evaluating the expectation value of the number of Rindler particles in the Minkowski vacuum one obtains a thermal distribution function, a consequence of the Unruh effect. This effect demonstrates manifestly that different observers may obey distinct particle interpretations. In consequence the vacuum may depend on the particular Killing vector.
One way (see also ) to calculate the expectation values explicitly is based on the Bogoliubov coefficients
$`\beta _{\mathrm{\Gamma }\mathrm{\Lambda }}=i{\displaystyle 𝑑\mathrm{\Sigma }^\mu F_\mathrm{\Gamma }^\mathrm{M}}\underset{\mu }{\overset{}{}}F_\mathrm{\Lambda }^\mathrm{R}.`$ (100)
The generalised Minkowski eigenfunctions are labelled by $`\mathrm{\Gamma }=(\xi ,\omega )`$
$`F_\mathrm{\Gamma }^\mathrm{M}(\underset{¯}{x})=F_{(\xi ,\omega )}^\mathrm{M}(x,t)=𝒩_\mathrm{M}{\displaystyle \frac{\mathrm{exp}(i\omega t)}{\sqrt{\omega }}}e^{i\xi \omega x},`$ (101)
where $`\xi =\pm 1`$ distinguishes the left-moving and right-moving particles. ($`𝒩_\mathrm{M}`$ denotes a normalisation factor.) In analogy the generalised Rindler eigenfunctions read
$`F_\mathrm{\Lambda }^\mathrm{R}(\underset{¯}{x})=F_{(\xi ^{},\omega ^{})}^\mathrm{R}(\rho ,\tau )=𝒩_\mathrm{R}{\displaystyle \frac{\mathrm{exp}(i\omega ^{}\tau )}{\sqrt{\omega ^{}}}}(\kappa \rho )^{i\xi ^{}\omega ^{}/\kappa }.`$ (102)
With the choice for the surface $`\mathrm{\Sigma }=\{\tau =0,0<\rho <\mathrm{}\}`$ the surface element takes the form $`d\mathrm{\Sigma }^0=g^{00}\sqrt{g}d\rho =d\rho /(\kappa \rho )`$. At this surface the Minkowski coordinates are simply given by $`t=0`$, $`x=\rho `$ and the derivative transforms according to $`_\tau =\kappa \rho _t`$. Putting all this together, the $`\beta `$-coefficient transforms into
$`\beta _{(\xi ,\omega );(\xi ^{},\omega ^{})}=𝒩_{\mathrm{MR}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d\rho }{\kappa \rho }}{\displaystyle \frac{\omega ^{}\kappa \rho \omega }{\sqrt{\omega \omega ^{}}}}e^{i\xi \omega \rho }(\kappa \rho )^{i\xi ^{}\omega ^{}/\kappa }.`$ (103)
This integral involves generalised eigenfunctions (corresponding to $`\widehat{A}_\mathrm{\Gamma }`$) and has to be understood in a distributional sense. For well-defined expressions (such as $`\widehat{a}_I`$) we have to insert a convergence factor, for instance $`(\kappa \rho )^\epsilon \mathrm{exp}(\epsilon \kappa \rho )`$. The existence of the limit $`\epsilon 0`$ confirms the possibility of approximating the singular eigenfunctions by regular quantities. After this procedure we may make use of the formula
$`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑ue^{uw}u^{z1}=w^z\mathrm{\Gamma }(z),`$ (104)
which holds for $`\mathrm{}(w)>0`$ and $`\mathrm{}(z)>0`$, and we arrive at
$`\beta _{(\xi ,\omega );(\xi ^{},\omega ^{})}=𝒩_{\mathrm{MR}}{\displaystyle \frac{1+\xi ^{}\xi }{\kappa }}\sqrt{{\displaystyle \frac{\omega ^{}}{\omega }}}\mathrm{\Gamma }(i\xi ^{}\omega ^{}/\kappa )\left(i\xi \omega /\kappa +\epsilon \right)^{i\xi ^{}\omega ^{}/\kappa }.`$ (105)
Calculating the remaining Bogoliubov coefficient $`\alpha _{(\xi ,\omega );(\xi ^{},\omega ^{})}`$ one gets nearly the same expression but with a positive sign in front of the term $`i\xi \omega /\kappa `$. Therefore both coefficients merely contribute for particles moving in the same ”direction” $`\beta _{(\xi ,\omega );(\xi ^{},\omega ^{})}\delta _{\xi ,\xi ^{}}`$, respectively, $`\alpha _{(\xi ,\omega );(\xi ^{},\omega ^{})}\delta _{\xi ,\xi ^{}}`$. Now it is possible to compare both quantities. As said before, the only difference between $`\alpha `$ and $`\beta `$ is the sign in front of the term $`\omega /\kappa `$. Dividing the two coefficients all other terms cancel and the convergence factor $`\epsilon `$ determines the side of the branch cut of the logarithm in the complex plane. Hence we find
$`\beta _{(\xi ^{},\omega );(\xi ^{},\omega ^{})}=\mathrm{exp}(\pi \omega ^{}/\kappa )\alpha _{(\xi ^{},\omega );(\xi ^{},\omega ^{})}.`$ (106)
An alternative way to obtain this important result is based on the analytic continuation into the complex plane. For that purpose we define slightly modified Bogoliubov coefficients via
$`\beta _{\xi ,\xi ^{}}^\mathrm{c}(\omega ,\omega ^{})=\sqrt{\omega \omega ^{}}\beta _{(\xi ,\omega );(\xi ^{},\omega ^{})},`$ (107)
and in analogy the $`\alpha `$-coefficient. In view of Eqs. (100)–(103) the modified Bogoliubov coefficients can be analytically continued to the complex $`\omega `$-plane where the relation
$`\alpha _{\xi ,\xi ^{}}^\mathrm{c}(\omega ,\omega ^{})=\beta _{\xi ,\xi ^{}}^\mathrm{c}(\omega ,\omega ^{})`$ (108)
holds. Inserting this equality into Eq. (105) reproduces Eq. (106). In order to evaluate the absolute value squared of the $`\beta `$-coefficient we may utilise the identity
$`\mathrm{\Gamma }(z)\mathrm{\Gamma }(z)={\displaystyle \frac{\pi }{z\mathrm{sin}\pi z}}`$ (109)
to obtain the final result
$`\left|\beta _{(\xi ,\omega );(\xi ^{},\omega ^{})}\right|^2={\displaystyle \frac{8\pi 𝒩_{\mathrm{MR}}^2}{\kappa \omega }}{\displaystyle \frac{\delta _{\xi ,\xi ^{}}}{\mathrm{exp}(2\pi \omega ^{}/\kappa )1}}.`$ (110)
In view of the remaining $`\omega `$-integration the number of Rindler particles in the Minkowski vacuum diverges. This result can also be re-derived using the well-known unitarity relation
$`\begin{array}{c}\\ {\displaystyle }\\ \mathrm{\Gamma }\end{array}\alpha _{\mathrm{\Lambda }\mathrm{\Gamma }}\alpha _{\mathrm{{\rm Y}}\mathrm{\Gamma }}^{}\beta _{\mathrm{\Lambda }\mathrm{\Gamma }}\beta _{\mathrm{{\rm Y}}\mathrm{\Gamma }}^{}=\delta (\mathrm{\Lambda },\mathrm{{\rm Y}}),`$ (114)
where $`\mathrm{\Gamma }`$ symbolises the Minkowski index. This equality reflects the completeness of the Minkowski solutions. Special care is required concerning the derivation of an analogue expression involving the Rindler functions since these solutions are restricted to the Rindler wedge and thereby they are not complete in the full Minkowski space-time.
Inserting Eq. (106) and considering the singular coincidence $`\mathrm{\Lambda }=\mathrm{{\rm Y}}`$ it follows
$`N_{(\xi ^{},\omega ^{})}`$ $`=`$ $`0_\mathrm{M}\left|\widehat{N}_{(\xi ^{},\omega ^{})}^\mathrm{R}\right|0_\mathrm{M}=\begin{array}{c}\\ {\displaystyle }\\ (\xi ,\omega )\end{array}\left|\beta _{(\xi ,\omega );(\xi ^{},\omega ^{})}\right|^2`$ (118)
$`=`$ $`{\displaystyle \frac{\delta (\omega ,\omega )}{\mathrm{exp}(2\pi \omega ^{}/\kappa )1}}={\displaystyle \frac{𝒩_VV}{\mathrm{exp}(2\pi \omega ^{}/\kappa )1}}.`$ (119)
In the last step we have used the results of Section 3.8. Recalling the argumentation made there we come to the conclusion that the divergence of $`N_\mathrm{\Gamma }`$ is necessary for a thermal behaviour.
The same calculation can be performed with well-defined operators $`\widehat{n}_I`$ corresponding to localised wave packets. For an appropriately chosen basis $`e_I(𝒓)`$ the coefficients $`\beta _{IJ}`$ are up to normalisation factors approximately the same as the $`\beta _{\mathrm{\Gamma }\mathrm{\Lambda }}`$ evaluated above. But in this case the results for $`\widehat{n}_I=_J|\beta _{IJ}|^2`$ are finite owing to $`\delta _{II}=1`$. Another explanation is the fact, that the $`\beta _{IJ}`$ for arbitrary frequencies $`\omega `$ do not coincide with the $`\beta _{\mathrm{\Gamma }\mathrm{\Lambda }}`$ (due to the localised character of the wave packets) which makes the $`\omega _I`$-summation finite. One way to perform technically such a calculation involving localised quantities is to insert a convergence factor with a finite $`\epsilon `$ similar to the comment after Eq. (103). Accordingly this finite $`\epsilon `$ enters the $`\beta `$-coefficients and causes a finite result of the $`\omega _I`$-summation and thereby a finite number of created particles as well. Omitting the corresponding normalisation factor the infinite volume divergence can be restored in the limit $`\epsilon 0`$.
### 4.4 KMS condition
From a strictly axiomatic point of view the divergent result in Eq. (118) in the last section may not be completely convincing. However, it is possible to show more rigorously that the Minkowski vacuum indeed behaves as a thermal state when analysed by a Rindler observer. This can be done by employing the Kubo-Martin-Schwinger (KMS) condition . A KMS state $`_T`$ is defined as a time-translationally invariant state which satisfies the following condition
$`\widehat{U}(t)\widehat{V}(t^{})_T=\widehat{V}(t^{})\widehat{U}(t+i/T)_T`$ (120)
for all observables $`\widehat{U}`$ and $`\widehat{V}`$ and some temperature $`T`$. It can be shown that if the (irreducible) algebra of observables possesses a well-defined matrix-representation then the KMS state corresponds to the usual canonical ensemble
$`\widehat{U}_T=\mathrm{Tr}\left\{\widehat{U}{\displaystyle \frac{\mathrm{exp}(\widehat{H}/T)}{Z}}\right\}.`$ (121)
One might wonder at the fact that the Minkowski vacuum, i.e. a pure state, displays thermal features – usually connected with mixed states. This can be explained by the thermo-field formalism, see e.g. . I.e., a pure state of a quantum system transforms into a mixed state after averaging over a subsystem owing to the correlations between the different subsystems. As a result of the particle horizon at $`\rho =0`$ the Rindler observer is causally separated from a part of the Minkowski space-time and does therefore indeed regard the Minkowski vacuum as a mixed state.
To show that the Minkowski vacuum displays the temperature $`T=\kappa /(2\pi )`$ we consider the corresponding two-point Wightman function. The Wightman axioms (in particular the spectral condition) imply that this bi-distribution can be considered as the boundary value of an analytic function. Hence we may restrict to the space-like region for reasons of simplicity where the two-point function assumes the form
$`W(\underset{¯}{x},\underset{¯}{x}^{})=\widehat{\mathrm{\Phi }}(\underset{¯}{x})\widehat{\mathrm{\Phi }}(\underset{¯}{x}^{})={\displaystyle \frac{(2\pi )^2}{(\underset{¯}{x}\underset{¯}{x}^{})^2}}`$ (122)
for 3+1 dimensions. In 1+1 dimensions it behaves as $`\mathrm{ln}[(\underset{¯}{x}\underset{¯}{x}^{})^2]`$ which does not alter the following considerations. Since for a free field all $`n`$-point functions can be derived from this 2-point function it contains all information about the theory.
Now we may consider the two-point function in terms of Rindler coordinates. The $`t,x`$-contribution to the geodesic distance transforms according to Eq. (4.2) into
$`(tt^{})^2(xx^{})^2=2\rho \rho ^{}\mathrm{cosh}\left(\kappa [\tau \tau ^{}]\right)\rho ^2\rho _{}^{}{}_{}{}^{2}.`$ (123)
As a result the two-point function is periodic along the imaginary Rindler time axis and thus satisfies the KMS condition for the temperature $`T=\kappa /(2\pi )`$. This result confirms the considerations in the previous section and justifies the identification of the UV-divergence occurring there with the infinite volume divergence of the Rindler space.
It can be shown quite generally that only the KMS state corresponding to the temperature $`T=\kappa /(2\pi )`$ satisfies the Hadamard condition (local stability) in the complete Rindler space-time (and in particular at the horizon), see . The Hadamard requirement demands the singularity of the two-point function ($`1/s^2`$ and $`\mathrm{ln}s^2`$) to be independent of the state, i.e. it is only determined by the structure of the space-time, see also . This property ensures the validity of the point-splitting renormalisation technique, cf. . As it will become more evident in Section 5.1, an analogue idea can be employed to derive the Hawking effect.
## 5 Black holes
In this Section we are going to apply the formalism presented in Sec. 3 to one of the most fascinating curved space-time structures, the black hole. Various coordinate systems which represent this object are known. For our purpose we have to demand a static metric with a time coordinate $`t`$ corresponding to a Killing vector. Because the black hole space-time becomes asymptotically flat, another requirement is the coincidence of this time coordinate $`t`$ with the usual Minkowski time of an observer at spatial infinity. All this requisites are fulfilled by the Schwarzschild coordinates $`t,r,\vartheta ,\phi `$ for which the black hole metric reads
$`ds^2=h(r)dt^2{\displaystyle \frac{dr^2}{h(r)}}r^2d\vartheta ^2r^2\mathrm{sin}^2\vartheta d\phi ^2.`$ (124)
Other coordinates, e.g. Kruskal, Eddington-Finkelstein, etc. are not suitable for the above reasons. As the Schwarzschild coordinates measure time and length scales with respect to an observer at fixed spatial distance to the black hole all results obtained later refer to this observer.
In order to describe a (non-extreme) black hole with a horizon at $`r=R`$ and a surface gravity $`\kappa >0`$ the function $`h`$ obeys the properties (see e.g. )
$`h(R)=0,\kappa ={\displaystyle \frac{1}{2}}h^{}(R)`$ (125)
and also $`h(r>R)>0`$ together with $`h(r\mathrm{})=1`$. With the aid of this function $`h`$ it is possible to consider the rather general case of a static black hole, for example the Schwarzschild metric with $`h=1R/r`$.
Using these coordinates the canonical conjugate momenta turn out to be $`\widehat{\mathrm{\Pi }}=_t\widehat{\mathrm{\Phi }}/h`$. In terms of these momenta the formal expression for the Hamiltonian density can be cast into the following form
$`\widehat{}={\displaystyle \frac{h}{2}}\widehat{\mathrm{\Pi }}^2+{\displaystyle \frac{h}{2}}\left(_r\widehat{\mathrm{\Phi }}\right)^2+{\displaystyle \frac{1}{2r^2}}\left(_\vartheta \widehat{\mathrm{\Phi }}\right)^2+{\displaystyle \frac{1}{2r^2\mathrm{sin}^2\vartheta }}\left(_\phi \widehat{\mathrm{\Phi }}\right)^2.`$ (126)
The fields $`\widehat{\mathrm{\Phi }}(𝒓,t)`$ as well as their momenta $`\widehat{\mathrm{\Pi }}(𝒓,t)`$ are operator-valued distributions (see also Section 3). Consequently, the Hamiltonian density above is not well-defined. In analogy to Sec. 3.5 it may only be considered as a formal expression until an appropriate regularisation method, for instance the point-splitting technique (see e.g. ), has been applied.
It is possible to split up the Hamiltonian $`\widehat{H}`$ of the field $`\widehat{\mathrm{\Phi }}`$ into two parts $`\widehat{H}=\widehat{H}_>+\widehat{H}_<`$ that account for the interior $`\widehat{H}_<`$ and the exterior $`\widehat{H}_>`$ region of the black hole, respectively
$`\widehat{H}_>`$ $`=`$ $`{\displaystyle d^3𝒓\widehat{}\mathrm{\Theta }(rR)},`$
$`\widehat{H}_<`$ $`=`$ $`{\displaystyle d^3𝒓\widehat{}\mathrm{\Theta }(Rr)}`$ (127)
with the Heaviside step function $`\mathrm{\Theta }`$ and the volume element $`d^3𝒓=\sqrt{g}d^3x=r^2\mathrm{sin}\vartheta drd\vartheta d\phi `$. Employing the equal time commutation relations
$`[\widehat{\mathrm{\Phi }}(𝒓,t),\widehat{\mathrm{\Phi }}(𝒓^{},t)]=[\widehat{\mathrm{\Pi }}(𝒓,t),\widehat{\mathrm{\Pi }}(𝒓^{},t)]=0;[\widehat{\mathrm{\Phi }}(𝒓,t),\widehat{\mathrm{\Pi }}(𝒓^{},t)]=i\delta ^3(𝒓𝒓^{}),`$ (128)
where $`t`$ denotes the Schwarzschild time and represents a Killing vector, one observes that the two parts of the Hamiltonian commute
$`[\widehat{H}_>,\widehat{H}_<]=0.`$ (129)
In the language of point-splitting, cf. and the remarks in Section 4.4, the divergent terms of the Hamiltonian density are independent of the state and therefore pure $``$-numbers which do not contribute to the commutator. The remaining (convergent) operator-valued components commute because of $`h(r=R)=0`$. The same result can be obtained by means of normal ordering or the regularisation described in Eq. (58). Due to $`h(r=R)=0`$ the $`𝒦`$-operator and the operators projecting onto the interior, respectively, exterior domain commute. Hence it is possible to select a basis $`b_I`$ for the inside and outside region separately such that the Hamiltonian possesses no mixing terms.
Accordingly, the separation $`\widehat{H}=\widehat{H}_>+\widehat{H}_<`$ represents two independent systems. This fact displays one advantage of the Schwarzschild coordinates because there is a horizon at $`r=R`$. The consistency with the results of Section 2.1 can be demonstrated if one considers the spatial surface term $`dS_ig^{ij}=d^2x\sqrt{g}n_ig^{ij}=d\vartheta d\phi r^2\mathrm{sin}\vartheta g^{rr}`$ which indeed vanishes for $`r=R`$. As a consequence, it is impossible to transport matter (energy or information) across the horizon, nothing can come out or fall into the black hole. Of course, this holds only for a fixed metric, i.e. if one neglects the back-reaction. Without this restriction it is possible that the horizon increases due to the in-falling matter, and swallows it. It should be emphasised again that all of our assertions refer to an observer at a fixed spatial distance to the black hole and therefore not necessarily to a free falling one.
### 5.1 Black hole exterior
In the following we restrict our considerations to the domain outside the black hole governed by $`\widehat{H}_>`$. The properties of the interior will be discussed in the next Section. The exterior region $`G=\{r>R\}`$ fulfils the conditions imposed in Sec. 2.1 which allows for a particle definition. As a result, the $`\widehat{H}_>`$-part of the Hamiltonian can be diagonalised formally via
$`\widehat{H}_>=\begin{array}{c}\\ {\displaystyle }\\ \mathrm{\Gamma }\end{array}\omega _\mathrm{\Gamma }\widehat{N}_\mathrm{\Gamma }^{\mathrm{BH}}+E_{\mathrm{}},`$ (133)
where $`E_{\mathrm{}}`$ denotes the divergent zero-point energy.
In order to isolate the features that are specific for black holes, the most interesting region is the neighbourhood of the horizon $`rR`$. To investigate the behaviour in this zone we introduce a dimensionless variable $`\chi `$ with
$`\chi =2\kappa (rR)h=\chi \left(1+𝒪(\chi )\right).`$ (134)
Without loosing the static character of the metric it is possible to perform a radial coordinate transformation for $`r>R`$ via
$`r_{}={\displaystyle \frac{dr}{h}}={\displaystyle \frac{\mathrm{ln}\chi }{2\kappa }}+𝒪(\chi ).`$ (135)
The new radial $`r_{}`$ coordinate is called the the Regge-Wheeler tortoise coordinate. According to the above arguments it is sufficient to cover the region outside the horizon by the new coordinate. The function $`h`$ and the original radial variable $`r`$ have to be considered as functions of the introduced coordinate: $`r=r(r_{})`$ and $`h=h(r_{})=h(r[r_{}])`$. The tortoise coordinate has the advantage of a very simple form of the $`𝒦`$-operator
$`𝒦={\displaystyle \frac{1}{r^2}}{\displaystyle \frac{}{r_{}}}r^2{\displaystyle \frac{}{r_{}}}h_{\vartheta \phi }^2={\displaystyle \frac{^2}{r_{}^2}}+𝒪(\chi ),`$ (136)
together with $`G=\{r_{}\}`$. The generalised eigenfunctions $`F_\mathrm{\Gamma }^{\mathrm{BH}}(\underset{¯}{x})`$ of this operator behave $`𝒪(\chi )`$-approximately as $`\mathrm{exp}(\pm i\omega r_{})`$ and after the separation of the angular variables they can be written as follows
$`F_\mathrm{\Gamma }^{\mathrm{BH}}(\underset{¯}{x})=F_{\xi \omega \mathrm{}m}^{\mathrm{BH}}(t,\chi ,\vartheta ,\phi )=𝒩_\omega \mathrm{}^{\mathrm{BH}}{\displaystyle \frac{e^{i\omega t}}{\sqrt{\omega }}}\chi ^{i\xi \omega /(2\kappa )}𝒴_\mathrm{}m(\vartheta ,\phi )\left(1+𝒪(\chi )\right).`$ (137)
$`𝒩_\omega \mathrm{}^{\mathrm{BH}}`$ symbolises a normalisation factor which may without any loss of generality chosen to be independent of $`\xi `$. These eigenfunctions are rapidly oscillating near the horizon.
By inspection, we recognise the occurrence of two generalised eigenfunctions for a given frequency $`\omega `$ and fixed angular dependence $`\mathrm{},m`$ distinguished by $`\xi =\pm 1`$. Thus the definition of ingoing and outgoing particles is possible in this case. (It should be mentioned that potential scattering effects cause slight deviations from the purely ingoing and outgoing behaviour in Eq. (137) at $`r\mathrm{}`$. However, this way of definition does not alter the conclusions.) This – perhaps surprising – fact can be elucidated in the following way. The horizon separates the space into two independent domains (interior and exterior) and prevents the field modes outside from being influenced by the effective ”boundary condition” at $`r=0`$. In view of the study of the $`𝒦`$-operator in terms of the tortoise coordinate $`r_{}`$ one may consider the horizon as some new kind of spatial infinity ($`r_{}\mathrm{}`$) in addition to $`r\mathrm{}`$.
Also for the black hole example the horizon acts opposite to a mirror, cf. Section 4.2. Even for the scenario of a black hole which is enclosed in a large box with Dirichlet boundary conditions the spectrum of the operator $`𝒦`$ is still continuous – but now not degenerated.
For a black hole in an asymptotically flat (unbounded) space-time there are two contributions to the infinite volume divergence (see Sec. 3.8) $`\delta (\mathrm{\Gamma },\mathrm{\Gamma })=𝒩_VV`$: firstly, the usual infinity $`r,r_{}\mathrm{}`$ and secondly, the effective infinity at the horizon $`rR`$ respectively $`r_{}\mathrm{}`$. The former divergence $`\delta _+(\mathrm{\Gamma },\mathrm{\Gamma })`$ does also arise in the (unbounded) Minkowski space-time – but not inside a finite box (e.g. with Dirichlet boundary conditions) – whereas the latter divergence $`\delta _{}(\mathrm{\Gamma },\mathrm{\Gamma })`$ is restricted to the scenario of a black hole, but it is not affected by a finite box. E.g., the expectation value of the number ”operator” $`\widehat{N}_\mathrm{\Gamma }`$ in any KMS state (with a non-vanishing temperature) contains the complete divergence $`\widehat{N}_\mathrm{\Gamma }_T\delta _+(\mathrm{\Gamma },\mathrm{\Gamma })+\delta _{}(\mathrm{\Gamma },\mathrm{\Gamma })`$. One important example is the Israel-Hartle-Hawking state, the KMS state corresponding to the Hawking temperature $`T=\kappa /(2\pi )`$. For large radial distances to the black hole the (renormalised) expectation value of the energy-momentum tensor evaluated in this state approaches a constant value (proportional to $`T^4`$). In contrast, for the Unruh state – the state describing the black hole evaporation – the (renormalised) energy density decreases with $`1/r^2`$ for large $`r`$. As a consequence, the expectation value of the number of particles in this state does not display the complete divergence $`\delta _+(\mathrm{\Gamma },\mathrm{\Gamma })`$.
It might be interesting to illustrate the point above with the aid of the Bogoliubov coefficients: If we consider the spherically symmetric collapse of a star to a black hole the metric outside the initial radius of the star does not change (Birkhoff theorem). Ergo the behaviour of the modes at very large radial distances $`r`$ is not affected by the collapse. Accordingly, this region does not contribute to the $`\beta _{\omega ,\omega ^{}}`$-coefficients and generates a $`\delta (\omega \omega ^{})`$-term for the $`\alpha _{\omega ,\omega ^{}}`$-coefficients, cf. also . Recalling the unitarity relation for the Bogoliubov coefficients in Eq. (114) we arrive at the conclusion that exactly this term generates the $`\delta _+(\mathrm{\Gamma },\mathrm{\Gamma })`$-part of the infinite volume divergence. Following Hawking we assume that – for large initial frequencies $`\omega `$ – the Bogoliubov coefficients are related via Eq. (106) in analogy to Section 4.3. Proceeding in the same way as in that Section we observe that the $`\omega `$-integration of the absolute value squared of the $`\beta _{\omega ,\omega ^{}}`$-coefficients is UV-divergent again. But in contrast to the Unruh effect this divergence does not contain $`\delta _+(\mathrm{\Gamma },\mathrm{\Gamma })`$, but only $`\delta _{}(\mathrm{\Gamma },\mathrm{\Gamma })`$ owing to the unitarity relation (114). Consequently the Minkowski vacuum is a KMS state with respect to the Rindler observer, but it does not transform into a KMS state during the collapse to a black hole (if we assume the space-time to be asymptotically flat and therefore unbounded). Hawking derived the relation (106) only for the finally outgoing particles. But even if this relation would hold for both, the (finally) ingoing and outgoing particles, the state would still contain less particles than the corresponding KMS state.
As it became evident in the previous considerations, the vicinity of the horizon of a black hole displays many similarities to the scenario of the Unruh effect in Section 4.3. Indeed, with $`\chi =\kappa ^2\rho ^2`$ the black hole metric approaches the Rindler metric in that region
$`ds^2=\left(\kappa ^2\rho ^2dt^2d\rho ^2R^2d\mathrm{\Omega }^2\right)\left(1+𝒪(\chi )\right),`$ (138)
together with the angular part $`d\mathrm{\Omega }^2`$. This observation motivates an argumentation analogue to that at the end of Sec. 4.4, cf. also . Indeed, it is possible to prove that for a state fulfilling the Hadamard requirement (among other conditions, see ) throughout the complete space-time (and in particular at the horizon) the asymptotic expectation values correspond to the Hawking temperature. The ground state of the quantum field (the Boulware state) as well as every KMS state (with an arbitrary temperature) obey the Hadamard singularity structure away from the horizon, see . But only that KMS state that corresponds to the Hawking temperature $`T=\kappa /(2\pi )`$ – i.e. the Israel-Hartle-Hawking state – matches the Hadamard condition at the horizon. (The same holds true for the Unruh state.) It can be shown that the Hadamard condition is conserved during the dynamics of a $`C^{\mathrm{}}`$ space-time. Accordingly, if the collapse of a star to a black hole can be described by a $`C^{\mathrm{}}`$-metric, the consideration above can be used to deduce the Hawking effect. (The Minkowski vacuum is of course also a Hadamard state.) However, dropping the assumption of a $`C^{\mathrm{}}`$ space-time the situation becomes less clear.
If we compare the outcome of this Section with the Minkowski example, we arrive at the conclusion that the formation of the horizon causes a bifurcation in a double sense:
The total Hamiltonian of the field $`\widehat{H}`$ splits up into two commuting parts $`\widehat{H}_<`$ and $`\widehat{H}_>`$ which account for two independent (physical complete) regions $`r<R`$ and $`r>R`$, respectively.
Before the horizon has been formed there exists only one generalised eigenfunction for every given frequency and fixed angular behaviour. This property forbids the definition of ingoing and outgoing particles (see Section 4.1). After the horizon has been formed the spectrum is twice degenerated and the definition of ingoing and outgoing particles becomes possible.
### 5.2 Black hole interior
Our previous investigations focused on the exterior of the black hole. As indicated before we shall now take the interior region into account. Inside the (non-extreme) black hole it yields $`h(r)<0`$ and therefore $`g_{tt}<0`$, $`g_{rr}>0`$, $`g_{\vartheta \vartheta }<0`$ and $`g_{\phi \phi }<0`$. As a consequence the signature of the metric is degenerated and thus the particle and vacuum definition proposed in Section 3 does not apply. However, it is still possible to obtain a self-adjoint $`𝒦`$-operator governing the dynamics of the system. But for this purpose some modifications are necessary with the result that $`𝒦`$ is not given by $`𝒟^{}\overline{𝒟}`$ and therefore not non-negative. As we shall see later the negative parts of $`𝒦`$ correspond to unstable fields modes.
At first the scalar product of the interior region $`\{|\}_1^<`$ has to be defined with $`|g^{00}|`$ instead of $`g^{00}`$ in Eq. (14) in order to obtain a positive-definite bilinear form. For reasons of simplicity we restrict our further considerations to the Schwarzschild metric $`h(r)=1R/r`$ and start with the functions
$`C_0^{\mathrm{}}(G_<)=\mathrm{lin}\left\{C_0^{\mathrm{}}(0<r<R)C^{\mathrm{}}(𝕊_2)\right\},`$ (139)
where $`𝕊_2`$ denotes the 2-sphere of $`\vartheta `$ and $`\phi `$. Again the Hilbert space $`L_2(G_<,|g^{00}|)`$ is given by the completion of all these functions with respect to the (redefined) scalar product.
The degenerated signature permits the definition of a scalar product containing $`g^{ij}`$. Accordingly, the subsequent steps in Sec. 3 cannot be adopted here. In particular we cannot introduce an operator $`𝒟`$ such that the self-adjoint $`𝒦`$-operator is represented by the absolute value squared of $`𝒟`$. Instead we may define an operator $`𝒦_0`$ via
$`𝒦_0:C_0^{\mathrm{}}(G_<)L_2(G_<,|g^{00}|)`$ $``$ $`L_2(G_<,|g^{00}|)`$
$`\varphi `$ $``$ $`{\displaystyle \frac{h}{r^2}}{\displaystyle \frac{}{r}}hr^2{\displaystyle \frac{\varphi }{r}}h_{\vartheta \phi }^2\varphi .`$ (140)
The second term at the r.h.s. of the above expression for $`𝒦_0`$ generates the negative parts of this operator. These negative parts originate from the angular derivatives and cannot be obtained in a purely radial symmetric consideration.
Obviously $`𝒦_0`$ is Hermitian with respect to the scalar product containing the weight $`|1/h|`$ (and $`\sqrt{g}=r^2\mathrm{sin}\vartheta `$). In addition – since $`\mathrm{Def}(𝒦_0)=C_0^{\mathrm{}}(G_<)`$ is dense in the underlying Hilbert space $`L_2(G_<,|g^{00}|)`$ – it is densely defined and therefore symmetric.
Now we can make use of the following theorem (see e.g. ): Every symmetric and $``$-real operator acting on a complex Hilbert space possesses (at least one) self-adjoint extension(s). As a result we will always find a self-adjoint operator $`𝒦`$ (as an appropriate extension of $`𝒦_0`$) governing the dynamics of the field. In terms of $`𝒦`$ the Lagrange function in Eq. (13) for the interior domain assumes the simple form
$`L_<={\displaystyle \frac{1}{2}}\left\{\dot{\mathrm{\Phi }}\right|\dot{\mathrm{\Phi }}\}_1^<+{\displaystyle \frac{1}{2}}\left\{\mathrm{\Phi }\right|𝒦|\mathrm{\Phi }\}_1^<.`$ (141)
Note that in contrast to Eq. (20) the global sign has changed. However, this global sign does not affect the equation of motion, but – as it will become evident later – the negative parts of the $`𝒦`$-operator do so.
Since the self-adjoint $`𝒦`$-operator represents an extension of the original operator $`𝒦_0`$ these two operators have to coincide on the subspace $`C_0^{\mathrm{}}(G_<)`$. Accordingly, it is possible to construct test functions $`w(r,\vartheta ,\phi )=w(r)𝒴_\mathrm{}m(\vartheta ,\phi )C_0^{\mathrm{}}(G_<)`$ generating negative expectation values of the $`𝒦`$-operator via
$`\left\{w\right|𝒦|w\}={\displaystyle \underset{0}{\overset{R}{}}}𝑑r\left(|h|r^2\left|_rw\right|^2\mathrm{}(\mathrm{}+1)\left|w\right|^2\right)<0.`$ (142)
If we choose the angular quantum number $`\mathrm{}`$ very large the expectation value $`\left\{w\right|𝒦|w\}`$ equals negative numbers of arbitrarily large absolute values, even for normalised test functions $`\left\{w\right|w\}=1`$. Hence the spectrum of $`𝒦`$ is unbounded from below. (Of course it is also unbounded from above.) Diagonalising the Hamiltonian by means of a quasi-unitary transformation $`𝒱`$ in analogy to Sec. 3 yields
$`\widehat{H}_<={\displaystyle \frac{1}{2}}\begin{array}{c}\\ {\displaystyle }\\ \mathrm{\Gamma }\end{array}\left(\widehat{P}_\mathrm{\Gamma }^2+\lambda _\mathrm{\Gamma }\widehat{Q}_\mathrm{\Gamma }^2\right).`$ (146)
The interior Hamiltonian is still self-adjoint (by Stone’s theorem) – but it is not bounded from above and below. Ergo it does not possess a ground state, and a definition of particles as excitations over the ground state is impossible.
As mentioned before, the global sign does not affect the equations of motion, but the occurring negative eigenvalues $`\lambda _\mathrm{\Gamma }`$ do so: The modes $`\mathrm{\Gamma }`$ corresponding to negative eigenvalues $`\lambda _\mathrm{\Gamma }`$ obey the following equations of motion
$`{\displaystyle \frac{d^2}{dt^2}}\widehat{Q}_\mathrm{\Gamma }=|\lambda _\mathrm{\Gamma }|\widehat{Q}_\mathrm{\Gamma }.`$ (147)
Their solutions $`\mathrm{exp}(\pm \sqrt{|\lambda _\mathrm{\Gamma }|}t)`$ display a highly (linear) unstable behaviour.
This instability cannot be avoided by introducing an indefinite metric of the Fock space if we assume the black hole to be formed by a collapse because in this case the Fock space is initially well-defined and obeys a positive definite metric $`\mathrm{\Psi }|\mathrm{\Psi }0`$.
One might suspect that the initial conditions are just in such a way that the exponentially increasing solutions do not occur, cf. also . Employing an analogue from classical mechanics this situation corresponds to a point-mass moving on the top of a parabolic hill which just comes to rest at the zenith of the parabola. However, within quantum theory no (regular) stationary state exists in such a scenario (Heisenberg uncertainty relation). Even if the expectation value of the amplitude $`\widehat{Q}_\mathrm{\Gamma }`$ vanishes for all times, its variance increases exponentially (for late times).
Since the unstable behaviour described above accounts for the time-evolution of the (global) modes $`\mathrm{\Gamma }`$ it describes a global instability which should not be confused with the concept of local stability usually associated with the Hadamard condition, cf. and .
It should be mentioned here that potential terms (which we have omitted in Sec. 2) may also give raise to negative parts of the $`𝒦`$-operator. E.g., if the assumptions in Section 3.2 fail and the scalar curvature $``$ assumes negative values over a large enough volume the operators $``$ and $`𝒦`$ are not non-negative. However, in this situation the $`𝒦`$-operator is still bounded from below (if $``$ does not diverge). Hence only modes up to a certain quantum number are unstable. These modes are strongly correlated to the global structure of the space-time. Special care is required concerning the interpretation of the instability caused by mass terms. Mass terms that are generated by the Higgs mechanism occur in the effective Lagrangian for low excitations and cannot be extrapolated to large amplitudes. Restricting ourselves to the massless and minimally coupled scalar field (as a model for the photon field) only the instability due to the angular derivatives remains where all these objections do not apply.
In order to interpret the instability it might be interesting to investigate the corresponding proper or generalised eigenfunctions. Near the horizon (inwards), the modes behave as
$`f_\mathrm{\Gamma }\mathrm{exp}\left(r_{}\sqrt{\lambda _\mathrm{\Gamma }}\right)\left(2\kappa [Rr]\right)^{\sqrt{\lambda _\mathrm{\Gamma }}/(2\kappa )}.`$ (148)
Depending on the behaviour at the origin $`r=0`$ one might expect the existence of proper eigenfunctions $`f_\mathrm{\Gamma }`$ at some points of the negative part of the spectrum.
However, even if no proper and (pointwise/locally defined) generalised eigenfunctions exist, one may still construct suitable distributions $`f_\mathrm{\Gamma }`$ with analogous properties : Considering the Schwartz/Sobolev space $`𝔖_1(\sigma ,𝒱)L_2(\sigma ,𝒱)`$ of all continuous functions over the spectrum $`\sigma `$ of the $`𝒦`$-operator we may define a Dirac $`\delta `$-distribution as a linear functional over this space. This distribution $`\delta _\mathrm{\Gamma }=\delta (\lambda ,\lambda _\mathrm{\Gamma })`$ is then defined within the dual space $`𝔖_1(\sigma ,𝒱)`$. It represents a generalised eigendistribution of the diagonalised $`𝒦`$-operator $`𝒱𝒦𝒱^{}\delta _\mathrm{\Gamma }=\delta _\mathrm{\Gamma }=\lambda _\mathrm{\Gamma }\delta _\mathrm{\Gamma }`$. Hence its spatial representation $`f_\mathrm{\Gamma }=𝒱^{}\delta _\mathrm{\Gamma }`$ exists at least as a distribution over $`𝒱^{}𝔖_1(\sigma ,𝒱)`$ (which is dense in $`L_2(G_<,|g^{00}|)`$) and describes an eigendistribution of $`𝒦`$. The construction described above generates non-vanishing eigendistributions $`f_\mathrm{\Gamma }`$ for all non-singular points (of the spectral measure) $`\mathrm{\Gamma }`$ of the spectrum $`\sigma `$. Since every open interval of $`\sigma `$ contains non-singular points we can always find an appropriate mode $`\mathrm{\Gamma }`$ where $`f_\mathrm{\Gamma }`$ exists.
Using these eigendistributions $`f_\mathrm{\Gamma }(𝒓)`$ we can construct solutions of the Klein-Fock-Gordon equation of the form
$`F_\mathrm{\Gamma }(𝒓,t)=\mathrm{exp}\left(\pm \sqrt{|\lambda _\mathrm{\Gamma }|}t\right)f_\mathrm{\Gamma }(𝒓),`$ (149)
if we choose a mode $`\mathrm{\Gamma }`$ from the negative part of the spectrum. As a consequence, the equation of motion does not only possess unstable solutions – even the degree of the instability $`\sqrt{|\lambda _\mathrm{\Gamma }|}`$ can be arbitrarily large. In a vivid description one may speak about an explosion interiorly.
It should be mentioned here that a (partial) negative Hamiltonian, i.e. a (partial) negative generator of the time-evolution, is not sufficient for the prediction of an instability. As a counter-example we may consider a 1+1 dimensional black hole with $`ds^2=hdt^2dr^2/h`$. In this situation there are no angular terms and thus the interior as well as the exterior $`𝒦`$-operator are both non-negative. Consequently the equation of motion is completely stable. Of course, the interior Hamiltonian $`\widehat{H}_<`$ displays a global minus sign, but this does not affect the equation of motion
$`\widehat{H}=\widehat{H}_>+\widehat{H}_<={\displaystyle \frac{1}{2}}\begin{array}{c}\\ {\displaystyle }\\ \mathrm{\Gamma },>\end{array}\left(\widehat{P}_{\mathrm{\Gamma },>}^2+\mathrm{\Omega }_{\mathrm{\Gamma },>}^2\widehat{Q}_{\mathrm{\Gamma },>}^2\right){\displaystyle \frac{1}{2}}\begin{array}{c}\\ {\displaystyle }\\ \mathrm{\Gamma },<\end{array}\left(\widehat{P}_{\mathrm{\Gamma },<}^2+\mathrm{\Omega }_{\mathrm{\Gamma },<}^2\widehat{Q}_{\mathrm{\Gamma },<}^2\right).`$ (156)
Moreover, although the total Hamiltonian is unbounded from above and below, it splits up into two independent parts which are bounded. The existence of a horizon is essential for this bifurcation. In a flat space-time the Wightman axioms (spectral condition) demand a non-negative generator for stability.
The Schwarzschild metric $`ds^2=hdt^2dr^2/h`$ or $`ds^2=hdt^2dr^2/hr^2d\mathrm{\Omega }^2`$ possesses a unique analytic continuation to values of $`r`$ beyond the horizon $`r<R`$. In contrast the analytic continuation of the Rindler metric to negative values of $`\rho `$ does not lead to a degenerated signature and complex values of $`\rho `$ and/or $`\tau `$ do not describe a physical sheet of the space-time. As a consequence one observes no instability in the Rindler metric – i.e. the scenario of the Unruh effect.
The notion of the unstable behaviour obtained above refers to the time $`t`$ measured by an (outside) observer at a fixed spatial distance to the black hole. One might argue that this time coordinate is not capable for describing effects inside the black hole due to the coordinate singularity at $`r=R`$. However, the instability obtained above is not restricted to the Schwarzschild coordinates – it occurs in other coordinate systems as well: By virtue of the transformation
$`dtdt\pm {\displaystyle \frac{\sqrt{R/r}}{1R/r}}dr.`$ (157)
the metric of the black hole can be cast into the Painlevé-Gullstrand-Lemaître form
$`ds^2=\left(1{\displaystyle \frac{R}{r}}\right)dt^2\pm 2\sqrt{{\displaystyle \frac{R}{r}}}dtdrdr^2r^2d\mathrm{\Omega }^2.`$ (158)
This metric is regular everywhere except at the singularity at $`r=0`$. The transformation of the unstable solutions in Eq. (149) into this coordinate system via Eq. (157), i.e. $`tt\pm \mathrm{\Xi }(r)`$, merely results in a simple $`r`$-dependent factor
$`F_\mathrm{\Gamma }(𝒓,t)=\mathrm{exp}\left(\sqrt{|\lambda _\mathrm{\Gamma }|}t\right)f_\mathrm{\Gamma }(𝒓)\mathrm{exp}\left(\pm \sqrt{|\lambda _\mathrm{\Gamma }|}\mathrm{\Xi }(r)\right),`$ (159)
while the unstable behaviour persists. The same holds true for the Eddington-Finkelstein coordinates with $`v=t+r_{}`$
$`ds^2=\left(1{\displaystyle \frac{R}{r}}\right)dv^22dvdrr^2d\mathrm{\Omega }^2.`$ (160)
Within these coordinates ingoing light rays are simply governed by $`v=\mathrm{const}`$. Both coordinate systems lead to a stationary (but not static) metric, i.e. the evolution parameter still coincides with a Killing vector. (This is not the case for the Kruskal metric.) In summary the instability of the field equation inside the black hole turns out to be a quite general phenomenon.
### 5.3 Back-reaction
The Eddington-Finkelstein metric in Eq. (160) allows for a demonstrative visualisation of the unstable behaviour: If one emits radially ingoing light pulses in uniform intervals these beams are labelled by equidistant values of $`v`$. According to the results of the previous Section the amplitude of the field $`\mathrm{\Phi }`$ inside the black hole increases exponentially with rising numbers $`v`$ of the light rays. Hence we may draw the conclusion that the instability is not just an artifact caused by an inappropriate description but a physical effect.
Nevertheless, for an eternal black hole the outside observer is completely causally separated from the region of the instability. Hence the interpretation of the unstable behaviour is not obvious in that case. But if one considers the possibility of the decay of the black hole (no matter whether via evaporation or explosion) and assumes that this decay can be described using one of the coordinates above the unstable behaviour should be relevant. (Of course, the assumption of an eternal black hole automatically excludes some of the scenarios where the instability may become relevant.)
In order to investigate the consequences of the instability one has to deal with the back-reaction problem. Within all of our previous considerations the quantum field was regarded as a test field, i.e. it did not influence the given (externally prescribed) space-time. It is known from classical field theory (see e.g. and references therein) that the formation of the horizon and the singularity may well be affected by the scalar field $`\mathrm{\Phi }`$. (For quantum field theory one expects that the back-reaction will become important at the Planck scale.) However, Ref. deals with radially symmetric fields only. For that reason the unstable behaviour was not obtained there. The correct implementation of the back-reaction of a quantum field has to be determined by an underlying theory unifying gravitational and quantum effects. Since we have no well-established solution to this problem, we may only speculate about the impact of the quantum field on the metric based on physical reasonable arguments. There are several possible consequences:
* The explosion of the complete black hole
The unstable field modes evolve as $`\mathrm{exp}(\sqrt{|\lambda _\mathrm{\Gamma }|}[tr_{}])`$. Hence they ”reach” after a finite period of time the Planck scale vicinity of the horizon, where the classical treatment of the gravitation is expected to break down. In that case one might imagine that the ”wave front” destroys the horizon and thus the complete black hole. (Such an event might perhaps be regarded as a toy candidate for the big bang.) In view of arguments concerning the time-reversal symmetry there is no obvious reason why the explosion of the complete black hole should be impossible.
As long as there is some matter falling into the black hole its horizon increases. Depending on the particular dynamics of the metric this may prevent the ”wave front” from ”reaching” the vicinity of the horizon. But for a static black hole there is no way to avert the impact.
One should be aware that most of the theorems of classical general relativity – e.g. the black hole analogues of the laws of thermodynamics – are based on appropriate energy conditions, cf. . But incorporating the expectation value of the energy-momentum tensor of the quantum field these energy conditions do not hold in general. In some cases one may employ averaged energy conditions instead, but even the validity of an averaged condition is by no means obvious in view of the unstable solutions of the field equation.
* The prevention of the singularity at $`r=0`$
One might expect that the impact of the instability is at the origin $`r=0`$ much stronger than at the horizon $`r=R`$. In fact, also those theorems of general relativity that predict a space-time singularity after a gravitational collapse are based on energy conditions. Accordingly, taking the back-reaction of the quantum field into account, the formation of the singularity may perhaps be avoided. Instead one might imagine some kind of quasi-oscillations: Impelled by the (exponentially large) amplitudes of the quantum field, the matter around the origin blows up, absorbs the excitations of the field, collapses (while the field repeatedly evolves exponentially), and eventually blows up again.
* The field does not affect the metric
This possibility cannot be excluded within the framework of quantum field theory in given (external) space-times. However, the situation of a completely static black hole (neglecting the Hawking effect, which is very small for macroscopic black holes) seems to be rather strange. In that case the amplitude of the field exceeds the Planck scale after a finite period of time (measured by an outside observer). Hence one would expect drastic modifications of the space-time.
### 5.4 Sonic analogue of black holes
In 1980 Unruh discovered a very interesting model for the kinematics of fields in curved space-times. He considered the propagation of sound waves in flowing fluids where the effective equation of motion assumes the same form as the Klein-Fock-Gordon equation in curved space-times. The effective metric depends on the particular flow profile. Many investigations have been devoted to this topic during the last years, see e.g. , the recent work , and references therein.
Before discussing the consequences of the results of the previous section within this scenario we shall repeat the basic ideas: The flow of a fluid can be described by its local velocity field $`𝒗`$, its density $`\varrho `$, and the pressure $`p`$. The dynamics of the fluid is governed by the non-linear Euler equation
$`\dot{𝒗}+(𝒗)𝒗+{\displaystyle \frac{p}{\varrho }}=𝒇_{\mathrm{ext}},`$ (161)
if we neglect the viscosity, and the equation of continuity
$`\dot{\varrho }+(\varrho 𝒗)=0.`$ (162)
For reasons of simplicity we restrict our further considerations to a constant speed of sound $`c_\mathrm{s}`$. This implies the very simple relation between the density and the pressure $`p=c_\mathrm{s}^2\varrho `$. If we assume an irrotational flow $`\times 𝒗=0`$, we may introduce a generating scalar field $`𝒗=\mathrm{\Phi }`$. Now we linearise the non-linear system of the two equations above around a fixed background solution via
$`\mathrm{\Phi }`$ $`=`$ $`\mathrm{\Phi }_0+\epsilon \mathrm{\Phi }_1+𝒪(\epsilon ^2),`$
$`𝒗`$ $`=`$ $`𝒗_0+\epsilon 𝒗_1+𝒪(\epsilon ^2),`$
$`p`$ $`=`$ $`p_0+\epsilon p_1+𝒪(\epsilon ^2),`$
$`\varrho `$ $`=`$ $`\varrho _0+\epsilon \varrho _1+𝒪(\epsilon ^2).`$ (163)
This enables us to consider the propagation of small perturbations – i.e. sound waves – within a given flow profile. It turns out that the potential $`\mathrm{\Phi }_1`$ of the fluctuations satisfies the Klein-Fock-Gordon equation with the effective (acoustic) metric
$`g_{\mu \nu }={\displaystyle \frac{\varrho _0}{c_\mathrm{s}}}\left(\begin{array}{cc}c_\mathrm{s}^2𝒗_0^2& \hfill 𝒗_0\\ 𝒗_0& \hfill \mathrm{𝟏}\end{array}\right).`$ (166)
Ergo sound waves in flowing fluids share a lot of interesting features with fields in curved space-times. E.g., the surface of transition from subsonic to supersonic flow represents the acoustic analogue of a horizon. For a stationary and radially symmetric flow this surface possesses even the properties of an event and an apparent horizon. (Unfortunately this scenario exhibits the problem of fluid conservation at $`r=0`$ which has to be evaded in some way.)
Selecting a particular velocity profile $`𝒗=\pm 𝒓\sqrt{R/r^3}`$ it is possible to simulate a space-time which obeys – up to a conformal factor $`r^{3/2}`$ – the Painlevé-Gullstrand-Lemaître metric in Eq. (158). According to the results of the previous section the Klein-Fock-Gordon equation possesses unstable solutions inside the black hole. Consequently, also the sound waves within the supersonic region obey an instability. The conformal factor mentioned above and the coordinate transformation in Eq. (157) do not alter this conclusion – see the remarks in the previous Section.
In contrast to the ”real” black hole, where the consequences of the instability are not a priori clear (back-reaction problem), there is no possibility to avoid the instability for the acoustic black hole models since in that case $`t`$ denotes the appropriate time also for an inside observer and the sound waves affect the fluid directly.
In the theory of fluid dynamics, such an instability is a well-known indicator for the breakdown of the laminar (irrotational) flow, see e.g. . I.e., that flow does not represent a stable fixed point of the non-linear equation of motion. Accordingly, any small disturbance will grow up exponentially until the non-linear regime has been reached. (It should be mentioned here that the unstable behaviour obtained above is not a downstream instability, cf. , since the perturbation increases exponentially also at a fixed radius $`r`$.) In order to investigate the behaviour of the flow after leaving the unstable fixed point – e.g. pattern formation or turbulence – one has to consider the non-linear region. For quantum fields in curved space-times one expects to reach the non-linear regime at the Planck scale where the back-reaction strongly contributes.
Recalling the outcome of the previous Section the unstable behaviour of the equation of motion results from the angular derivatives of the $`𝒦`$-operator. Ergo we may draw the conclusion that the quantum field inside the black hole as well as the supersonically flowing fluid favour a spontaneous breaking of the radial symmetry, similar to the formation of a vortex in the drain of a basin.
## 6 Conclusions
### 6.1 Summary
For a minimally coupled, massless and neutral scalar quantum field $`\widehat{\mathrm{\Phi }}`$ propagating in an arbitrary physical complete and causal space-time $`M`$ that possesses a static metric of non-degenerated signature it is possible to perform a particle definition via diagonalisation of the Hamiltonian.
Application of this method to the 3+1 dimensional Minkowski space-time yields the nonexistence of radial ingoing and outgoing particles. For the 1+1 dimensional Rindler metric we exactly recover the well-known Unruh effect.
If we employ the same formalism in order to investigate a black hole the associated space-time splits up into two independent domains, inside and outside the horizon, respectively. Within the presented approach a particle definition can be accomplished for the exterior region only.
The quantum field inside the black hole possesses a highly unstable behaviour. The corresponding Hamiltonian is unbounded from above and below. Accordingly, it is not possible to define a vacuum as its ground state and particles as excitations over this state.
This instability is not a remnant of an inappropriate description but a physical effect. Due to our lack of understanding the unification of quantum theory and gravity the consequences of this effect are not altogether clear (back-reaction problem). In view of the sonic analogues of black holes – where the unstable solutions go along with the breakdown of the laminar flow – one might expect that the instability indicates (at least) the breakdown of the treatment of quantum fields in given (externally prescribed) space-times.
### 6.2 Discussion
In order to elucidate the outcome of the formalism presented in this article it might be interesting to discuss the main statements together with their relations to other approaches:
As we have observed in Section 4.1, the particle definition via diagonalisation of the Hamiltonian (equivalent to the energy) does not allow for the introduction of ingoing and/or outgoing particles in the Minkowski space-time. The same holds true for more general regular space-times. As a consequence, the vacuum coinciding with the ground state cannot be defined as that state that is annihilated by the ”operators” $`\widehat{A}_\mathrm{\Gamma }`$ corresponding to purely ingoing (and/or outgoing) components $`\mathrm{exp}(i\omega v)/r`$ (and/or $`\mathrm{exp}(i\omega u)/r`$), with $`v=t+r`$ and $`u=tr`$. Instead the ground state gets annihilated by ”operators” (strictly speaking, operator-valued distributions) corresponding to standing waves, i.e. superpositions of ingoing and outgoing components with equal weights. Ergo, considering the collapse of a star to a black hole the initial ground state cannot be uniquely and consistently defined by the requirement ”no ingoing/incoming particles/radiation”. Ref. states explicitly: Note that we have not defined the vacuum by minimizing some positive-definite-operator expectation value (e.g. the Hamiltonian), but we have defined the vacuum as the state with no incoming particles. In order to investigate the relationship of the state defined in this way and the initial ground state additional considerations are necessary.
In contrast to the Minkowski case the ground state of the quantum field in the exterior black hole space-time – the Boulware state – has to be defined via demanding that the action of the annihilators for both, the ingoing and outgoing modes, yields zero: $`_{\xi \omega \mathrm{}m}:\widehat{A}_{\xi \omega \mathrm{}m}|\mathrm{\Psi }_\mathrm{B}=0`$. This fact illustrates the bifurcation caused by the formation of the horizon.
However, if we assume the black hole to be enclosed by a large sphere with e.g. Dirichlet boundary conditions then the definition of ingoing or outgoing particles is impossible again. This observation demonstrates manifestly that the particle interpretation is a global concept – it may be influenced by objects (e.g. the sphere) at arbitrarily large distances.
As another difference between the black hole and the Minkowski situation we may recall the fact that the $`𝒦`$-operator of the black hole possesses – even in the presence of a finite sphere – a continuous spectrum. Due to the additional effective infinity at the horizon the infinite volume divergence of the black hole space-time cannot be regularised by enclosing it by a finite box. (This regularisation applies only to space-time without any horizon.)
There are two main interpretations of the Hawking effect: The first view considers the particles to be produced by the dynamics of the space-time during the collapse while within the second view the radiation is created in a steady rate after the collapse. The observations in Section 5, i.e. the splitting of the total Hamiltonian into two independent parts and the diagonalisation of the exterior part by a suitable particle definition (where the number of particles is conserved), supports the former interpretation.
The Hawking effect may be regarded as the verification of the extension of the laws of thermodynamics to objects like black holes. This effect allows us to assign a temperature to the black hole via $`T=1/(8M)`$ for the Schwarzschild black hole with $`h=12M/r`$. As a result the associated heat capacity of the black hole turns out to be negative: If the mass/energy increases the temperature decreases. The classical laws of thermodynamics predict that an object obeying a negative heat capacity will be unstable. Accordingly, the instability of the black hole interior as observed in Section 5.2 might also be regarded as a verification of the application of thermodynamics to black holes.
The consequences of the unstable behaviour of the Klein-Fock-Gordon equation in the interior of the black hole cannot be deduced rigorously within the framework of quantum fields in (externally prescribed) space-times. The evaluation of the impact of this instability demands the knowledge of the back-reaction which has to be determined by a unifying theory. Nevertheless, if the underlying theory possesses an evolution parameter corresponding to the Schwarzschild time $`t`$ (or one of the other coordinates discussed in Sec. 5.2) and contains the treatment of quantum fields and external metrics in some limiting case, then one would expect that the representation of a black hole also obeys the linearly unstable behaviour. (This would be consistent with the frequently adopted interpretation that black holes are highly excited states of the unifying theory.)
For the situation of the acoustic black hole the interpretation of the unstable behaviour is more obvious. Without any mechanism preserving (enforcing) the radial symmetry (e.g. effects of super-fluids) it is probably impossible to realise the sonic analogue of a black hole experimentally.
### 6.3 Outline
The particle definition presented in this article is restricted to static space-times. This includes the Schwarzschild and the Reissner metric, but not the Kerr space-time describing a rotating black hole. Accordingly, further investigations should be devoted to the extension of the previous results to stationary metrics. (Without any Killing vector generating the time-translation symmetry it is probably impossible to perform a unique and physical reasonable particle definition.)
Another important extension of the provided formalism is given by the incorporation of the electromagnetic field
$`={\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\nu \mu }.`$ (167)
The Maxwell theory possesses primary and secondary constraints . These gauge-problems have to be solved before the quantisation and the particle definition becomes possible. One way to accomplish this – which seems to be suitable to the canonical approach – is the method of separation of variables . Nevertheless there is no obvious reason why the main results of this article should not persist. The equation of motion of the electromagnetic field is given by
$`_\mu F^{\mu \nu }={\displaystyle \frac{1}{\sqrt{g}}}_\mu \left(\sqrt{g}g^{\mu \rho }g^{\nu \sigma }_\rho A_\sigma \right){\displaystyle \frac{1}{\sqrt{g}}}_\mu \left(\sqrt{g}g^{\mu \rho }g^{\nu \sigma }_\sigma A_\rho \right)=0.`$ (168)
For a very rough estimate one may drop the second term, which is related to the longitudinal degrees of freedom. The remaining equation possesses unstable interior solutions similar to the scalar field scenario.
The investigation of the Dirac field
$`=\overline{\mathrm{\Psi }}\left({\displaystyle \frac{i}{2}}\gamma ^\mu \stackrel{}{D_\mu }m\right)\mathrm{\Psi }`$ (169)
around charged black holes creates some new kind of problems, see e.g. . Similar to the Schwinger mechanism in the semi-classical description a tunnelling process is possible. This tunnelling probability gives raise to the question of whether a stable vacuum in the quantum field theoretical treatment exists.
Having obtained a linear instability of the linear equations of motion one may ask whether the unstable behaviour persists for non-linear equations of motion including interaction terms, for instance $`\mathrm{\Phi }^4`$. One might suspect that the non-linear terms generate new stable fixed points of the equation of motion – i.e. a non-perturbative stabilisation of the black hole. However, in this case the amplitude of the field has to be located at some fixed scale while the (linear) instability exists for arbitrary large scales $`|\lambda _\mathrm{\Gamma }|`$. This might be an argument for the dominance of the unstable linear contribution in this region. In order to elucidate this point it is necessary to consider the scale behaviour of the interacting theory.
This article considers the propagation of quantised fields in a given (i.e. externally prescribed) space-time. To examine how the quantum fields influence the metric one has to deal with the back-reaction problem. Within the canonical (operator) quantisation one usually employs the renormalised expectation value of the energy-momentum tensor as the source of the Einstein equations , and within the path-integral approach one may integrate out the quantum field in order to obtain an effective action (accounting for the degrees of freedom associated with the dynamics of the space-time). However, a complete solution to this question probably requires the knowledge of the unification of general relativity and quantum field theory.
## Acknowledgement
The author is indebted to A. Calogeracos, K. Fredenhagen, I. B. Khriplovich, G. Plunien, G. Soff, and R. Verch for fruitful conversations and helpful criticism. Discussions with R. Picard concerning questions of functional analysis are also gratefully acknowledged. This work was partially supported by BMBF, DFG and GSI.
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# Dust Echos from Gamma Ray Bursts
## 1. Introduction
The relationship between Gamma Ray Bursts (GRBs) and supernovae has become increasingly interesting over the past year. Though exploding massive stars have long been considered as possible progenitors of GRBs (e.g. Woosley (1993)), no evidence existed to support these theories until observations of the afterglow of GRB 980425 suggested an association of the burst with an unusual supernova 1998bw (Galama et al. (1998), Kulkarni et al. (1998)). Later reanalysis of the optical afterglow lightcurves of two other bursts, GRB 970228 (Fruchter et al. (1999)) and GRB 980326 (Bloom et al. (1999)) showed a deviation from the power-law decline expected if the emission is due to synchrotron radiation from electrons accelerated by the blast wave. In both cases a significant excess emission was observed around $`30`$ days after the gamma-ray burst, with simultaneous reddening of the spectrum. Bloom et al. (1999), Galama et al. (1999), and Reichart (1999) attribute this excess to the emission from an underlying supernova event.
The relationship of GRBs to SN explosions is a question of great importance, since it provides a powerful clue to the fundamental nature of these objects. However, the evidence presented so far is circumstantial – the association of GRB 980425 with SN 1998bw is unproven and the excess emission seen from GRB 970228 and GRB 980326 is based upon relatively few actual measurements – and possible alternative explanations need to be seriously considered, if only to strengthen the case for the SN explanation. In this spirit, Waxman & Draine (1999) suggested that the red excess emission observed in GRB 970228 and GRB 980326 is due to dust in the vicinity of the burst progenitor absorbing and then re-radiating the optical/UV flash observed shortly after the recent GRB 990123 (Akerlof et al. (1999)) and generally attributed to the reverse shock which propagates into the fireball ejecta (Mészáros, Rees & Papathanassiou (1994), Mészáros & Rees (1997), Panaitescu & Mészáros (1998), Sari & Piran (1999)). However, the Waxman & Draine (1999) scenario has two shortcomings. Firstly, the equilibrium temperature of dust is limited to $`2300`$ K and so the emission should peak at $`2(1+z)\mu \mathrm{m}`$ (where $`z`$ is the GRB redshift), although a small amount of higher temperature emission may be produced by the dust as it is subliming. Secondly, the optical flash is so powerful that the sublimation radius lies beyond $`10`$ pc from the GRB. Thus, in this picture it is rather difficult to reproduce the observed flux in the $`0.40.8\mu \mathrm{m}`$ band with a time delay of order a few weeks.
In this letter we propose an alternative explanation, which relies on the scattering of the direct optical transient emitted in the first day by dust as the primary source of excess optical radiation. The fundamental point is that in the two observed cases, assuming isotropic emission, the fluence of the observed transient exceeds that of the reported excess and the unobserved transient is even larger if we extrapolate to earlier times. A fraction of this emission scattered from a radius where dust can outlive the optical transient should therefore produce a delayed echo. As dust absorbs selectively as well as scatters, the echo is likely to be significantly redder than the original optical transient, as reported.
In the next section, we describe our model for the dust scattering properties and then present the results in the context of the observed GRBs in §3. Implications for future tests of our scenario are discussed in §4. We assume $`h=0.6`$, $`\mathrm{\Omega }_\mathrm{M}=0.3`$, and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ so that the angular diameter distance of the GRB is $`D_\mathrm{A}=1.52`$ Gpc for $`0.5\begin{array}{c}<\hfill \\ \hfill \end{array}z\begin{array}{c}<\hfill \\ \hfill \end{array}3`$.
## 2. Dust Echos
### 2.1. Sublimation Radius
Waxman & Draine (1999) estimate that dust grains in the path of the optical/UV flash will be effectively sublimed out to a distance
$$R_{\mathrm{sub}}1(Q_{\mathrm{abs}}L_{47}/a_1)^{1/2}\mathrm{pc},$$
(2-1)
where $`Q_{\mathrm{abs}}1`$ is the absorption efficiency factor for optical/UV photons, $`L_{47}𝑑\nu L_\nu /10^{47}\mathrm{erg}\mathrm{s}^1`$ is the unbeamed luminosity of the optical transient (OT) in the $`17.5\mathrm{eV}`$ energy band, $`a_11`$ is the dust grain size in units of $`0.1\mu \mathrm{m}`$. Beyond $`R_{\mathrm{sub}}`$, only the most refractory grains, like silicates, can survive. Note that the thermal time for a typical dust particle is of order $`10^410^2\mathrm{s}`$, much shorter that the duration of the optical transient, so that we can treat grains as being in thermal equilibrium with the incident radiation, (which has a pressure $`P_{\mathrm{sub}}0.03`$ dyne cm<sup>-2</sup> at $`R_{\mathrm{sub}}`$).
The extinction properties of silicate dust particles were computed by Draine & Lee (1984) (see their Fig. 10), for a power-law distribution of particle sizes proposed by Mathis, Rumpl & Nordsieck (1977) to explain interstellar starlight extinction. Based on their results, we take the ratio of the scattering and absorption efficiency factors to be of order $`Q_{\mathrm{sc}}/Q_{\mathrm{abs}}4`$, and the average scattering angle to be $`\mathrm{cos}\theta \mu 0.5`$ for observed wavelengths $`0.21(1+z)\mu \mathrm{m}`$.
### 2.2. Source Geometry
Fig. 2.1 shows a schematic picture of the GRB environment observed at time $`t`$ after the detection of $`\gamma `$-rays. The incident optical transient emission is supposed to be limited to an interval $`\mathrm{\Delta }t^{\mathrm{OT}}\mathrm{\Delta }t_{\mathrm{ob}}^{\mathrm{OT}}/(1+z)<<t`$ after the GRB, and scattered by dust beyond $`R_{\mathrm{sub}}`$. We specialize immediately to the case when the dust is associated with an outflowing spherical wind, and the OT is isotropic. (It is straightforward to modify our formalism to accommodate other reasonable assumptions, as discussed in Madau et al. (1999).) As the dust density declines with distance as $`R^2`$, the light “echo” observed at time $`t_{\mathrm{ob}}=t(1+z)`$ will be scattered by dust concentrated in a ring located at the intersection of the sphere $`R=R_{\mathrm{sub}}`$ and the paraboloid
$$t=\frac{R}{c}(1\mu ),$$
(2-2)
where $`\mu =\mathrm{cos}\theta `$ (see Fig. 2.1). In our model, it is adequate to ignore a finite re-processing time and the radial distribution of the dust. The dust only has to survive for a time $`\mathrm{\Delta }t^{\mathrm{OT}}`$. We expect that, in practice, it will be quickly destroyed by the effects of secondary cosmic ray electrons created through electron scattering of the GRB so that the observed optical afterglow need not necessarily be subject to the same extinction as the echo.
### 2.3. Optical Scattering
The optical echo flux density $`F_{\nu _{\mathrm{ob}}}^\mathrm{E}`$, observed at frequency $`\nu _{\mathrm{ob}}=\nu (1+z)^1`$ is
$`F_{\nu _{\mathrm{ob}}}^\mathrm{E}(\nu _{\mathrm{ob}},t_{\mathrm{ob}})={\displaystyle \frac{L_\nu (\nu ,t)}{2(1+z)^3D_\mathrm{A}^2}}{\displaystyle \frac{c\mathrm{\Delta }t^{\mathrm{OT}}}{R}}{\displaystyle \frac{d𝒫^{\mathrm{sc}}}{d\mathrm{\Omega }}}(\nu ,\mu );`$ (2-3)
$`0<t_{\mathrm{ob}}<2R_{\mathrm{sub}}(1+z)/c,`$
where $`d𝒫^{\mathrm{sc}}(\nu ,\mu )/d\mathrm{\Omega }`$, is the probability of escape along the direction defined by angle $`\theta =\mathrm{cos}^1\mu `$ for a photon of frequency $`\nu `$.
From Eq. (2-3) it is clear that the only time dependence comes from the angular dependence of the escape probability, $`d𝒫^{\mathrm{sc}}/d\mathrm{\Omega }`$; $`F_{\nu _{\mathrm{ob}}}^{\mathrm{sc}}(t_{\mathrm{ob}})`$ is simply a step function for isotropic scattering. We adopt a Henyey-Greenstein function (e.g. White (1979)) to describe the differential cross section for dust scattering:
$$\frac{d\sigma }{d\mathrm{\Omega }}\frac{1\mu ^2}{[1+\mu ^22\mu \mu ]^{3/2}},$$
(2-4)
with $`\mu =0.5`$ (Draine & Lee (1984)). We then use Eq. (2-4) to compute $`d𝒫^{\mathrm{sc}}(\lambda ,\mu )/d\mathrm{\Omega }`$ numerically for a slab-like dust cloud. The results for different observer angles (w.r.t. to the slab normal vector) and two different values of the total extinction, $`\tau =\tau _{\mathrm{abs}}+\tau _{\mathrm{sc}}`$ (measured at $`\lambda =0.3\mu \mathrm{m}`$) are shown in Fig. 2.3. The differential escape probability is normalized so that the integral $`_{4\pi }\frac{d𝒫^{\mathrm{sc}}}{d\mathrm{\Omega }}(\lambda ,\mu )𝑑\mathrm{\Omega }`$ is equal to the escape probability from the dust cloud. Fig. 2.3 shows that at low optical depth $`\tau _{0.3}\tau (0.3\mu \mathrm{m})\begin{array}{c}<\hfill \\ \hfill \end{array}3`$, the echo should have a similar color to the OT whereas at larger optical depth, the echo will be much redder due to absorption.
To illustrate that our calculation of the escape probability is not overly simplified (though it ignores wavelength dependence of the functional form for $`d\sigma /d\mathrm{\Omega }`$, (White (1979)) in Fig. 2.3 we show one curve (thin long-dashed line) computed using $`d\sigma /d\mathrm{\Omega }1+2\mu +\mu ^2`$. This expression gives the same value of $`\mu `$ but is less strongly peaked at $`\mu =1`$ than Eq. (2-4). The resulting $`d𝒫^{\mathrm{sc}}(\lambda ,\mu )/d\mathrm{\Omega }`$ is very similar to what we use in our calculations.
The angular dependence of the escape probability is exhibited in Fig. 2.3 for a dust cloud with $`\tau _{0.3}=7`$ and three values of the incident photon wavelength. Note that $`d𝒫^{\mathrm{sc}}(\mu )/d\mathrm{\Omega }`$ remains relatively flat for $`\theta \begin{array}{c}<\hfill \\ \hfill \end{array}\theta _{\mathrm{sc}}20^{}`$ and decreases exponentially at larger angles.
### 2.4. Infrared Echo
Hot dust will also emit an isotropic infrared echo due to thermal emission from dust at the rapid sublimation temperature $`2300`$ K, peaking at an observed wavelength $`\lambda 2(1+z)\mu \mathrm{m}`$. Waxman & Draine (1999) argue that only the photons in the $`17.5\mathrm{eV}`$ range will contribute to dust heating. For $`\tau _{0.3}7`$ the absorption efficiency for photons in this energy range is $`>0.8`$; and moreover, such photons are likely to carry a considerable fraction of the total OT emission. Therefore the integrated infrared flux is
$$F_{\mathrm{IR}}^\mathrm{E}=\frac{L\mathrm{\Delta }t^{\mathrm{OT}}}{8\pi D_\mathrm{A}^2}\frac{c}{R}\frac{1}{(1+z)^4}.$$
(2-5)
## 3. Comparison with Observations
### 3.1. OT-Echo-Redshift Relations
Adopting our simple model of dust scattering, Eqs. (2-1, 2-2) allow us to relate the sublimation radius and OT power, $`10^{47}L_{47}`$ erg s<sup>-1</sup>, to the observed echo delay, $`t_{\mathrm{ob}}^\mathrm{E}10^6t_{\mathrm{ob}\mathrm{\hspace{0.17em}6}}^\mathrm{E}`$ s.
$`R_{\mathrm{sub}}`$ $``$ $`0.2C_1^1C_2^1t_{\mathrm{ob},6}^\mathrm{E}(1+z)^1\mathrm{pc},`$ (3-1)
$`L_{47}`$ $``$ $`0.03(1+z)^2C_1^2C_2^2(t_{\mathrm{ob},6}^\mathrm{E})^2,,`$ (3-2)
where $`C_1=(1\mu )/0.06`$ allows for beaming or characteristic scattering angles different from $`20^{}`$, and $`C_2=R/R_{\mathrm{sub}}`$ should be used if the dust is located beyond $`R_{\mathrm{sub}}`$.
For simplicity, we now suppose that the spectral index of the OT is $`\alpha 1`$. This is quite close to the spectral index of the observed afterglows. We can then use Eqs(2-3) to relate the R-band ($`0.65\mu \mathrm{m}`$) echo flux density to the escape probability
$`F_\nu ^\mathrm{E}[0.65\mu \mathrm{m}]0.4{\displaystyle \frac{dP^{\mathrm{sc}}}{d\mathrm{\Omega }}}\left({\displaystyle \frac{t_{\mathrm{ob},6}^\mathrm{E}\mathrm{\Delta }t_{\mathrm{ob},3}^{\mathrm{OT}}}{C_1C_2^2}}\right)\times `$ (3-3)
$`\left({\displaystyle \frac{D_\mathrm{A}}{1.5\mathrm{Gpc}}}\right)^2(1+z)^6\mu \mathrm{Jy},`$
where the observed duration of the optical transient is $`10^3\mathrm{\Delta }t_{\mathrm{ob},3}^{\mathrm{OT}}`$ s. Note the strong dependence on redshift which implies that accurate measurements of both the optical transient and the echo flux could lead to a fairly precise redshift prediction.
The ratios of the optical transient flux density, $`F_{\nu _{\mathrm{ob}}}^{\mathrm{OT}}=L_\nu f_{\mathrm{ns}}(4\pi )^1D_\mathrm{A}^2(1+z)^3`$, and infrared echo flux density to the optical echo flux density are likewise given by
$`{\displaystyle \frac{F_\nu ^{\mathrm{OT}}[0.65\mu \mathrm{m}]}{F_\nu ^\mathrm{E}[0.65\mu \mathrm{m}]}}3000\left({\displaystyle \frac{f_{\mathrm{ns}}}{C_1dP^{\mathrm{sc}}/d\mathrm{\Omega }}}\right)\left({\displaystyle \frac{t_{\mathrm{ob},6}^\mathrm{E}}{\mathrm{\Delta }t_{\mathrm{ob},3}^{\mathrm{OT}}}}\right),`$ (3-4)
$`{\displaystyle \frac{F_\nu ^\mathrm{E}[2(1+z)\mu \mathrm{m}]}{F_\nu ^\mathrm{E}[0.65\mu \mathrm{m}]}}0.5\left({\displaystyle \frac{dP^{\mathrm{sc}}}{d\mathrm{\Omega }}}\right)^1(1+z),`$ (3-5)
where $`f_{\mathrm{ns}}`$ is the fraction of incident OT photons, emerging unscattered from the dust cloud.
### 3.2. GRB 980326
For GRB 980326, an excess R flux $`F_{\nu _{\mathrm{ob}}}^\mathrm{E}[0.65\mu \mathrm{m}]0.4\mu \mathrm{Jy}`$ was measured a time $`t_{\mathrm{ob}}^\mathrm{E}20`$ d (Bloom et al. (1999)). If we make the simplest assumptions, $`a_1Q_{\mathrm{abs}}C_1C_21`$, then $`R_{\mathrm{sub}}0.3(1+z)^1\mathrm{pc}`$ and $`L9\times 10^{45}(1+z)^2\mathrm{erg}\mathrm{s}^1`$. Comparing the reported spectral slope, $`\alpha 2.8`$ of the putative echo to that of the afterglow ($`\alpha 0.8`$), we estimate that $`\tau _{0.3}7`$ (cf Fig. 2). This, in turn, implies that $`dP^{sc}/d\mathrm{\Omega }`$ in the observed R band $`0.2(1+z)^1`$ and $`f_{\mathrm{ns}}0.05(1+z)^4`$ (see Fig. 2.3). We can then use Eq.(3-3) to deduce that $`\mathrm{\Delta }t_{\mathrm{ob},3}^{\mathrm{OT}}3(1+z)^7`$ and $`F_{\nu _{\mathrm{ob}}}^{\mathrm{OT}}[0.65\mu \mathrm{m}]200(1+z)^{10}\mu \mathrm{Jy}`$. If $`z0.4`$, then the energy associated with the first optical measurement of the afterglow ($`F_{\nu _{\mathrm{ob}}}^{\mathrm{OT}}[0.65\mu \mathrm{m}]10\mu \mathrm{Jy}`$ after 0.5 d), suffices to account for the observed excess after 20 d as a dust echo. If $`z>0.4`$, then the optical transient would have had to be present and create a larger fluence at earlier times. This is not unreasonable as the OT flux was measured to satisfy $`F^{\mathrm{OT}}t^2`$. In view of the large number of simplifying assumptions that we have made, this estimate can only be regarded as illustrative. However it suffices to demonstrate that dust scattering is consistent with all of the available data.
### 3.3. GRB 970228
A somewhat similar story can be told for GRB 970228, where the redshift, $`z=0.695`$, is known (Djorgovski et al. (1999)). The earliest R-band measurement is $`30\mu \mathrm{Jy}`$ 0.7 d after the GRB; and after $`30`$ d there red excess flux $`0.3\mu \mathrm{Jy}`$ was observed, with the spectral slope ($`\alpha 3.0`$) very similar to that seen in GRB 980326 (Galama et al. (1999)). For this object, again, within the uncertainties, the fluence measured in the first stages of the optical transient is sufficient to account for the energy in the optical excess.
### 3.4. Dust Origin
In both examples above, the mass of dust required to produce an optical depth $`\tau _{0.3}7`$ with our simplest assumptions and assuming that it is spherically symmetrically distributed with respect to the GRB is $`0.1`$ M. This amount of dust could form in an expanding high-metallicity wind associated with an earlier stage in the evolution of the GRB progenitor as we have assumed in our simple model. Alternatively the dust might be associated with a molecular cloud if GRBs are associated with massive star formation or a molecular torus should they be located in obscured galactic nuclei.
## 4. Discussion
In this letter we present an alternative explanation for the reddened excess emission observed in GRB 970228 and GRB 980326, which we attribute to dust scattering of the early-time, afterglow emission. This scenario is predictive enough to be confirmed or ruled out with observations of future GRBs. In particular, in contrast to the supernova explanation (Bloom et al. (1999); Galama et al. (1999); Reichart (1999)), if the excess emission is due to dust scattering, then its properties will depend on the luminosity of the optical transient. HETE II (http://space.mit.edu/HETE/) scheduled to be launched in early 2000 and Swift (http://swift.gsfc.nasa.gov/homepage.html), scheduled for 2003, should provide real-time localization of GRB X-ray afterglows with sufficient precision to permit faster follow-up and better measurements of its total fluence. Infrared observations may discover the expected thermal emission from hot subliming dust (cf Waxman & Draine (1999)). In fact dust emission might be the correct explanation for the “near-IR” bump seen in the spectrum of the GRB 991216 afterglow (Frail et al. (2000)). Note that as most GRBs are at redshifts $`\begin{array}{c}>\hfill \\ \hfill \end{array}0.5`$, 3 $`\mu \mathrm{m}`$, (as as opposed to the more common 2 $`\mu \mathrm{m}`$) photometry may be necessary to see this emission.
In those GRBs, where it is also possible to measure a redshift, the the simplest model of dust-scattering is over-contrained and therefore refutable. Beaming and dust inhomogeneity introduce additional uncertainty but such models may also be excludable. For example, if ROTSE (Akerlof et al. (1999)) were to detect another optical flash in a GRB as luminous as that seen in GRB 990123, which had an isotropic luminosity $`L10^{51}`$ erg s<sup>-1</sup>, then dust should be physically sublimed out to a distance $`R_{\mathrm{sub}}100`$ pc along the line of sight. Unreasonably large beaming would then be required to explain a dust echo with a delay of only a few weeks. Alternatively, if the radio light curve in an afterglow tracked the optical light curve, then this would be incompatible with both dust scattering and a supernova.
A further prediction of the dust echo model is that, unless the dust and OT are both arranged axisymmetrically with respect to the line of sight, we expect there to be linear polarization associated with dust echos and this may be measurable in bright examples. (1.7 percent polarization has been reported in the optical transient associated with GRB 990510 by Covino et al. (1999) but this is unlikely to be due to scattering.)
In conclusion, we have demonstrated that dust scattering can account for the excess optical emission observed in the afterglows of two GRBs as an alternative to an underlying supernova explosion. Future observations should be able to rule out or confirm this explanation.
We thank J. Bloom for helpful comments. This work was supported by NASA grant 5-2837, NSF grant AST 99-00866 and a Chandra Postdoctoral Fellowship grant #PF8-10002 awarded by the Chandra X-Ray Center, which is operated by the SAO for NASA under contract NAS8-39073.
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# 1 Introduction
## 1 Introduction
M-theory compactified on
$$O_7=X_6\times S^1/𝐙_2,$$
(1.1)
where $`X_6`$ is a Calabi-Yau three-fold, leads to a four-dimensional theory with $`N=1`$ local supersymmetry. At present, our knowledge of the resulting effective supergravity is based on the few aspects of M-theory which are quantitatively understood and on the small-orbifold limit<sup>1</sup><sup>1</sup>1The limit in which the length of the $`S^1/𝐙_2`$ orbifold direction is small compared to the size of the Calabi-Yau manifold. which is the perturbative heterotic $`E_8\times E_8`$ string compactified on the Calabi-Yau space $`X_6`$. In the low-energy limit, M-theory information can be organized as an expansion in powers of the eleven-dimensional gravitational constant $`\kappa _{11}`$ . The lowest order $`\kappa _{11}^2`$ is eleven-dimensional supergravity . In the case of a compactification on $`S^1/𝐙_2`$ only, the next orders in $`\kappa _{11}`$ are known to include orbifold plane contributions (super-Yang-Mills terms) as well as gauge and gravitational anomaly-cancelling terms .
Similarly, the effective four-dimensional supergravity<sup>2</sup><sup>2</sup>2Which is a low-energy description. can be formulated as an expansion in the four-dimensional gravitational constant $`\kappa `$, even if a more common choice suggested by string theory is to use the dilaton as expansion parameter. The two options coincide provided the appropriate supersymmetric description of the dilaton is adopted in a Wilson effective Lagrangian formulation. The lowest order $`\kappa ^2`$ is the $`O_7`$ truncation of eleven-dimensional supergravity. The next order includes super-Yang-Mills, charged matter kinetic and superpotential contributions. Then come sigma-model anomaly-cancelling terms contributing in particular to gauge threshold corrections. These first corrections to the low-energy limit of compactified M-theory are identical to those obtained from heterotic compactifications on Calabi-Yau. This is certainly expected since the information content is identical. The literature gives a detailed description of these results, with particular attention paid to the ‘strong-coupling’ heterotic limit in which the size of the Calabi-Yau space is smaller than the orbifold length, supersymmetry breaking by gaugino condensation and non-standard embeddings .
Our goal in this paper is to provide a derivation of the effective supergravity which explicitly relates four-dimensional supergravity statements with M-theory aspects like Bianchi identities modified at singularities and anomaly cancellation. We reformulate these basic facts of M-theory on $`O_7`$ directly in terms of four-dimensional supermultiplets and equations. For instance, Bianchi identities from M-theory are promoted to field equations, as constraints on multiplets which are massless modes of M-theory bulk fields. With this formulation, we expect to obtain a clean, direct derivation of the effective supergravity suitable to cases more subtle than the universal modes of M-theory on $`O_7`$. A first use of our formalism will be the coupling of five-brane moduli supermultiplets in the $`O_7`$ case .
The organization of the paper is as follows. In Section 2, we establish our basic supergravity formulation, using the bosonic bulk dynamics as starting point. The resulting Lagrangian is the lowest order in the $`\kappa `$-expansion. It is essentially defined by a dynamical Lagrangian involving tensor fields supplemented by Bianchi identities which are field equations of the theory. We discuss in detail the bosonic component expansion of the supermultiplet action, the question of the gravity frame and the generation of a superpotential. In Section 3, we introduce the next order corrections: gauge multiplets and charged matter contributions. We show that their introduction is controlled by a simple modification of the four-dimensional Bianchi identities, in analogy with the appearance of $`𝐙_2`$ fixed planes contributions in the M-theory Bianchi identities. Section 4 discusses anomaly-cancelling terms. We begin by modifying the four-dimensional effective supergravity by adding terms similar to those appearing for gauge threshold corrections in $`(2,2)`$ compactifications of the heterotic string. These modifications can be formulated in terms of our particular set of multiplets directly related to M-theory bulk degrees of freedom. We then directly compute these anomaly-cancelling terms by Kaluza-Klein reduction of the ten-dimensional Green-Schwarz counterterms arising from M-theory on $`S^1/𝐙_2`$. Section 5 gives some final comments and conclusions, and an appendix contains our notations and conventions.
## 2 Four-dimensional effective M-theory <br>supergravities: bulk dynamics
Our concern is compactifications of M-theory to four dimensions preserving $`N=1`$ supersymmetry. Or compactifications in which supersymmetry would break spontaneously or dynamically at solutions of the effective field equations. As a consequence, the light (massless) modes can be described by a local effective $`N=1`$ supergravity Lagrangian, to be understood in the sense of Wilson.
These solutions are not as well understood as, for instance, Calabi-Yau compactifications of heterotic strings and it will prove useful to give a precise description of the aspects which are better known. We will then focus on a precise description of two aspects which are of importance in compactifications of M-theory. Firstly, we will use a supersymmetric description of the Bianchi identities verified by antisymmetric tensors, as arising from M-theory or higher-dimensional supergravities. This procedure allows to avoid multiplet ambiguities arising when duality transformations are performed at the bosonic level only. Secondly, we will use a formulation which leaves explicitly the choice of gravity frame open. This can be a relevant issue since an expansion, perturbative or not, is performed around a gravitational background which selects a gravity frame. Standard Poincaré supergravity is usually written in the Einstein frame, in which the gravitational Lagrangian is $`\frac{1}{2\kappa ^2}eR`$. Corrections to the lowest order effective action, which includes this gravitational term, induce in general (but not always) corrections to the gravitational Lagrangian which affect the Einstein frame condition.
Keeping open the choice of gravity frame suggests to use superconformal supergravity. And displaying Bianchi identities explicitly in an effective Lagrangian requires using chiral, linear or vector supermultiplets with constraints.
The purpose of this first section is to establish our procedure by considering the well-known ‘bulk dynamics’, which follows from $`O_7`$ compactification of eleven-dimensional supergravity.
### 2.1 Superconformal formalism
We use the superconformal formulation of $`N=1`$ supergravity with a chiral compensating multiplet $`S_0`$ to generate Poincaré theories by gauge fixing<sup>3</sup><sup>3</sup>3This is ‘old minimal’ Poincaré supergravity .. Its conformal and chiral weights are taken as $`w=1`$ and $`n=1`$. This formalism is particularly convenient to keep control of a change of frame<sup>4</sup><sup>4</sup>4Mostly the so-called Einstein or string frames. which corresponds to a different Poincaré gauge condition applied on the modulus of the scalar compensator $`z_0`$, which fixes dilatation symmetry. Up to two derivatives, a supergravity Lagrangian is written as<sup>5</sup><sup>5</sup>5Except otherwise mentioned, our notation is as in ref. where also reference to the original literature can be found.
$$=\left[S_0\overline{S}_0\mathrm{\Phi }\right]_D+\left[S_0^3W\right]_F+\frac{1}{4}\left[f_{ab}𝒲^a𝒲^b\right]_F.$$
(2.1)
The symbols $`[\mathrm{}]_D`$ and $`[\mathrm{}]_F`$ denote respectively the invariant $`D`$\- and $`F`$-density formulæ given by (all fermion contributions are omitted)
$$\begin{array}{ccc}\hfill [𝒱]_D& =& e(d+\frac{1}{3}cR),\hfill \\ \multicolumn{3}{c}{}\\ \hfill [𝒮]_F& =& e(f+\overline{f}),\hfill \end{array}$$
(2.2)
where $`𝒱`$ is a vector multiplet with components $`(c,\chi ,m,n,b_\mu ,\lambda ,d)`$ and $`𝒮`$ a chiral multiplet with components $`(z,\psi ,f)`$. The real vector multiplet $`\mathrm{\Phi }`$ with zero weights is a function (in the sense of tensor calculus) of the multiplets present in the theory, including in general the compensating multiplet. The holomorphic function $`W`$ of the chiral multiplets is the superpotential. The last term contributes to gauge kinetic terms (the chiral multiplet $`𝒲`$ is the gauge field strength for the gauge multiplets) and involves a holomorphic gauge kinetic function $`f_{ab}`$ of the chiral multiplets<sup>6</sup><sup>6</sup>6 Gauge kinetic terms may also arise from the first term.. Expression (2.1) provides the most general supergravity Lagrangian up to terms with more than two derivatives and up to terms which would contribute to kinetic terms in a fermionic background only . Besides $`S_0`$ and $`𝒲`$, we will use chiral multiplets with zero weights and neither $`W`$ nor $`f_{ab}`$ will depend on the compensator.
The chiral $`U(1)`$ symmetry of the superconformal algebra can be extended to
$$S_0,W,\mathrm{\Phi }\mathrm{\Lambda }S_0,\mathrm{\Lambda }^3W,(\mathrm{\Lambda }\overline{\mathrm{\Lambda }})^1\mathrm{\Phi }|_{S_0\mathrm{\Lambda }S_0},$$
(2.3)
with an arbitrary chiral multiplet $`\mathrm{\Lambda }`$. This symmetry is at the origin of Kähler invariance of Poincaré supergravity. The last transformation suggests that $`\mathrm{log}\mathrm{\Phi }`$ transforms as the corresponding gauge connection. Choosing $`\mathrm{\Lambda }=W^{1/3}`$ eliminates the superpotential except if it vanishes. One can then use a $`U(1)`$/Kähler gauge fixing in which the supergravity Lagrangian (2.1) reads
$$=\left[S_0\overline{S}_0\mathrm{\Phi }\right]_D+c\left[S_0^3\right]_F+\frac{1}{4}\left[f_{ab}𝒲^a𝒲^b\right]_F,$$
(2.4)
with an arbitrary constant $`c`$ as superpotential and two arbitrary functions $`\mathrm{\Phi }`$ and $`f_{ab}`$.
The real function $`\mathrm{\Phi }`$ depends on matter multiplets, which will be either chiral multiplets like the Calabi-Yau universal modulus $`T`$, or real linear multiplets (with weights $`w=2`$, $`n=0`$) like the dilaton multiplet in the version of the theory with an antisymmetric tensor, or real vector multiplets ($`n=0`$, $`w`$ arbitrary) like the multiplet of gauge potentials ($`w=0`$). Vector multiplets will appear as essential ingredients in the effective description of $`(N=1)`$-preserving M-theory five-branes .
### 2.2 Bulk Lagrangian
The lowest order (in the $`\kappa `$ expansion) effective four-dimensional supergravity of M-theory compactified on $`O_7`$ describes Kaluza-Klein massless modes of eleven-dimensional supergravity. It is the $`S^1/𝐙_2`$ truncation of eleven-dimensional supergravity on a Calabi-Yau three-fold.
In the version given by Cremmer, Julia and Scherk (CJS) , the Lagrangian of eleven-dimensional supergravity can be written as<sup>7</sup><sup>7</sup>7We use the flat space-time metric $`(,+,+,\mathrm{},+)`$. The gravitational constant $`\kappa _{11}`$ has dimension (mass)<sup>-9/2</sup> and the three-index tensor field $`C_{MNP}`$ is dimensionless. Notations and conventions are defined in an appendix.
$$\begin{array}{ccc}\hfill e^1_{\mathrm{CJS}}& =& \frac{1}{2\kappa _{11}^2}[R\frac{1}{24!}G_{M_1M_2M_3M_4}G^{M_1M_2M_3M_4}\hfill \\ \multicolumn{3}{c}{}\\ & & \frac{1}{6}\frac{1}{4!4!3!}e^1ϵ^{M_1\mathrm{}M_{11}}G_{M_1M_2M_3M_4}G_{M_5M_6M_7M_8}C_{M_9M_{10}M_{11}}]\hfill \\ \multicolumn{3}{c}{}\\ & & +\mathrm{fermionic}\mathrm{terms}.\hfill \end{array}$$
(2.5)
Omitting all fields related to the detailed geometry of the Calabi-Yau manifold, the particle content of the four-dimensional theory is the $`N=1`$ supergravity multiplet, with metric tensor $`g_{\mu \nu }`$, and matter multiplets including (on-shell) four bosons and four fermions. Two bosons are scalars and correspond to the dilaton and the ‘universal modulus’ of the Calabi-Yau space, the massless volume mode. Two bosons are Kaluza-Klein modes of the four-form field $`G`$, with Bianchi identity $`dG=0`$. Explicitly, these fields and their Bianchi identities can be written as
$$\begin{array}{ccc}G_{\mu \nu \rho 4},\hfill & & _{[\mu }G_{\nu \rho \sigma 4]}=0,\hfill \\ \multicolumn{3}{c}{}\\ G_{\mu j\overline{k}4}=iT_\mu \delta _{j\overline{k}},\hfill & & _{[\mu }T_{\nu ]}=0.\hfill \end{array}$$
(2.6)
It will prove useful to identify these fields with the vector components of two real vector multiplets $`V`$ and $`V_T`$, and to impose the Bianchi identities as field equations using the appropriate multiplets as Lagrange multipliers. The bulk supergravity Lagrangian takes then the form
$$_\mathrm{B}=\left[\frac{1}{\sqrt{2}}(S_0\overline{S}_0V_T)^{3/2}V^{1/2}(S+\overline{S})V+L_TV_T\right]_D.$$
(2.7)
The vector superfield $`V`$ (weights $`w=2`$, $`n=0`$) includes in its components the vector field $`v^\mu ϵ^{\mu \nu \rho \sigma }G_{\nu \rho \sigma 4}`$ for which the Bianchi identity is $`_\mu v^\mu =0`$. This condition is a component of the (super)field equation of the chiral $`S`$ ($`w=n=0`$) which imposes $`V=L`$, a real linear multiplet ($`w=2`$, $`n=0`$). Secondly, the vector multiplet $`V_T`$ ($`w=n=0`$) includes in its components $`T_\mu `$. The supersymmetric extension of its Bianchi identity $`_{[\mu }T_{\nu ]}=0`$ is enforced by the real linear multiplet $`L_T`$, which implies $`V_T=T+\overline{T}`$, with a chiral weightless multiplet $`T`$. The usefulness of obtaining Bianchi identities via field equations will become apparent with the introduction of higher orders in the $`\kappa `$ expansion. At this stage of the discussion however, it gives a formulation of the familiar duality relating scalars and antisymmetric tensors or, for superfields, chiral and linear multiplets.
Solving in Eq. (2.7) for the Lagrange multipliers $`S`$ and $`L_T`$ leads to the ‘standard form’ of the bulk four-dimensional Lagrangian
$$_{\mathrm{B},\mathrm{l}}=\frac{1}{\sqrt{2}}\left[\left(S_0\overline{S}_0e^{\widehat{K}/3}\right)^{3/2}L^{1/2}\right]_D,$$
(2.8)
or, as defined in Eq. (2.1),
$$\mathrm{\Phi }=\left(\frac{2L}{S_0\overline{S}_0}\right)^{1/2}e^{\widehat{K}/2}.$$
(2.9)
The Kähler potential for the volume modulus $`T`$ is
$$\widehat{K}=3\mathrm{log}(T+\overline{T}).$$
(2.10)
We will see again below that this standard form is naturally obtained by direct reduction of the CJS version of eleven-dimensional supergravity on $`O_7`$. Clearly, theory (2.8) is also the Calabi-Yau truncation of ten-dimensional $`N=1`$ pure supergravity . Notice that $`\widehat{K}`$ can be regarded as the Kähler connection for symmetry (2.3), with transformation
$$\widehat{K}\widehat{K}+3\mathrm{log}\mathrm{\Lambda }+3\mathrm{log}\overline{\mathrm{\Lambda }},$$
(2.11)
such that $`S_0\overline{S}_0e^{\widehat{K}/3}`$ is chiral/Kähler invariant.
A theory with a linear multiplet is in principle dual to an equivalent Lagrangian with the linear multiplet replaced by a chiral one. In our case, solving for $`V`$ and $`L_T`$ in expression (2.7) leads to
$$\begin{array}{ccc}\hfill _{\mathrm{B},\mathrm{c}}& =& \frac{3}{2}\left[S_0\overline{S}_0e^{K/3}\right]_D,\hfill \\ \multicolumn{3}{c}{}\\ \hfill K& =& \mathrm{log}(S+\overline{S})+\widehat{K}=\mathrm{log}(S+\overline{S})3\mathrm{log}(T+\overline{T}).\hfill \end{array}$$
(2.12)
This familiar chiral form is not the most useful as long as one insists on the four-dimensional translation of eleven-dimensional Bianchi identities.
Notice that one can obtain another equivalent form of the Lagrangian (2.7) by choosing to solve for $`S`$ and $`V_T`$. In this case, the Calabi-Yau modulus is described by a linear multiplet $`L_T`$. This form will not be useful since it is known that one-loop string corrections in general break the chiral-linear duality for this modulus: they involve holomorphic functions of $`T`$ in a $`F`$-density which are intrinsically chiral . Finally, there is an obstruction when trying to solve for $`V`$ and $`V_T`$ and one cannot write an expression in terms of the chiral $`S`$ and the linear $`L_T`$.
Before turning to explicit component expressions, we should discuss the choice of Poincaré frame, and introduce the expansion in the four-dimensional gravity coupling $`\kappa `$, which effectively corresponds to the low-energy expansion of M-theory in powers of the eleven-dimensional gravitational constant $`\kappa _{11}`$.
#### 2.2.1 Einstein frame
To gauge-fix dilatations, we impose as usual a condition on the Einstein term appearing in the superconformal supergravity Lagrangian. According to the component expression for the $`D`$-density and the tensor calculus of superconformal multiplets , the Einstein term included in $`[S_0\overline{S}_0\mathrm{\Phi }]_D`$ is
$$_\mathrm{E}=\frac{1}{2}eR\left[\frac{2}{3}z_0\overline{z}_0\left(\mathrm{\Phi }\frac{1}{2}\underset{i}{}w_iC_i\frac{\mathrm{\Phi }}{C_i}\right)\right],$$
(2.13)
where $`w_i`$ is a Weyl weight, the sum is taken over the linear ($`w_i=2`$) and vector ($`w_i`$ arbitrary) multiplets, and $`z_0`$, $`\mathrm{\Phi }`$ and $`C_i`$ are the lowest, scalar components of respectively $`S_0`$, $`\mathrm{\Phi }`$ and of the vector or linear multiplets<sup>8</sup><sup>8</sup>8We use in general the same notation for the lowest component of $`\mathrm{\Phi }`$ and the multiplet itself. Also, in expression (2.13), we use the component expansion of vector multiplets indicated below \[Eq.(2.23)\], which differs in its highest component from ref. ..
Applied to the bulk Lagrangian (2.7), expression (2.13) leads to
$$_\mathrm{E}=\frac{1}{2}eR\left[(z_0\overline{z}_0C_T)^{3/2}(2C)^{1/2}\right].$$
(2.14)
As they should, the terms introduced to impose Bianchi identities do not contribute. The Einstein frame is then selected by the dilatation gauge condition
$$\kappa ^2=(z_0\overline{z}_0C_T)^{3/2}(2C)^{1/2}.$$
(2.15)
It will be convenient to introduce the (composite) real vector multiplet
$$\mathrm{{\rm Y}}=(S_0\overline{S}_0V_T)^{3/2}(2V)^{1/2},$$
(2.16)
with conformal weight two. In the Poincaré theory and in the Einstein frame, its lowest component is precisely equal to $`\kappa ^2`$. With this definition, the bulk Lagrangian becomes simply
$$_\mathrm{B}=\left[\mathrm{{\rm Y}}(S+\overline{S})V+L_TV_T\right]_D,$$
(2.17)
and the equation of motion for $`V`$ (the chiral-linear duality equation)
$$2V(S+\overline{S})=\mathrm{{\rm Y}}$$
(2.18)
indicates that the Einstein Lagrangian also reads
$$_\mathrm{E}=(2CRes)eR.$$
In the Einstein frame (Planck units), $`Res=(4\kappa ^2C)^1`$.
Eq. (2.18) is compatible with the standard relation of heterotic strings<sup>9</sup><sup>9</sup>9 We use $`\mathrm{}`$ for a background value. $`2\kappa ^2Res=\alpha ^{}`$ if one identifies $`2C=1/\alpha ^{}`$. This equation defines string units, in which
$$_\mathrm{E}=\frac{e^{2\phi }}{\alpha ^{}}eR,$$
(2.19)
with a dilaton given by $`e^{2\phi }=Res`$.
#### 2.2.2 Modified Bianchi identities and $`\kappa `$-expansion
Compactification of M-theory on $`S^1/𝐙_2`$ is commonly discussed in an expansion in powers of $`\kappa _{11}`$. Compactification on $`O_7`$ can similarly be formulated with $`\kappa `$ as expansion parameter. In the upstairs version, Bianchi identities are modified at the ten-dimensional planes fixed by $`S^1/𝐙_2`$. Suppose now that we modify the four-dimensional supersymmetric Bianchi identities of the bulk Lagrangian in the following way:
$$_\mathrm{B}=\left[\mathrm{{\rm Y}}(S+\overline{S})(V+\mathrm{\Delta }_V)+L_T(V_T+\mathrm{\Delta }_T)\right]_D,$$
(2.20)
with two composite vector multiplets $`\mathrm{\Delta }_V`$ ($`w=2`$, $`n=0`$) and $`\mathrm{\Delta }_T`$ ($`w=n=0`$). Solving for the Lagrange multipliers now leads to
$$V=L\mathrm{\Delta }_V,V_T=T+\overline{T}\mathrm{\Delta }_T.$$
The Lagrangian to first order in these modifications reads then
$$\begin{array}{ccc}\hfill & =& _\mathrm{B}\left[\frac{\mathrm{{\rm Y}}}{2V}\mathrm{\Delta }_V\frac{3}{2}\frac{\mathrm{{\rm Y}}}{V_T}\mathrm{\Delta }_T\right]_D\hfill \\ \multicolumn{3}{c}{}\\ & =& _\mathrm{B}\left[(S+\overline{S})\mathrm{\Delta }_V\frac{3}{2V_T}(\mathrm{{\rm Y}}\mathrm{\Delta }_T)\right]_D,\hfill \end{array}$$
(2.21)
with $`V`$ and $`V_T`$ respectively replaced by $`L`$ and $`T+\overline{T}`$ in these expressions. The multiplets $`\mathrm{\Delta }_V`$ and $`\mathrm{{\rm Y}}\mathrm{\Delta }_T`$, with ‘canonical’ dimension $`w=2`$, appear at order $`\mathrm{{\rm Y}}^0\kappa ^0`$, in comparison with bulk terms of order $`\mathrm{{\rm Y}}\kappa ^2`$. This is the relation with the expansion in powers of $`\kappa _{11}`$ of M-theory in the low-energy limit. In M-theory compactification, the multiplets $`\mathrm{\Delta }_V`$ and $`\mathrm{\Delta }_T`$ can thus be obtained either by considering the modified Bianchi identities on $`O_7`$, formulated as in Eq. (2.20), or from corrections to the Lagrangian of eleven-dimensional supergravity on $`O_7`$, as in expression (2.21).
#### 2.2.3 Component expressions
To analyze the Lagrangian (2.7), we will need to define some notations. Since we will only explicitly consider the bosonic sector of the theory, all fermions in the $`N=1`$ supermultiplets will be omitted. Since also we are concerned with Poincaré supergravity, we will immediately gauge-fix the superconformal symmetries not contained in $`N=1`$ Poincaré supersymmetry, with one exception, dilatation symmetry: we want to keep the freedom of a frame choice as explicit as possible. These assumptions imply in particular that superconformal covariant derivatives reduce in general to $`D_\mu ^c\varphi =D_\mu \varphi \frac{i}{2}nA_\mu \varphi `$ for a complex field with chiral weight $`n`$ and to
$$\mathrm{}^c\phi =\mathrm{}\phi +\frac{1}{6}w\phi R,$$
(2.22)
if $`w`$ is the Weyl weight and $`\phi `$ real. The gauge boson $`A_\mu `$ of chiral $`U(1)`$ symmetry is auxiliary and $`D_\mu `$ and $`\mathrm{}`$ would be covariantized with respect to Poincaré symmetries.
We will use the following components for the various superconformal multiplets appearing in Lagrangian (2.7):
$$\begin{array}{ccc}\hfill V& =& (C,0,H,K,v_\mu ,0,d\mathrm{}C\frac{1}{3}CR),\hfill \\ \multicolumn{3}{c}{}\\ \hfill V_T& =& (C_T,0,H_T,K_T,T_\mu ,0,d_T\mathrm{}C_T),\hfill \\ \multicolumn{3}{c}{}\\ \hfill S& =& (s,0,f,if,i_\mu s,0,0),\hfill \\ \multicolumn{3}{c}{}\\ \hfill L_T& =& (\mathrm{}_T,0,0,0,t_\mu ,0,\mathrm{}\mathrm{}_T\frac{1}{3}\mathrm{}_TR),t_\mu =\frac{e}{2}ϵ_{\mu \nu \rho \sigma }^\nu t^{\rho \sigma },\hfill \\ \multicolumn{3}{c}{}\\ \hfill S_0& =& (z_0,0,f_0,if_0,iD_\mu ^cz_0,0,0).\hfill \end{array}$$
(2.23)
The role of the Lagrange multipliers $`S`$ and $`L_T`$ follows from
$$\begin{array}{ccc}\hfill e^1[(S+\overline{S})V]_D& =& 2(^\mu Ims)v_\mu 2^\mu (Res_\mu C)+2dRes\hfill \\ \multicolumn{3}{c}{}\\ & & f(HiK)\overline{f}(H+iK)\hfill \\ \multicolumn{3}{c}{}\\ & =& 2Ims^\mu v_\mu +2dResf(HiK)\overline{f}(H+iK)\hfill \\ \multicolumn{3}{c}{}\\ & & +\mathrm{derivative},\hfill \\ \multicolumn{3}{c}{}\\ \hfill e^1[L_TV_T]_D& =& \mathrm{}_T(d_T\mathrm{}C_T)\frac{e}{2}ϵ_{\mu \nu \rho \sigma }(^\mu T^\nu )t^{\rho \sigma }+\mathrm{derivative}.\hfill \end{array}$$
(2.24)
Solving for the components of $`S`$ leads to $`^\mu v_\mu =d=H=K=0`$, and $`V`$ is a linear multiplet. Solving for the components of $`L_T`$ leads to $`d_T\mathrm{}C_T=_{[\mu }T_{\nu ]}=0`$, and $`V_T`$ can be written as $`T+\overline{T}`$, with a chiral multiplet $`T`$ (zero weights)<sup>10</sup><sup>10</sup>10 The components are: $`C_T=2ReT`$, $`T_\mu =2_\mu ImT`$, $`H_T=2Ref_T`$, $`K_T=2Imf_T`$..
Since one can always write $`v_\mu =\frac{e}{6}ϵ_{\mu \nu \rho \sigma }v^{\nu \rho \sigma }`$, we have generated with $`Ims`$ and $`t_{\mu \nu }`$ the Bianchi identities
$$_{[\mu }v_{\nu \rho \sigma ]}=_{[\mu }T_{\nu ]}=0.$$
A modification of these Bianchi identities, as induced by $`S^1/𝐙_2`$ compactification or by five-brane couplings will then be phrased as a modification of the supermultiplets appearing multiplied by $`S+\overline{S}`$ or $`L_T`$ in Eqs. (2.24).
To complete the identification of the four-dimensional supergravity (2.7) with the modes of eleven-dimensional supergravity we need its complete bosonic expansion, which after solving for the chiral $`U(1)`$ auxiliary field $`A_\mu `$ reads:
$$\begin{array}{ccc}\hfill e^1_\mathrm{B}& =& \frac{1}{2}\mathrm{{\rm Y}}R+\frac{1}{4}\mathrm{{\rm Y}}C^2v_\mu v^\mu \frac{3}{4}\mathrm{{\rm Y}}C_T^2T_\mu T^\mu t_\mu T^\mu +2Ims_\mu v^\mu \hfill \\ \multicolumn{3}{c}{}\\ & & \frac{1}{4}\mathrm{{\rm Y}}C^2(_\mu C)(^\mu C)\frac{3}{4}\mathrm{{\rm Y}}C_T^2(_\mu C_T)(^\mu C_T)\hfill \\ \multicolumn{3}{c}{}\\ & & +d(\frac{1}{2}\mathrm{{\rm Y}}C^12Res)+(d_T\mathrm{}C_T)(\mathrm{}_T\frac{3}{2}\mathrm{{\rm Y}}C_T^1)\hfill \\ \multicolumn{3}{c}{}\\ & & +\frac{1}{2}(_\mu \mathrm{{\rm Y}})[^\mu \mathrm{log}C+^\mu \mathrm{log}\mathrm{{\rm Y}}]\hfill \\ \multicolumn{3}{c}{}\\ & & +e^1_{\mathrm{AUX}.}+\mathrm{derivative},\hfill \end{array}$$
(2.25)
where $`\mathrm{{\rm Y}}=(z_0\overline{z}_0C_T)^{3/2}(2C)^{1/2}`$, and
$$\begin{array}{ccc}\hfill e^1_{\mathrm{AUX}.}& =& \frac{1}{4}\mathrm{{\rm Y}}C^2(H+iK)(HiK)+\frac{3}{4}\mathrm{{\rm Y}}C_T^2(H_T+iK_T)(H_TiK_T)\hfill \\ \multicolumn{3}{c}{}\\ & & +f(HiK)+\overline{f}(H+iK).\hfill \end{array}$$
(2.26)
The last equality is obtained after solving for the $`f_0`$ component of $`S_0`$ <sup>11</sup><sup>11</sup>11Note that the result would be different with a superpotential..
The above component expansion of the bosonic Lagrangian is useful because it explicitly displays the dependence on the gauge choice for dilatation symmetry. The gravitational constant is the field-dependent quantity $`\mathrm{{\rm Y}}`$. Choosing a Poincaré frame amounts to impose the value of this quantity, and to use this condition to eliminate $`z_0`$. Notice also that the choice of the phase of $`z_0`$ is a gauge condition for the chiral internal $`U(1)`$ superconformal symmetry. With the exception of $`z_0`$, all bosons are $`U(1)`$-neutral. As a consequence, the bosonic Lagrangian only depends on the modulus of $`z_0`$.
Taking the Einstein frame, $`\mathrm{{\rm Y}}=\kappa ^2`$, and solving for the components of $`S`$ and $`L_T`$ leads to
$$\begin{array}{ccc}\hfill e^1_\mathrm{B}& =& \frac{1}{2\kappa ^2}R\frac{1}{4\kappa ^2}C^2[(_\mu C)(^\mu C)v_\mu v^\mu ]\hfill \\ \multicolumn{3}{c}{}\\ & & \frac{3}{4\kappa ^2}C_T^2[(_\mu C_T)(^\mu C_T)+T_\mu T^\mu ],\hfill \end{array}$$
(2.27)
with $`v_\mu =\frac{e}{2}ϵ_{\mu \nu \rho \sigma }^\nu b^{\rho \sigma }`$ since $`V`$ is a linear multiplet, $`C_T=2ReT`$ and $`T_\mu =2_\mu ImT`$ since $`V_T=T+\overline{T}`$. This Lagrangian is to be compared with the reduction of the CJS version of eleven-dimensional supergravity (2.5). The $`𝐙_2`$ orbifold projection eliminates all states which are odd under $`x^4x^4`$, and since we disregard massless modes related to the detailed Calabi-Yau geometry, the reduction of the $`D=11`$ space-time metric is
$$g_{MN}=\left(\begin{array}{ccc}e^\gamma e^{2\sigma }g_{\mu \nu }& 0& 0\\ 0& e^{2\gamma }e^{2\sigma }& 0\\ 0& 0& e^\sigma \delta _{i\overline{j}}\end{array}\right).$$
(2.28)
The $`SU(3)`$–invariant tensor $`\delta _{i\overline{j}}`$ refers to complex coordinates on the Calabi-Yau space. The surviving components of the four-index tensor $`G_{MNPQ}`$ are only $`G_{\mu \nu \rho 4}`$ and $`G_{\mu i\overline{j}4}`$, with
$$G_{\mu \nu \rho 4}=3_{[\mu }C_{\nu \rho ]4},G_{\mu i\overline{j}4}=_\mu C_{i\overline{j}4},C_{i\overline{j}4}=ia(x)\delta _{i\overline{j}},$$
and the four-dimensional Lagrangian for these fields reads
$$\begin{array}{ccc}\hfill e^1_{\mathrm{CJS}}& =& \frac{1}{2\kappa ^2}R\frac{1}{4\kappa ^2}\left[9(_\mu \sigma )(^\mu \sigma )+\frac{1}{6}e^{6\sigma }G_{\mu \nu \rho 4}G^{\mu \nu \rho 4}\right]\hfill \\ \multicolumn{3}{c}{}\\ & & \frac{3}{4\kappa ^2}\left[(_\mu \gamma )(^\mu \gamma )+e^{2\gamma }(_\mu a)(^\mu a)\right].\hfill \end{array}$$
(2.29)
In this expression, $`\kappa `$ is the four-dimensional gravitational coupling
$$\kappa ^2=\frac{\kappa _{11}^2}{V_7},$$
$`V_7=V_1V_6`$ being the volume of the compact space $`S^1\times X_6`$.
At this stage, the identification of the bosonic components $`C`$ and $`b_{\mu \nu }`$ of $`V`$, $`C_T`$ and $`T_\mu `$ of $`V_T`$ with the bulk fields $`\sigma `$, $`\gamma `$, $`C_{\mu \nu 4}`$ and $`a`$ can only be determined up to a proportionality constant for each multiplet. We will define these constants later on from the couplings of $`C`$ and $`C_T=2ReT`$ to charged matter and gauge fields, to obtain:
$$\begin{array}{cc}4\kappa ^2C=\frac{\lambda ^2}{V_6}e^{3\sigma },\hfill & 4\kappa ^2b_{\mu \nu }=\frac{\lambda ^2}{V_6}C_{\mu \nu 4},\hfill \\ \multicolumn{2}{c}{}\\ C_T=2\frac{\lambda ^2}{V_6}e^\gamma =\mathrm{\hspace{0.17em}\hspace{0.17em}2}ReT,\hfill & T_\mu =2\frac{\lambda ^2}{V_6}_\mu a=2_\mu ImT.\hfill \end{array}$$
(2.30)
The quantity $`\lambda `$ is the gauge coupling constant on the $`𝐙_2`$ fixed planes and $`\lambda ^2/V_6`$ is a dimensionless number which will actually never appear in the four-dimensional effective theory.
#### 2.2.4 Bianchi identities and symmetries
In the bulk Lagrangian (2.7), the terms $`[(S+\overline{S})V+L_TV_T]_D`$ impose in particular the Bianchi identities (2.6). They are certainly invariant under
$$\begin{array}{cccc}\hfill V& & V+L,\hfill & L\mathrm{linear},\hfill \\ \multicolumn{4}{c}{}\\ \hfill V_T& & V_T+T+\overline{T},\hfill & T\mathrm{chiral}.\hfill \end{array}$$
These symmetries are the supersymmetric extensions of the gauge invariances of Bianchi identities, $`\delta G_{\mu \nu \rho 4}=3_{[\mu }\mathrm{\Lambda }_{\nu \rho ]}`$ and $`\delta G_{\mu i\overline{j}4}=i_\mu \mathrm{\Lambda }\delta _{i\overline{j}}`$. Solving for $`S`$ and $`L_T`$ implies then that $`V`$ and $`V_T`$ are ‘pure gauge’, $`V=L`$ and $`V_T=T+\overline{T}`$. The last equation defines $`V_T`$ up to a holomorphic redefinition of $`T`$, $`Tf(T)`$. This redefinition is a symmetry of the bulk Lagrangian if the function $`\mathrm{\Phi }`$ simultaneously transforms as in (2.3), with $`\mathrm{\Lambda }`$ a holomorphic function of $`T`$. The equation for the invariance of $`S_0\overline{S}_0V_T`$ is
$$f(T)+\overline{f}(\overline{T})=\frac{T+\overline{T}}{\mathrm{\Lambda }(T)\overline{\mathrm{\Lambda }}(\overline{T})},$$
and its solution is clearly $`Sl(2,𝐑)`$ symmetry,
$$Tf(T)=\frac{aTib}{icT+d},adbc=1,$$
(2.31)
the modular invariance of $`T`$ (T-duality), extended to a continuous symmetry at the lowest order.
This chiral symmetry is generically anomalous: in the presence of a $`N=1`$ super-Yang-Mills sector, with or without chiral matter, mixed anomalies arise in the triangle diagram for two gauge bosons and one connection $`3\mathrm{log}(T+\overline{T})`$ for $`Sl(2,𝐑)`$ symmetry. This anomaly is cancelled in particular by a Green-Schwarz mechanism as was demonstrated in the effective Lagrangian description of gauge thresholds calculated at one-loop for $`(2,2)`$ compactifications of the heterotic strings . We will see below that this phenomenon is also a useful tool in the construction of effective supergravities of M-theory compactifications.
#### 2.2.5 Superpotential
The standard reduction of eleven-dimensional supergravity with unbroken $`N=1`$ supersymmetry does not generate a superpotential. This fact is however not a direct consequence of the eleven-dimensional Bianchi identity or of the Calabi-Yau and $`S^1/𝐙_2`$ symmetries. In principle, the Bianchi identity $`_{[M}G_{NPQR]}=0`$ allows a solution
$$G_{ijk4}=2i\kappa ^1hϵ_{ijk},G_{\overline{ijk}4}=2i\kappa ^1hϵ_{\overline{ijk}}.$$
(2.32)
In these equations, $`h`$ is a constant chosen real and $`ϵ_{ijk}`$ is the $`SU(3)`$–invariant Calabi-Yau tensor. The Lagrangian term $`\frac{e}{2\kappa _{11}^2}\frac{1}{48}G_{MNPQ}G^{MNPQ}`$ leads then to a contribution
$$\frac{e}{\kappa ^4}C_T^3(2\kappa ^2C)h^2$$
in the four-dimensional effective supergravity. This contribution corresponds to the addition of a superpotential term
$$[ihS_0^3]_F$$
to the bulk Lagrangian, a contribution which however breaks supersymmetry . Since we have insisted in writing Lagrangians in which all Bianchi identities are field equations, we prefer instead to use
$$[U(W+\overline{W})]_D+[S_0^3W]_F.$$
(2.33)
The field equation of the vector multiplet $`U`$ (weights $`w=2`$, $`n=0`$) implies that the chiral multiplet $`W`$ ($`w=n=0`$) is an arbitrary imaginary constant, which can be zero and supersymmetry stays unbroken, or non-zero.
With the addition (2.33) of a superpotential, the bulk Lagrangian takes its final ‘off-shell’ form
$$_\mathrm{B}=\left[\mathrm{{\rm Y}}(S+\overline{S})V+L_TV_T+U(W+\overline{W})\right]_D+[S_0^3W]_F,$$
(2.34)
in which the Bianchi identities of eleven-dimensional supergravity are translated into field equations of the Lagrange multipliers $`S`$, $`L_T`$ and $`U`$. At this stage, the introduction of these multiplets is not fascinating. This approach encodes simply the (Poincaré) dualities relating antisymmetric tensors (in linear multiplets) and scalars (in chiral multiplets), and the Bianchi identity for the superpotential is trivial. But this procedure will prove useful and informative in the forthcoming sections.
## 3 Fixed planes: gauge and matter contributions
In this section, we start with the well-known effective $`N=1`$ four-dimensional supergravity for symmetric $`(2,2)`$ compactifications of heterotic strings. We then rewrite this theory in a form where explicit Bianchi identities allow a direct comparison with $`O_7`$ compactification of M-theory, in the so-called upstairs formulation .
The dependence on charged matter (in chiral multiplets collectively denoted by $`M`$, with $`w=n=0`$) and gauge multiplets (vector multiplet $`A`$, in the adjoint representation, with $`w=n=0`$) of the effective supergravity theory for Calabi-Yau compactifications of heterotic strings is well-known , at least for the ‘universal’ matter multiplets arising from the simplest Calabi-Yau modes of the ten-dimensional super-Yang-Mills fields. Information on the non-trivial harmonic modes is more subtle , as for generic Calabi-Yau moduli. In the chiral formulation, Eq. (2.12) becomes
$$\begin{array}{ccc}\hfill _\mathrm{c}& =& \frac{3}{2}\left[S_0\overline{S}_0e^{K/3}\right]_D+[\frac{1}{4}S𝒲𝒲+S_0^3W]_F,\hfill \\ \multicolumn{3}{c}{}\\ \hfill K& =& \mathrm{log}(S+\overline{S})3\mathrm{log}(T+\overline{T}2\overline{M}e^AM),\hfill \\ \multicolumn{3}{c}{}\\ \hfill W& =& \alpha M^3.\hfill \end{array}$$
(3.1)
For notational simplicity, we omit traces over the gauge group representation and their normalisation factors. The chiral multiplet $`𝒲`$ is the gauge field-strength for $`A`$ ($`w=n=3/2`$). Since the gauge group is in general not simple,
$$𝒲𝒲=\underset{a}{}c^a𝒲^a𝒲^a,$$
(3.2)
with a (real) coefficient $`c^a`$ for each simple or abelian factor<sup>12</sup><sup>12</sup>12Corresponding to Kac-Moody levels in superstrings. All coefficients can be equal to one, as with the ‘standard embedding’, but our discussion is not affected by their presence.. The superpotential should be understood as a gauge invariant trilinear interaction with coupling constant $`\alpha `$ defined as an integral over the Calabi-Yau space. In the linear multiplet version, the equivalent expression is
$$_\mathrm{l}=\frac{1}{\sqrt{2}}\left[(S_0\overline{S}_0)^{3/2}\widehat{L}^{1/2}e^{\widehat{K}/2}\right]_D+[\alpha S_0^3M^3]_F.$$
(3.3)
With respect to Eq. (2.8), gauge and matter dependence arises in modifications of the linear multiplet $`L`$ (to $`\widehat{L}`$) and of $`\widehat{K}`$: the new modulus and matter Kähler potential is
$$\widehat{K}=3\mathrm{log}(T+\overline{T}2\overline{M}e^AM),$$
(3.4)
instead of Eq. (2.10) and
$$\widehat{L}=L2\mathrm{\Omega },$$
(3.5)
where $`\mathrm{\Omega }(A)`$ is the Chern-Simons vector multiplet ($`w=2`$, $`n=0`$), defined by <sup>13</sup><sup>13</sup>13 In global Poincaré supersymmetry, $`\mathrm{\Sigma }(\mathrm{\Omega })=\frac{1}{4}\overline{DD}\mathrm{\Omega }`$. A linear multiplet is defined by the condition $`\mathrm{\Sigma }(L)=0`$.
$$\mathrm{\Omega }=\underset{a}{}c^a\mathrm{\Omega }^a,\mathrm{\Sigma }(\mathrm{\Omega }^a)=\frac{1}{16}𝒲^a𝒲^a.$$
(3.6)
Insisting as before on Bianchi identities, both forms (3.1) and (3.3) are equivalent to
$$\begin{array}{ccc}\hfill & =& [(S_0\overline{S}_0V_T)^{3/2}(2V)^{1/2}(S+\overline{S})(V+2\mathrm{\Omega })\hfill \\ \multicolumn{3}{c}{}\\ & & +L_T(V_T+2\overline{M}e^AM)+\{U(W\alpha M^3)+\mathrm{c}.\mathrm{c}.\}]_D+[S_0^3W]_F\hfill \\ \multicolumn{3}{c}{}\\ & =& \left[\mathrm{{\rm Y}}(S+\overline{S})(V+2\mathrm{\Omega })+L_T(V_T+2\overline{M}e^AM)\right]_D+[S_0^3(ih+\alpha M^3)]_F.\hfill \end{array}$$
(3.7)
Supersymmetric vacua have $`h=0`$. As before, solving in the last expression for $`S`$ and $`L_T`$ imposes respectively $`V=L2\mathrm{\Omega }=\widehat{L}`$ and $`V_T=T+\overline{T}2\overline{M}e^AM`$, leading to Eq. (3.3). Alternatively, solving for $`V`$ and $`L_T`$ leads back to the chiral form (3.1), with the tensor calculus identity
$$2[(S+\overline{S})\mathrm{\Omega }]_D=\frac{1}{4}\underset{a}{}c^a[S𝒲^a𝒲^a]_F+\mathrm{derivative},$$
(3.8)
which follows from Eq. (3.6) and the definition of the $`F`$-density, $`[\mathrm{\Sigma }(\mathrm{})]_F=[\mathrm{}]_D`$.
This reformulation of the gauge invariant Lagrangian suggests some remarks. Firstly, it enhances the importance of Chern-Simons multiplets in superstring effective actions: gauge fields and matter fields couple to the bulk Lagrangian using a Chern-Simons multiplet. The gauge Chern-Simons multiplet $`\mathrm{\Omega }`$ is defined by Eq. (3.6), which indicates that its chiral projection $`\mathrm{\Sigma }(\mathrm{\Omega })`$ is the chiral multiplet for the kinetic super-Yang-Mills Lagrangian. Similarly for chiral matter, the kinetic Wess-Zumino Lagrangian can be written as $`[S_0\overline{S}_0\overline{M}e^AM]_D=[\mathrm{\Sigma }(S_0\overline{S}_0\overline{M}e^AM)]_F`$, defining
$$\mathrm{\Omega }_M=S_0\overline{S}_0\overline{M}e^AM$$
(3.9)
as a matter Chern-Simons multiplet ($`w=2`$, $`n=0`$) which then couples to $`S_0\overline{S}_0V_T`$ as $`\mathrm{\Omega }`$ couples to $`V`$.
Secondly, the Chern-Simons vector multiplet $`\mathrm{\Omega }(A)`$ is not gauge invariant: its variation is a linear multiplet. Then, the variation of $`[(S+\overline{S})\mathrm{\Omega }]_D`$ is a derivative and $`V`$ remains gauge invariant. When solving for $`S`$, it simply follows that $`\widehat{L}`$ is gauge invariant and that the linear multiplet transforms as
$$\delta L=2\delta \mathrm{\Omega }.$$
(3.10)
Finally, expression (3.7) shows that all gauge and chiral matter contributions can be viewed as the supersymmetrization of modified Bianchi identities imposed by $`S`$, $`L_T`$ and $`U`$. This is equally true in the ten-dimensional supergravity–Yang-Mills system: the curl of the antisymmetric tensor field is modified by Chern-Simons contributions which are supersymmetry partners of the super-Yang-Mills Lagrangian . This observation provides the link to the approach based on M-theory on $`O_7`$, in which the $`𝐙_2`$–fixed planes carrying the Yang-Mills fields induce because of supersymmetry modifications to the Bianchi identity of the four-form field strength of eleven-dimensional supergravity.
In the effective supergravity of M-theory on $`O_7`$ (‘upstairs formulation’), the various components of the Lagrangian (3.7),
$$\begin{array}{ccc}\hfill & =& [(S_0\overline{S}_0V_T)^{3/2}(2V)^{1/2}(S+\overline{S})(V+2\mathrm{\Omega })\hfill \\ \multicolumn{3}{c}{}\\ & & +L_T(V_T+2\overline{M}e^AM)+\{U(W\alpha M^3)+\mathrm{c}.\mathrm{c}.\}]_D+[S_0^3W]_F,\hfill \end{array}$$
have the following origin. As already discussed at length, the first term is the bulk supergravity contribution. Then $`[(S+\overline{S})(V+2\mathrm{\Omega })]_D`$ is the supersymmetrization of the Bianchi identity verified by the component $`G_{\mu \nu \rho 4}`$ of the four-form field, modified by gauge contributions on the fixed planes. Similarly, $`[L_T(V_T+2\overline{M}e^AM)]_D`$ and $`[U(W\alpha M^3)+\mathrm{c}.\mathrm{c}.]_D`$ are respectively the supersymmetric extensions of the Bianchi identities of $`G_{\mu j\overline{k}4}`$ and $`G_{ijk4}`$, when fixed plane contributions are included. Thus, all fixed plane contributions are given at this order by the supersymmetrization of Bianchi identities, as obtained by direct $`O_7`$ truncation of the eleven-dimensional identities .
At this point, the gauge coupling constant for each simple or abelian factor $`a`$ in the gauge group appears to be
$$\frac{1}{g_a^2}=c^aRes=\frac{c^a\mathrm{{\rm Y}}}{4C},$$
(3.11)
$`s`$ and $`C`$ being respectively the lowest scalar component of the chiral $`S`$ and the vector $`V`$ (or the linear $`L`$). At this order, $`g_a`$ is the tree-level wilsonnian and physical<sup>14</sup><sup>14</sup>14The coefficient of $`\frac{1}{4}F_{\mu \nu }^aF^{a\mu \nu }`$ in the generating functional of one-particle irreducible Green’s functions. gauge coupling. The second equality is the lowest component of the equation of motion of the vector multiplet $`V`$, Eq. (2.18). In the Einstein frame, $`Res=(4\kappa ^2C)^1`$.
It is clear, as already observed , that as far as the structure of the four-dimensional effective supergravity is concerned, the same information follows from $`O_7`$ compactification of M-theory at the next to lowest order in the $`\kappa `$ expansion and from Calabi-Yau compactifications of the heterotic strings, at zero string loop order.
Notice that Eq. (3.11) defines $`c^aRes`$ as the coefficient of gauge kinetic terms. It defines then this field in terms of the gauge kinetic action on the ten-dimensional $`𝐙_2`$ fixed planes,
$$𝒮_{\mathrm{gauge}}=\frac{1}{4\lambda ^2}_{M_{10}}d^{10}xe_{10}\mathrm{tr}F_{AB}F^{AB},$$
(3.12)
reduced on $`X_6`$. This action is also at the origin of charged matter kinetic terms, which in the effective supergravity read
$$\frac{e}{\kappa ^2}\frac{^2\widehat{K}}{\overline{M}M}(D_\mu \overline{M})(D^\mu M)=\frac{6e}{\kappa ^2C_T}(D_\mu \overline{M})(D^\mu M)+\mathrm{}$$
The Calabi-Yau reduction of the action (3.12) provides then the identification of $`C`$ and $`C_T`$ in terms of the bulk fields $`\sigma `$ and $`\gamma `$ appearing in the metric tensor (2.28). These results have already been displayed in Eq. (2.30).
## 4 Anomaly-cancelling terms
In the ten-dimensional heterotic string, cancellation of gauge and gravitational anomalies is a one-loop effect in string or effective supergravity perturbation theory. In the low-energy effective action description, we should then distinguish the Wilson effective supergravity from the standard effective action $`𝒮_\mathrm{\Gamma }`$, defined as the generating functional of one-particle irreducible Green’s functions. The latter action can be obtained in a diagrammatic expansion built from the Wilson Lagrangian $``$, itself obtained from string perturbation theory as an expansion
$$=^{(0)}+^{(1)}+\mathrm{},$$
the subscript being the string-loop order. The expressions given in the previous sections were for $`^{(0)}`$, or for the tree-level $`𝒮_\mathrm{\Gamma }`$. At the string one-loop level, $`𝒮_\mathrm{\Gamma }`$ includes tree and one-loop diagrams generated by the Feynman rules of the tree-level Wilson Lagrangian $`^{(0)}`$. These include anomalous loop diagrams. It also includes tree diagrams generated by the Green-Schwarz counterterm introduced in $`^{(1)}`$ to cancel the anomalies generated by $`^{(0)}`$. The mechanism for symmetry restoration implies that $`^{(1)}`$ is not invariant under the restored symmetry.
In four space-time dimensions, the nature of the cancelled anomalies is known from studies of $`(2,2)`$ compactifications of heterotic strings in the Yang-Mills sector : target-space duality of the modulus $`T`$ has a one-loop anomaly which is cancelled by a counterterm in $`^{(1)}`$, in a generalization to sigma-model anomalies of the Green-Schwarz mechanism . The derivation of the complete counterterm requires a calculation to all orders in the modulus $`T`$ . However, at the present stage of understanding, the M-theory approach should be regarded as a large-$`T`$ limit in which T-duality reduces to a shift symmetry in the imaginary part of $`T`$.
Our goal in this section is to obtain some or all counterterms in $`^{(1)}`$ associated with anomaly-cancellation in the low-energy description. We are particularly interested in contributions to gauge kinetic terms, the so-called threshold corrections. And we want to formulate these terms using the ‘M-theory multiplets’ $`V`$, $`V_T`$ and $`W`$ corresponding to the surviving components of $`G`$, in contrast to the ‘heterotic multiplets’ $`S`$ (or $`L`$) and $`T`$. We begin by obtaining the relevant information from the case of heterotic $`(2,2)`$ symmetric orbifolds.
### 4.1 Information from symmetric (2,2) orbifolds
Retaining only the universal modulus $`T`$ and a linear dilaton multiplet $`L`$, the Wilson one-loop Lagrangian for heterotic symmetric $`(2,2)`$ orbifolds includes a term<sup>15</sup><sup>15</sup>15In this paragraph, we use $`c^a=1`$ and $`\widehat{L}=L2_a\mathrm{\Omega }^a`$.
$$^{(1)}=2\delta _{GS}\left[\widehat{L}\mathrm{log}(T+\overline{T})\right]_D+\underset{a}{}b^a[\mathrm{log}\eta (iT)𝒲^a𝒲^a]_F,$$
(4.1)
where $`\delta _{GS}`$ and $`b^a`$ are numbers depending on the orbifold and $`\eta (iT)`$ is the Dedekind function. Under $`Sl(2,𝐙)`$ T-duality (2.31), the variation of $`^{(1)}`$ is
$$\begin{array}{ccc}\hfill \delta ^{(1)}& =& 2\delta _{GS}\left[\widehat{L}\{\mathrm{log}\phi (T)+\mathrm{log}\overline{\phi }(\overline{T})\}\right]_D+\frac{1}{2}b^a\left[\mathrm{log}\phi (T)𝒲^a𝒲^a\right]_F\hfill \\ \multicolumn{3}{c}{}\\ & =& \frac{1}{2}(\delta _{GS}+b^a)\left[\mathrm{log}\phi (T)𝒲^a𝒲^a\right]_F,\hfill \end{array}$$
with $`\phi (T)=icT+d`$. On the other hand, the triangle one-loop diagram for two gauge fields and one Kähler connection $`3\mathrm{log}(T+\overline{T})`$ is anomalous. Its variation is
$$\delta \mathrm{\Delta }=\frac{1}{2}A^a\left[\mathrm{log}\phi (T)𝒲^a𝒲^a\right]_F,$$
where $`A^a`$ is the chiral-anomaly coefficient, as obtained from the expression of the diagram. The anomaly cancels since one finds that $`b^a+\delta _{GS}+A^a=0`$ for all factors in the gauge group (the index $`a`$). The one-loop correction to gauge kinetic terms obtained from the component expansion of $`^{(1)}`$ and from the triangle diagram reads
$$\frac{1}{4}F_{\mu \nu }^aF^{a\mu \nu }\left[(\delta _{GS}+A^a)\mathrm{log}(T+\overline{T})+b^a\mathrm{log}|\eta (iT)|^4\right].$$
(4.2)
It is modular invariant since the anomaly is cancelled, and its value is controlled by the $`F`$-density contribution to $`^{(1)}`$, with coefficients $`b^a`$.
The Wilson Lagrangian depends on the coefficients $`\delta _{GS}`$ and $`b^a`$. But the information on gauge thresholds is in the numbers $`b^a`$. In general, the parameters of the Wilson Lagrangian computed at a non-trivial loop order are not of direct physical significance and this is here the case of $`\delta _{GS}`$ in the sector of gauge kinetic terms.
In the large-$`T`$ limit, T-duality reduces to $`ImTImT+\text{constant}`$, the Kähler connection $`3\mathrm{log}(T+\overline{T})`$ is invariant and, strictly speaking, no anomaly survives to be cancelled. In addition, $`\mathrm{log}|\eta (iT)|^4\frac{\pi }{3}(T+\overline{T})`$ dominates the logarithmic contributions. The threshold correction is then of the simple form
$$\frac{1}{4}\underset{a}{}\left[\frac{\pi b^a}{3}(T+\overline{T})\right]F_{\mu \nu }^aF^{a\mu \nu },$$
invariant under the imaginary shift symmetry of $`T`$. Its supersymmetrization is
$$\frac{\pi }{12}\underset{a}{}b^a[T𝒲^a𝒲^a]_F=\frac{2\pi }{3}\underset{a}{}b^a[(T+\overline{T})\mathrm{\Omega }^a]_D.$$
As long as $`T`$ only is considered, there seems to be no way to identify the $`D`$-density contribution to $`^{(1)}`$ in the large-$`T`$ limit. The introduction of the matter multiplet $`M`$ changes the picture. In the lowest order Wilson Lagrangian $`^{(0)}`$, the Kähler connection $`3\mathrm{log}(T+\overline{T})`$ is modified to $`3\mathrm{log}(T+\overline{T}2\overline{M}e^AM)`$. As a consequence, the contribution with coefficient $`A^a`$ to the gauge threshold (4.2) involves the quantity $`\mathrm{log}(T+\overline{T}2\overline{M}M)`$. Then, either the one-loop term $`^{(1)}`$ is accordingly modified to $`2\delta _{GS}[\widehat{L}\mathrm{log}(T+\overline{T}2\overline{M}e^AM)]_D`$ and the parameter $`\delta _{GS}`$ disappears from gauge thresholds, or, less plausibly, this is not the case and $`\delta _{GS}`$ acquires a physical significance in $`M`$-dependent thresholds. In any case, since the holomorphic Dedekind function cannot depend on $`\overline{M}e^AM`$, a calculation of the correction $`^{(1)}`$ to the Wilson Lagrangian to first order in $`\overline{M}e^AM`$ will give access to the $`D`$-density Green-Schwarz term and to the parameter $`\delta _{GS}`$.
From the point of view of M-theory on $`O_7`$ or heterotic strings on Calabi-Yau, $`\overline{M}e^AM`$ arises from the ten-dimensional Chern-Simons terms, at the same order as gauge kinetic terms. This indicates that the $`D`$-density contribution to $`^{(1)}`$ should be visible in a reduction of the ten-dimensional Green-Schwarz terms.
### 4.2 The case of M-theory on $`O_7`$
Before embarking in a derivation of the anomaly-cancelling terms from the low-energy limit of M-theory on $`O_7`$, we consider the problem at the level of four-dimensional supergravity only.
In the large-$`T`$ limit, as discussed in the previous paragraph, the $`T`$-dependent corrections to gauge kinetic terms are of the form
$$\frac{1}{4}\underset{a}{}\beta ^a\left[T𝒲^a𝒲^a\right]_F,$$
(4.3)
with coefficients which are in principle calculable in heterotic strings. To rewrite them in terms of the ‘M-theory multiplets’, we first note that
$$\frac{1}{4}\underset{a}{}\beta ^a\left[T𝒲^a𝒲^a\right]_F=2\left[(T+\overline{T})\underset{a}{}\beta ^a\mathrm{\Omega }^a\right]_D.$$
(4.4)
Since the field equation of the Lagrange multiplier $`L_T`$ implies that $`V_T=T+\overline{T}2\overline{M}e^AM`$, this expression can also be written
$$2\left[(V_T+2\overline{M}e^AM)\underset{a}{}\beta ^a\mathrm{\Omega }^a\right]_D.$$
(4.5)
The right-hand side of Eq. (4.4) is gauge invariant because $`\delta \mathrm{\Omega }^a`$ is a linear multiplet and therefore $`\left[(T+\overline{T})\delta \mathrm{\Omega }^a\right]_D`$ is a derivative. To ensure gauge invariance of expression (4.5), we add the term $`\left[L_T(V_T+2\overline{M}e^AM)\right]_D`$ included in Lagrangian (3.7). We obtain
$$\left[(L_T2\underset{a}{}\beta ^a\mathrm{\Omega }^a)(V_T+2\overline{M}e^AM)\right]_D,$$
(4.6)
which is gauge invariant if we postulate that
$$\delta L_T=2\underset{a}{}\beta ^a\delta \mathrm{\Omega }^a.$$
(4.7)
The correction (4.4) to the super-Yang-Mills Lagrangian is independent of the matter fields and has a holomorphic character (it can be seen as a correction to the holomorphic gauge kinetic function $`f_{ab}`$). To enumerate possible matter-dependent contributions to gauge kinetic terms, we consider for simplicity a single matter multiplet $`M`$ transforming in some unspecified representation of the gauge group. The first candidate counterterm is a real density:
$$\underset{a}{}\gamma ^a\left[\overline{M}e^AM\mathrm{\Omega }^a\right]_D.$$
Gauge invariance requires however its appearance in the combination
$$2\delta \left[\overline{M}e^AM(L2\underset{a}{}c^a\mathrm{\Omega }^a)\right]_D+2\gamma \left[\overline{M}e^AM(L_T2\underset{a}{}\beta ^a\mathrm{\Omega }^a)\right]_D,$$
or
$$2\delta \left[\overline{M}e^AMV\right]_D+2\gamma \left[\overline{M}e^AM(L_T2\underset{a}{}\beta ^a\mathrm{\Omega }^a)\right]_D,$$
using ‘M-theory multiplet’ $`V`$. In the first counterterm, each gauge group factor contributes with weight $`c^a`$, as in lowest order terms. The second contribution can be combined with expression (4.6) into
$$\left[(L_T2\underset{a}{}\beta ^a\mathrm{\Omega }^a)(V_T+2(1+\gamma )\overline{M}e^AM)\right]_D.$$
(4.8)
We will also see below that the M-theory anomaly-cancelling terms generate a contribution of the form
$$ϵ\left[V|\alpha M^3|^2\right]_D,$$
(4.9)
involving the matter superpotential. Since the factor $`1+\gamma `$ in expression (4.8) can be eliminated by a rescaling of $`M`$, of $`\delta `$ and of the superpotential coupling constant $`\alpha `$, we may take $`\gamma =0`$ at our level of approximation, and the Wilson Lagrangian up to string one-loop order is expected to become
$$\begin{array}{ccc}\hfill & =& [\mathrm{{\rm Y}}(S+\overline{S})(V+2\mathrm{\Omega })+\{U(W\alpha M^3)+\mathrm{c}.\mathrm{c}.\}\hfill \\ \multicolumn{3}{c}{}\\ & & +(L_T2_a\beta ^a\mathrm{\Omega }^a)(V_T+2\overline{M}e^AM)\hfill \\ \multicolumn{3}{c}{}\\ & & +V(ϵ|\alpha M^3|^22\delta \overline{M}e^AM)]_D+[S_0^3W]_F.\hfill \end{array}$$
(4.10)
Notice that the contributions which correspond to string one-loop effects do not include any correction to the Einstein Lagrangian, which remains simply
$$\frac{1}{2}\mathrm{{\rm Y}}eR,\mathrm{{\rm Y}}=(z_0\overline{z}_0C_T)^{\frac{3}{2}}(2C)^{\frac{1}{2}},$$
and the Einstein frame condition remains $`\mathrm{{\rm Y}}=\kappa ^2`$. This is expected since the gravitational constant in the heterotic string is not corrected at string one-loop order.
From the general expression (4.10) in which Bianchi identities are field equations for $`S`$, $`L_T`$ and $`U`$, we can derive various equivalent forms. For instance, solving for $`S`$, $`L_T`$ and $`U`$ leads to the version of the effective supergravity in which the dilaton is described by a linear multiplet:
$$\begin{array}{ccc}\hfill _\mathrm{l}& =& \left[(S_0\overline{S}_0)^{3/2}\left(T+\overline{T}2\overline{M}e^AM\right)^{3/2}(2\widehat{L})^{1/2}\right]_D+\left[S_0^3(ih+\alpha M^3)\right]_F\hfill \\ \multicolumn{3}{c}{}\\ & & +\left[\widehat{L}(ϵ|\alpha M^3|^22\delta \overline{M}e^AM)\right]_D+[\frac{1}{4}T_a\beta ^a𝒲^a𝒲^a]_F.\hfill \end{array}$$
(4.11)
The second line is a one-loop correction in the perturbative expansion of the heterotic string in which $`\widehat{L}`$ is the string loop-counting field . Its $`T`$-dependent part corresponds to the Green-Schwarz counterterm found in for symmetric heterotic orbifolds. Each of these one-loop corrections, with coefficients $`ϵ`$, $`\delta `$ and $`\beta ^a`$, is related to a well-defined counterterm which can be easily identified in, for instance, the low-energy limit of M-theory on $`O_7`$. The Green-Schwarz counterterms controlled by $`\delta `$ and $`ϵ`$ are intrinsically real $`D`$-densities. They will appear as corrections to the Kähler potential, as matter-dependent ‘wave-function renormalisations’, in the dual version with a chiral dilaton multiplet. On the other hand, the holomorphic $`T`$-dependent terms are true threshold corrections.
Alternatively, solving for $`L_T`$, $`V`$ and $`U`$ leads to the version with a chiral dilaton multiplet:
$$_\mathrm{c}=\frac{3}{2}\left[S_0\overline{S}_0e^{K/3}\right]_D+\frac{1}{4}\left[\underset{a}{}(c^aS+\beta ^aT)𝒲^a𝒲^a\right]_F+\left[S_0^3(ih+\alpha M^3)\right]_F,$$
(4.12)
with the Kähler potential<sup>16</sup><sup>16</sup>16The superfield Kähler potential includes covariantization contributions $`e^A`$ which disappear in the bosonic expression used in component expansions.
$$\begin{array}{ccc}\hfill K& =& \mathrm{log}\left(S+\overline{S}+2\delta \overline{M}e^AMϵ|\alpha M^3|^2\right)\hfill \\ \multicolumn{3}{c}{}\\ & & 3\mathrm{log}\left(T+\overline{T}2\overline{M}e^AM\right),\hfill \end{array}$$
(4.13)
and the gauge kinetic functions $`f^a=c^aS+\beta ^aT`$. An ambiguity exists however because one can perform a holomorphic redefinition of the two chiral multiplets. For instance,
$$S=\stackrel{~}{S}\delta T,$$
(4.14)
leads to the equivalent Kähler potential
$$\begin{array}{ccc}\hfill K& =& \mathrm{log}\left(\stackrel{~}{S}+\overline{\stackrel{~}{S}}\delta (T+\overline{T}2\overline{M}e^AM)ϵ|\alpha M^3|^2\right)\hfill \\ \multicolumn{3}{c}{}\\ & & 3\mathrm{log}\left(T+\overline{T}2\overline{M}e^AM\right),\hfill \end{array}$$
(4.15)
with gauge kinetic functions $`f^a=c^a\stackrel{~}{S}+(\beta ^ac^a\delta )T`$. The origin of this ambiguity at the level of ‘M-theory multiplets’ is interesting. Suppose that we add the counterterm
$$\mathrm{\Delta }=A\left[(V+2\mathrm{\Omega })(V_T+2\overline{M}e^AM)\right]_D$$
to the fundamental Lagrangian (4.10), with an arbitrary constant $`A`$. The theory becomes then
$$\begin{array}{ccc}\hfill _A& =& [\mathrm{{\rm Y}}(S+\overline{S})(V+2\mathrm{\Omega })+\{U(W\alpha M^3)+\mathrm{c}.\mathrm{c}.\}\hfill \\ \multicolumn{3}{c}{}\\ & & +\left(L_T2_a(\beta ^aAc^a)\mathrm{\Omega }^a\right)(V_T+2\overline{M}e^AM)\hfill \\ \multicolumn{3}{c}{}\\ & & +V(ϵ|\alpha M^3|^2+AV_T+2(A\delta )\overline{M}e^AM)]_D+[S_0^3W]_F.\hfill \end{array}$$
(4.16)
It is gauge invariant provided the appropriate transformation of $`L_T`$ is postulated. We have apparently obtained a family of four-dimensional supergravities, depending on a new parameter $`A`$. This is however only true before solving for the Lagrange multiplier multiplets. Firstly, solving for $`S`$ and $`L_T`$ leads to the space-time derivative $`\mathrm{\Delta }=A\left[L(T+\overline{T})\right]_D`$. The counterterm $`\mathrm{\Delta }`$ is then irrelevant in the version of the theory with a linear dilaton. Secondly, if we instead solve for $`V`$ and $`L_T`$, we obtain the Kähler potential
$$\begin{array}{ccc}\hfill K& =& \mathrm{log}\left(S+\overline{S}A(T+\overline{T})+2\delta \overline{M}e^AMϵ|\alpha M^3|^2\right)\hfill \\ \multicolumn{3}{c}{}\\ & & 3\mathrm{log}\left(T+\overline{T}2\overline{M}e^AM\right),\hfill \end{array}$$
and the gauge kinetic functions $`f^a=c^aS+(\beta ^aAc^a)T`$. This theory is clearly related to Eq. (4.13) by the holomorphic redefinition $`SSAT`$, and the choice $`A=\delta `$ leads to theory (4.15).
This discussion shows that the counterterm $`\mathrm{\Delta }`$ is irrelevant in the four-dimensional effective supergravity, that all values of $`A`$ lead to equivalent Lagrangians, with the same dynamical equations. Further information due, for instance, to compactification of extra dimensions could however appear more natural with a specific value of $`A`$, if one insists to use the version of the effective supergravity with a chiral dilaton multiplet. For instance, all corrections linear in $`T`$ appear as gauge thresholds with the choice $`A=0`$. But one could as well use the version with a linear dilaton which is free of ambiguities.
In addition, the holomorphic redefinition (4.14) mixes terms of different orders in the string loop expansion. As a consequence, the distinction between terms in $`^{(0)}`$ and corrections in $`^{(1)}`$ becomes ambiguous in general in the large $`T`$ limit.
The values of the coefficients $`\delta `$ and $`ϵ`$ can be inferred from a direct calculation of the $`M`$-dependent anomaly-cancelling terms in M-theory on $`O_7`$. This is the subject of the next paragraph. Notice however that such a calculation only provides the terms of first order in the matter multiplets $`\overline{M}e^AM`$ and $`|\alpha M^3|^2`$. To this order, expression (4.13) becomes
$$\begin{array}{ccc}\hfill K& =& \mathrm{log}(S+\overline{S})3\mathrm{log}(T+\overline{T}2\overline{M}e^AM)\hfill \\ \multicolumn{3}{c}{}\\ & & \frac{1}{S+\overline{S}}\left[2\delta \overline{M}e^AMϵ|\alpha M^3|^2\right].\hfill \end{array}$$
(4.17)
The term with coefficient $`\delta `$ has been obtained in direct Calabi-Yau reductions of M-theory on $`S^1/𝐙_2`$ (see for instance <sup>17</sup><sup>17</sup>17It has also been obtained, in a quite different context, by Itoyama and Leon .). The charged matter contribution with coefficient $`ϵ`$ was not included in these analyses.
The gauge contributions appearing in Eq. (4.10) read
$$2\underset{a}{}\left[\{c^a(S+\overline{S})+\beta ^a(V_T+2\overline{M}e^AM)\}\mathrm{\Omega }^a\right]_D,$$
so that the gauge coupling constants are given by
$$\frac{1}{g_a^2}=c^aRes+\frac{1}{2}\beta ^a\left(C_T+2\overline{M}M\right).$$
(4.18)
This expression becomes harmonic once the Bianchi identity imposing $`C_T+2\overline{M}M=2ReT`$ has been used. Similarly, one obtains from Eq. (4.11)
$$\frac{1}{g_a^2}=\frac{c^a\mathrm{{\rm Y}}}{4C}+\beta ^aReTc^a\delta \overline{M}M+\frac{c^aϵ}{2}|\alpha M^3|^2,$$
(4.19)
with $`\mathrm{{\rm Y}}=2(z_0\overline{z}_0)^{3/2}(ReT\overline{M}M)^{3/2}C^{1/2}`$. This second form of the gauge couplings is never harmonic since it is obtained from a theory with a linear dilaton multiplet. Both expressions do however agree since the chiral-linear duality relation between $`Res`$ and $`C`$ is
$$Res=\frac{\mathrm{{\rm Y}}}{4C}\delta \overline{M}M+\frac{ϵ}{2}|\alpha M^3|^2.$$
### 4.3 On the eleven-dimensional origin of the <br>anomaly-cancelling terms
In ten dimensions, anomaly-cancelling terms for the $`E_8\times E_8`$ heterotic string are well-known. There are two terms. The first couples a gauge and Lorentz invariant eight-form $`\widehat{X}_8`$ to the two-form field $`\widehat{B}`$. The second one is proportional to $`(\mathrm{\Omega }_{3,1}+\mathrm{\Omega }_{3,2}\mathrm{\Omega }_{3\mathrm{L}})\widehat{X}_7`$, with $`d\widehat{X}_7=\widehat{X}_8`$.
The anomaly-cancelling terms arising from M-theory on $`S^1/𝐙_2`$ have recently been precisely computed . They arise from the following action terms in eleven dimensions :
$$\frac{\lambda ^2}{(4\pi )^2\kappa _{11}^2}GX_7\frac{1}{12\kappa _{11}^2}CGG,$$
(4.20)
where
$$X_7=\frac{1}{12(4\pi )^3}\left(\frac{1}{2}\mathrm{\Omega }_{7\mathrm{L}}\frac{1}{8}(\mathrm{tr}R^2)\mathrm{\Omega }_{3\mathrm{L}}\right),$$
(4.21)
$`\kappa _{11}`$ is the eleven-dimensional gravitational constant and $`\lambda `$ is the gauge coupling constant on both ten-dimensional $`𝐙_2`$-fixed planes, as defined by the gauge action (3.12). The Chern-Simons forms are defined by
$$d\mathrm{\Omega }_{7\mathrm{L}}=\mathrm{tr}R^4,d\mathrm{\Omega }_{3\mathrm{L}}=\mathrm{tr}R^2,d\mathrm{\Omega }_{3,i}=\mathrm{tr}F_i^2,i=1,2.$$
Cancellation of gauge and gravitational anomalies and coherence of the reduction to ten dimensions impose $`\lambda ^6=(4\pi )^5\kappa _{11}^4/12`$ . This condition relates the gauge coupling $`\lambda `$ and the $`S^1`$ radius. Solving the Bianchi identity verified by the four-form field $`G`$ on $`S^1/𝐙_2`$ and extracting the zero modes leads to the following Green-Schwarz terms
$$\begin{array}{ccc}\hfill S_{GS}& =& \frac{1}{48\pi }_{M_{10}}\widehat{B}\left[(I_{4,1})^2+(I_{4,2})^2I_{4,1}I_{4,2}+\frac{1}{(4\pi )^4}\left(\frac{1}{2}\mathrm{tr}R^4\frac{1}{8}(\mathrm{tr}R^2)^2\right)\right]\hfill \\ \multicolumn{3}{c}{}\\ & & \frac{1}{(4\pi )^2}_{M_{10}}(\mathrm{\Omega }_{3,1}+\mathrm{\Omega }_{3,2}\mathrm{\Omega }_{3\mathrm{L}})X_7.\hfill \end{array}$$
(4.22)
In this expression,
$$I_{4,i}=\frac{1}{(4\pi )^2}\left(\mathrm{tr}F_i^2\frac{1}{2}\mathrm{tr}R^2\right),i=1,2,$$
(4.23)
and<sup>18</sup><sup>18</sup>18Recall that $`x^4`$ is the $`S^1`$ coordinate.
$$\widehat{B}_{AB}=\frac{\lambda ^2}{\kappa _{11}^2}_{S^1}𝑑x^4C_{AB4}.$$
(4.24)
As usual, the two-form field $`\widehat{B}`$ couples to $`\widehat{X}_8`$. But the second term has a particular structure: $`\widehat{X}_7`$ is replaced by the purely gravitational seven-form (4.21): the anomaly-cancelling terms derived from M-theory differ from the standard expression of the heterotic string by a well-defined local counterterm, as permitted by the descent equations .
The Calabi-Yau compactification of $`\widehat{B}_{AB}`$ leads to two zero modes of the bulk fields,
$$\kappa ^2\widehat{B}_{\mu \nu }=\frac{\lambda ^2}{V_6}C_{\mu \nu 4},\kappa ^2\widehat{B}_{i\overline{j}}=\frac{\lambda ^2}{V_6}C_{i\overline{j}4},$$
since $`\frac{\lambda ^2}{\kappa _{11}^2}=\frac{1}{\kappa ^2}\frac{\lambda ^2}{V_6}\frac{1}{V_1}`$. By Eqs. (2.30), these states are related to our four-dimensional bulk multiplets $`L`$ (or $`S`$) and $`T`$ by
$$\widehat{B}_{\mu \nu }=4b_{\mu \nu },\kappa ^2\widehat{B}_{i\overline{j}}=iImT\delta _{i\overline{j}}.$$
(4.25)
The number $`\lambda ^2/V_6`$ disappears in this identification: it does not play any role in the four-dimensional effective theory.
Our task is then to compute the reduction of $`S_{GS}`$ on $`M_4\times `$(Calabi-Yau). Since we restrict ourselves to contributions with at most two derivatives, we only need the reduction of
$$\mathrm{\Delta }=\frac{1}{48\pi }_{M_{10}}\widehat{B}\left[(I_{4,1})^2+(I_{4,2})^2I_{4,1}I_{4,2}\right].$$
The global definition of the four-form field $`G`$ (‘cohomology condition’) implies
$$I_{4,1}=I_{4,2}$$
for the Calabi-Yau background, which is a $`(2,2)`$ form. The counterterm $`\mathrm{\Delta }`$ becomes then
$$\begin{array}{ccc}\hfill \mathrm{\Delta }& =& \frac{1}{4(4\pi )^3}_{M_{10}}\widehat{B}I_{4,1}\left(\mathrm{tr}F_1^2\mathrm{tr}F_2^2\right)\hfill \\ \multicolumn{3}{c}{}\\ & & \frac{1}{12(4\pi )^5}_{M_{10}}\widehat{B}\left[(\mathrm{tr}F_1^2)^2+(\mathrm{tr}F_2^2)^2(\mathrm{tr}F_1^2)(\mathrm{tr}F_2^2)\right]+\mathrm{},\hfill \end{array}$$
where gravitational contributions with more than two derivatives are omitted. Notice that the two $`E_8`$ factors contribute with opposite signs in the background-dependent term.
The standard embedding is defined by $`\mathrm{tr}F_2^2=0`$, while $`\mathrm{tr}F_1^2`$ is in the $`SU(3)`$ direction of the maximal embedding $`E_6\times SU(3)E_8`$. As a consequence,
$$\mathrm{tr}F_1^2=\mathrm{tr}R^2=2(4\pi )^2I_{4,1},$$
which in turn leads to
$$\begin{array}{ccc}\hfill \mathrm{\Delta }& =& \frac{1}{8(4\pi )^5}_{M_{10}}\widehat{B}\mathrm{tr}R^2(\mathrm{tr}F_1^2\mathrm{tr}F_2^2)\hfill \\ \multicolumn{3}{c}{}\\ & & \frac{1}{12(4\pi )^5}_{M_{10}}\widehat{B}\left[(\mathrm{tr}F_1^2)^2+(\mathrm{tr}F_2^2)^2(\mathrm{tr}F_1^2)(\mathrm{tr}F_2^2)\right]+\mathrm{}.\hfill \end{array}$$
(4.26)
To derive the zero modes of $`\mathrm{tr}F_1^2`$ and $`\mathrm{tr}F_2^2`$, it is simpler to consider the Chern-Simons forms, using the relation
$$(\mathrm{tr}F_i^2)_{ABCD}=4_{[A}(\mathrm{\Omega }_{3,i})_{BCD]},i=1,2.$$
In the standard embedding, the unbroken $`E_8`$ group generates a $`N=1`$ super-Yang-Mills multiplet only. The only massless mode is then
$$(\mathrm{tr}F_2^2)_{ABCD}(\mathrm{tr}F_{E_8}^2)_{\mu \nu \rho \sigma },(\mathrm{\Omega }_{3,2})_{ABC}(\mathrm{\Omega }_{E_8})_{\mu \nu \rho }.$$
Clearly, these massless modes of $`\mathrm{tr}F_2^2`$ and $`\mathrm{\Omega }_{3,2}`$ are respectively components of the four-dimensional $`N=1`$ multiplets $`𝒲^2𝒲^2`$ and $`\mathrm{\Omega }^2`$ used earlier. More precisely:
$$\begin{array}{ccc}\hfill [𝒲^2𝒲^2]_{\mathrm{f}\mathrm{component}}& =& \frac{1}{2}\mathrm{tr}F_{E_8\mu \nu }F_{E_8}^{\mu \nu }\frac{i}{4e}ϵ^{\mu \nu \rho \sigma }\mathrm{tr}F_{E_8\mu \nu }F_{E_8\rho \sigma }+\mathrm{},\hfill \\ \multicolumn{3}{c}{}\\ \hfill [\mathrm{\Omega }^2]_{\mathrm{d}\mathrm{component}}& =& \frac{1}{16}\mathrm{tr}F_{E_8\mu \nu }F_{E_8}^{\mu \nu }+\mathrm{},\hfill \\ \multicolumn{3}{c}{}\\ \hfill [\mathrm{\Omega }^2]_{\mathrm{vector}\mathrm{component}}& =& \frac{1}{8e}ϵ^{\mu \nu \rho \sigma }(\mathrm{\Omega }_{E_8})_{\nu \rho \sigma }+\mathrm{}\hfill \end{array}$$
(4.27)
The gauge fields of the $`E_8`$ group broken into $`E_6`$ generate $`E_6`$ gauge fields and the chiral matter multiplet $`M`$, transforming in representation $`\mathrm{𝟐𝟕}`$. Accordingly, the massless modes of $`\mathrm{tr}F_1^2`$ are:
$$\begin{array}{ccccc}\hfill (\mathrm{tr}F_1^2)_{ABCD}& & (\mathrm{tr}F_{E_6}^2)_{\mu \nu \rho \sigma },& & \\ \multicolumn{5}{c}{}\\ & & (\mathrm{tr}F_1^2)_{\mu \nu i\overline{j}}& =& 2_{[\mu }(\mathrm{\Omega }_{3,1})_{\nu ]i\overline{j}},\hfill \\ \multicolumn{5}{c}{}\\ & & (\mathrm{tr}F_1^2)_{\mu ijk}& =& _\mu (\mathrm{\Omega }_{3,1})_{ijk},\hfill \\ \multicolumn{5}{c}{}\\ & & (\mathrm{tr}F_1^2)_{\mu \overline{ijk}}& =& _\mu (\mathrm{\Omega }_{3,1})_{\overline{ijk}}.\hfill \end{array}$$
While as before $`(\mathrm{\Omega }_{3,1})_{ABC}(\mathrm{\Omega }_{E_6})_{\mu \nu \rho }`$, the other components of $`\mathrm{\Omega }_{3,1}`$ involve the scalar component of the matter multiplet $`M`$ and require more care since we have already precisely defined the four-dimensional field $`M`$ by its coupling to the bulk fields. To obtain the correct relations, a detour is helpful.
As already explained, the gauge kinetic action (3.12) also generates the four-dimensional kinetic terms for the matter multiplet $`M`$. In the four-dimensional effective Lagrangian, these contributions arise as the highest component of the ‘matter Chern-Simons multiplet’ $`\overline{M}e^AM`$, which includes $`2(D_\mu \overline{M})(D^\mu M)`$. This multiplet also contains in its vector component the matter Chern-Simons form
$$\mathrm{\Omega }_\mu ^M=i\overline{M}(D_\mu M)i(D_\mu \overline{M})M.$$
This is completely similar to the gauge Chern-Simons multiplet which includes gauge kinetic terms in its d-component and $`(\mathrm{\Omega }_3)_{\mu \nu \rho }`$ in its vector component, as indicated by expressions (4.27). A direct computation of the relation between kinetic terms due to the action (3.12), and the highest component of $`\overline{M}e^AM`$ delivers then the relation between $`(\mathrm{\Omega }_{3,1})_{\mu i\overline{j}}`$ and this multiplet, by four-dimensional supersymmetry. A similar operation gives the relation between $`(\mathrm{\Omega }_{3,1})_{ijk}`$ and the superpotential multiplet $`\alpha M^3`$. The relations are
$$\begin{array}{ccc}\hfill (\mathrm{\Omega }_{3,1})_{\mu i\overline{j}}& =& \frac{1}{6\kappa ^2}\left[(D_\mu \overline{M})M\overline{M}(D_\mu M)\right]\delta _{i\overline{j}}=\frac{i}{6\kappa ^2}\delta _{i\overline{j}}[\overline{M}e^AM]_{\mathrm{vector}\mathrm{component}},\hfill \\ \multicolumn{3}{c}{}\\ \hfill (\mathrm{\Omega }_{3,1})_{ijk}& =& \frac{1}{3\kappa ^3}\alpha \overline{M}^3ϵ_{ijk}=\frac{1}{3\kappa ^3}ϵ_{ijk}[\alpha \overline{M}^3]_{\mathrm{scalar}\mathrm{component}},\hfill \end{array}$$
(4.28)
together with the last equation (4.27) which applies to all gauge Chern-Simons forms.
With the complete identification of the massless modes of the Chern-Simons three-forms, we are equipped for translating the Calabi-Yau reduction of the Green-Schwarz counterterm (4.26) in a four-dimensional supergravity density formula. After straightforward manipulations of Eq. (4.26), we obtain
$$\begin{array}{ccc}\hfill _{GS}& =& \frac{1}{192(4\pi )^5}Iϵ^{\mu \nu \rho \sigma }\frac{\lambda ^2}{V_6}(_\mu a)(\mathrm{\Omega }_{E_6}\mathrm{\Omega }_{E_8})_{\nu \rho \sigma }\hfill \\ \multicolumn{3}{c}{}\\ & & +\frac{i}{384(4\pi )^5\kappa ^2}Iϵ^{\mu \nu \rho \sigma }\frac{\lambda ^2}{V_6}(_\mu C_{\nu \rho 4})(MD_\sigma \overline{M}\overline{M}D_\sigma M)\hfill \\ \multicolumn{3}{c}{}\\ & & +\frac{i\alpha ^2}{216(4\pi )^5}\frac{V_6}{\kappa ^6}ϵ^{\mu \nu \rho \sigma }\frac{\lambda ^2}{\kappa ^2V_6}(_\mu C_{\nu \rho 4})\left[\overline{M}^3(_\sigma M^3)(_\sigma \overline{M}^3)M^3\right]+\mathrm{}\hfill \\ \multicolumn{3}{c}{}\\ & =& \frac{1}{384(4\pi )^5}Iϵ^{\mu \nu \rho \sigma }T_\mu (\mathrm{\Omega }_{E_6}\mathrm{\Omega }_{E_8})_{\nu \rho \sigma }\hfill \\ \multicolumn{3}{c}{}\\ & & \frac{i}{96(4\pi )^5}Iϵ^{\mu \nu \rho \sigma }(_\nu b_{\rho \sigma })(MD_\mu \overline{M}\overline{M}D_\mu M)\hfill \\ \multicolumn{3}{c}{}\\ & & \frac{i\alpha ^2}{54(4\pi )^5}\frac{V_6}{\kappa ^6}ϵ^{\mu \nu \rho \sigma }(_\nu b_{\rho \sigma })\left[\overline{M}^3(_\mu M^3)(_\mu \overline{M}^3)M^3\right]+\mathrm{}\hfill \end{array}$$
(4.29)
where the dots indicate the terms required by $`N=1`$ supersymmetry and $`I`$ is the dimensionless integral
$$I=\kappa ^2_{CY}𝑑V_6\delta _{i\overline{i}}ϵ^{ijk}ϵ^{\overline{ijk}}\mathrm{tr}R^2_{jk\overline{jk}},$$
in terms of the $`(1,1)`$ (the metric tensor), $`(3,0)`$ and $`(0,3)`$ Calabi-Yau tensors and the background $`\mathrm{tr}R^2`$.
Using Eqs. (4.28), we can write the Green-Schwarz counterterm in superfield form as:
$$\begin{array}{ccc}\hfill _{GS}& =& \frac{I}{48(4\pi )^5}\left[(V_T+2\overline{M}e^AM)(\mathrm{\Omega }^1\mathrm{\Omega }^2)\right]_D\frac{I}{48(4\pi )^5}\left[V\overline{M}e^AM\right]_D\hfill \\ \multicolumn{3}{c}{}\\ & & +\frac{1}{27(4\pi )^5}\frac{V_6}{\kappa ^6}\left[V|\alpha M^3|^2\right]_D.\hfill \end{array}$$
(4.30)
Comparing the three terms of $`_{GS}`$ with the corresponding parts of (4.10), we can express the coefficients $`\beta ^a`$, $`\delta `$ and $`ϵ`$ in terms of the Calabi-Yau dependent integral $`I`$:
$$\delta =\beta ^1=\beta ^2=\frac{I}{96(4\pi )^5},ϵ=\frac{1}{27(4\pi )^5}\frac{V_6}{\kappa ^6}.$$
(4.31)
The calculation predicts then $`\beta ^1=\beta ^2=\delta `$, a result already obtained in references , for instance. Once again, we stress that we have obtained next-order corrections to the effective Wilson Lagrangian. As argued earlier in paragraph 4.1, while we expect $`\beta ^1`$ and $`\beta ^2`$ to have physical significance as coefficients of the modulus-dependent threshold corrections, the parameter $`\delta `$ is not necessarily a physical quantity. To decide of its relevance, a calculation at the same order of threshold corrections in the effective action should be performed, but this computation requires a detailed knowledge of the charged matter spectrum and couplings.
## 5 Conclusions
In this paper, we have deduced the structure of the four-dimensional $`N=1`$ effective (wilsonnian) supergravity describing the universal massless sector of M-theory compactified on $`(\text{CY})_3\times S^1/𝐙_2`$. The theory depends on three categories of multiplets: ‘M-theory multiplets’ $`V`$, $`V_T`$ and $`W`$ describe the degrees of freedom of the M-theory four-index tensor, ‘source multiplets’ $`\mathrm{\Omega }^a`$, $`\overline{M}e^AM`$ and $`\alpha M^3`$ are related to the source terms in M-theory Bianchi identities and ‘Lagrange multiplets’ $`S`$, $`L_T`$ and $`U`$ impose by their field equations the Bianchi identities. In addition, the multipliers $`S`$ and $`L_T`$ generate the four-dimensional axion-tensor duality which is known to be an important ingredient of the formulation of the string dilaton beyond the lowest order.
An effective supergravity similar to expression (4.10) is in principle valid for generic compactifications of M-theory with unbroken $`N=1`$ supersymmetry in four dimensions. One needs to identify the appropriate supermultiplets appearing as sources in the Bianchi identities generated by $`S`$, $`L_T`$ and $`U`$. Only the values of coefficients like $`\beta ^a`$, $`\delta `$ and $`ϵ`$ depend on the detailed geometry of the compact space. This method is especially useful in deriving contributions to the effective Lagrangian due to non-perturbative states like M-theory five-brane degrees of freedom or condensates. Some of these more general $`N=1`$ vacua will be studied in a forthcoming publication .
Acknowledgements
The authors have benefited from discussions with A. Bilal, C. Kounnas, A. Lukas, D. Lüst and B. Ovrut. This research was supported in part by the European Union under the TMR contract ERBFMRX-CT96-0045, the Swiss National Science Foundation and the Swiss Office for Education and Science.
## Appendix: Notations and conventions
Metrics and coordinates:
The space-time metric has signature $`(,+,+,\mathrm{},+).`$
For coordinates, our notation is:
| $`D=11`$ curved space-time: | $`x^M`$ | $`M=0,\mathrm{},10`$ |
| --- | --- | --- |
| $`D=10`$ curved space-time: | $`x^A`$ | |
| $`D=4`$ curved space-time: | $`x^\mu `$ | $`\mu =0,1,2,3`$ |
| $`S_1/𝐙_2`$ direction: | $`x^4`$ | |
| Calabi-Yau directions: | $`x^a`$ | $`a=5,\mathrm{},10`$ |
| Calabi-Yau complex (Kähler) coordinates: | $`z^i`$, $`\overline{z}^{\overline{i}}`$ | $`i=1,2,3`$ |
For reduction purposes, we simply use
$$z^l=\frac{1}{\sqrt{2}}\left(x^l+ix^{l+3}\right),\overline{z}^{\overline{l}}=\frac{1}{\sqrt{2}}\left(x^lix^{l+3}\right),l=1,2,3.$$
$`ϵ_{ijk}`$ is the $`SU(3)`$–invariant Calabi-Yau tensor such that $`ϵ_{123}=ϵ_{\overline{123}}=1`$.
Antisymmetric tensors:
Antisymmetrization of $`n`$ indices has unit weight:
$$A_{[M_1\mathrm{}M_n]}=\frac{1}{n!}\left(A_{M_1\mathrm{}M_n}\pm (n!1)\mathrm{permutations}\right).$$
Differential forms
For a $`p`$–index antisymmetric tensor, we define
$$A^{(p)}=\frac{1}{p!}A_{M_1\mathrm{}M_p}dx^{M_1}\mathrm{}dx^{M_p}.$$
Then,
$$\begin{array}{c}A^{(p)}B^{(q)}=\frac{1}{p!q!}A_{M_1\mathrm{}M_p}B_{M_{p+1}\mathrm{}M_{p+q}}dx^{M_1}\mathrm{}dx^{M_{p+q}}=C^{(p+q)},\hfill \\ \\ C_{M_1\mathrm{}M_{p+q}}=\frac{(p+q)!}{p!q!}A_{[M_1\mathrm{}M_p}B_{M_{p+1}\mathrm{}M_{p+q}]}.\hfill \end{array}$$
The exterior derivative is $`d=_Mdx^M`$. The curl $`F^{(p+1)}=dA^{(p)}`$ of a $`p`$-form reads then
$$\begin{array}{ccc}\hfill dA^{(p)}& =& \frac{1}{p!}(_MA_{N_1\mathrm{}N_p})dx^Mdx^{N_1}\mathrm{}dx^{N_p}\hfill \\ \multicolumn{3}{c}{}\\ & =& \frac{1}{(p+1)!}F_{M_1\mathrm{}M_{p+1}}dx^{M_1}\mathrm{}dx^{M_{p+1}},\hfill \\ \multicolumn{3}{c}{}\\ \hfill F_{M_1\mathrm{}M_{p+1}}& =& (p+1)_{[M_1}A_{M_2\mathrm{}M_{p+1}]}\hfill \\ \multicolumn{3}{c}{}\\ & =& _{M_1}A_{M_2\mathrm{}M_{p+1}}\pm p\mathrm{cyclic}\mathrm{permutations}.\hfill \end{array}$$
The volume form in $`D`$ space-time dimensions is $`dx^{M_1}\mathrm{}dx^{M_D}=ϵ^{M_1\mathrm{}M_D}d^Dx`$.
We use analogous conventions for forms in four space-time dimensions.
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# 1 Introduction
## 1 Introduction
The $`\varphi `$-meson was discovered many years ago as a $`K\overline{K}`$ resonance . Its decay is dominated by the two $`K\overline{K}`$ decay modes which proceed through Zweig-rule allowed strong interactions. The ratio $`R\varphi K^+K^{}/K^0\overline{K^0}`$ has been measured in a variety of independent experiments using different $`\varphi `$-production mechanisms. Among these, the cleanest one is electron-positron annihilation around the $`\varphi `$ resonance peak, i.e. the reactions $`e^+e^{}\varphi K^+K^{}/K^0\overline{K^0}`$, which have been accurately measured at Novosibirsk quite recently and are the object of intense investigation at the Frascati $`\mathrm{\Phi }`$-factory . With as much as $`8\times 10^6\varphi `$’s on tape, the KLOE experiment at DA$`\mathrm{\Phi }`$NE can be expected to measure the above ratio $`R`$ with a statistical accuracy of the order of the per mille. In view of this, we wish to discuss the theoretical expectations and compare them with the most recent determinations for this ratio.
In the following we shall first review the present experimental situation, then compare it with the naïve expectations from isospin symmetry and phase space considerations thus observing that a disagreement seems to exist. Contributions arising from electromagnetic radiative corrections and $`m_um_d`$ isospin breaking effects are analyzed and shown to bring the observed discrepancy to be more than three standard deviations. Various additional theoretical improvements on our analysis, such as the use of vector-meson dominated electromagnetic form-factors, the modification of the strong vertices and the inclusion of rescattering effects through the scalar resonances $`f_0(980)`$ and $`a_0(980)`$ using the charged kaon loop model, are also examined and shown not to change in any substantial way our results which imply a clear discrepancy between theory and the available data.
The first combined measurement of the four major $`\varphi `$ decay modes in a single $`e^+e^{}`$ dedicated experiment has been performed quite recently with the general purpose detector CMD-2 at the upgraded $`e^+e^{}`$ collider VEPP-2M at Novosibirsk . Having a single experiment normalized to almost 100% of decay modes implies a reduction of systematic errors, and the following branching ratios (BR) and errors from VEPP-2M are quoted:
$$\begin{array}{c}\mathrm{BR}(\varphi K^+K^{})=(49.2\pm 1.2)\%,\hfill \\ \mathrm{BR}(\varphi K^0\overline{K}^0)=(33.5\pm 1.0)\%,\hfill \end{array}$$
(1)
leading to
$$R_{\mathrm{exp}}\frac{\mathrm{BR}(\varphi K^+K^{})}{\mathrm{BR}(\varphi K^0\overline{K}^0)}=1.47\pm 0.06.$$
(2)
All these results were in agreement with the average values quoted in the then available PDG 1994 compilation :
$$\begin{array}{c}\mathrm{BR}(\varphi K^+K^{})=(49.1\pm 0.9)\%\hfill \\ \mathrm{BR}(\varphi K^0\overline{K}^0)=(34.3\pm 0.7)\%\hfill \end{array}\}R_{\mathrm{exp}}=1.43\pm 0.04.$$
(3)
The current PDG edition , now including the above VEPP-2M data, quotes
$$\begin{array}{c}\mathrm{BR}(\varphi K^+K^{})=(49.1\pm 0.8)\%\hfill \\ \mathrm{BR}(\varphi K^0\overline{K}^0)=(34.1\pm 0.6)\%\hfill \end{array}\}R_{\mathrm{exp}}=1.44\pm 0.04,$$
(4)
as a result of a global fit, which appears as a very stable result, established with a 3% error. In the same PDG edition, one can also find $`R_{\mathrm{direct}}=1.35\pm 0.06`$, as the averaged result of the various experiments measuring the ratio $`\varphi K^+K^{}/K^0\overline{K}^0`$ directly. A reduction of these errors can be expected from DA$`\mathrm{\Phi }`$NE, where the KLOE experiment has already collected $`8\times 10^6\varphi `$-mesons. Like in the case of the CMD-2 detector, all the main decay modes of the $`\varphi `$ will be measured by the same apparatus and this could bring the systematic errors to a minimum, while the statistics will allow to bring the statistical error well below the 1% level. Our discussion centers around this ratio $`R`$ and the possible interest in studying it with a much reduced experimental error.
We shall approach this discussion by starting with the most naïve result for the above ratio $`R`$, i.e. $`R=1`$, which follows from assuming that these $`\varphi K\overline{K}`$ decay modes proceed exclusively via the strong interaction dynamics in the good isospin limit $`m_u=m_d`$ and ignoring phase space differences. The mass difference between neutral and charged kaons —which includes both isospin breaking effects ($`m_um_d`$) and electromagnetic (photonic) contributions— considerably increases this too-naïve prediction via the (purely kinematical) phase-space factor. Assuming now perfect isospin symmetry only for the strong interaction dynamics (equal couplings for $`\varphi K^+K^{}`$ and $`\varphi K^0\overline{K}^0`$) and knowing that $`\varphi K\overline{K}`$ are $`P`$-wave decay modes of a narrow resonance, one necessarily has
$$R=\frac{\left(1\frac{4m_{K^+}^2}{M_\varphi ^2}\right)^{3/2}}{\left(1\frac{4m_{K^0}^2}{M_\varphi ^2}\right)^{3/2}}=1.528,$$
(5)
with negligible errors coming from the mass values quoted in the PDG. The phase-space correction thus pushes the ratio $`R`$ two standard deviations above its experimental value (4). This kinematical correction is exceptionally large because of the vicinity of the $`\varphi `$ mass to the $`K\overline{K}`$ thresholds, which translates into considerably large differences between the charged and neutral kaon momenta (or velocities, $`v_+/v_0=0.249/0.216=1.152`$), a difference which is further increased to its third power in such $`P`$-wave decay modes.
This two-$`\sigma `$ discrepancy between experiments and the theoretical tree level predictions obviously claims for further corrections. The most immediate of such corrections is due to electromagnetic radiative effects on the ratio $`R`$, which affect the numerator but not the denominator, and which will be discussed in the next section.
## 2 Electromagnetic radiative corrections
Electromagnetic radiative corrections are frequently ignored when dealing with strong decays. In our case, they could be relevant since, although small, they affect the charged decay mode but not the neutral one, and, in order to solve the discrepancy in the ratio $`R`$ under consideration, only a few per cent correction is needed. Many years ago they were already considered by Cremmer and Gourdin who found a positive correction of the order of 4% to the prediction in Eq. (5), thus enlarging that discrepancy. The dominant contribution was found to arise from the so-called Coulomb term which is positive for $`\varphi K^+K^{}`$ and rather large because of the small kaon velocities $`v_\pm =0.249`$. A similar increase of the ratio $`R`$ (some 5%) by radiative corrections is expected by the experimentalists at VEPP-2M , whose quoted result is inclusive of any vertex correction. If we include this correction in the theoretically predicted ratio, the final result for the radiatively corrected ratio is then $`R1.59`$ , in agreement with still another independent analysis by Pilkuhn leading to $`R`$ in the range $`1.52`$$`1.61`$ . To better qualify these statements, we shall now examine in detail the contribution of such corrections to the ratio $`R`$.
We have recalculated the electromagnetic radiative corrections to $`\varphi K^+K^{}`$ along the lines of Ref. . For the charged amplitude we start with the usual and simplest tree level expression $`A_0(\varphi K^+K^{})=g_0ϵ_\mu (p_+p_{})^\mu `$, where $`g_0`$ is the uncorrected strong coupling constant for $`\varphi K\overline{K}`$, $`ϵ_\mu `$ is the $`\varphi `$ polarization and $`p_\pm `$ are the kaon four-momenta. As is well known, the various contributions to the radiative corrections can be grouped in two parts. The first part comprises one-loop corrections to the uncorrected amplitude $`A_0(\varphi K^+K^{})`$. This part contains three vertex diagrams with one virtual photon exchanged between the two charged-kaons or between the $`\varphi K^+K^{}`$ vertex and each charged-kaon. In addition, it also contains wave-function renormalization of external kaon lines that render the whole amplitude ultraviolet finite<sup>1</sup><sup>1</sup>1 Notice that Eq. (19) in Ref. contains a small imaginary part while it is supposed to be the real part of the one-loop amplitude.. The second part is needed to cancel the infrared divergence. It contains three real-photon emission diagrams which are order $`\sqrt{\alpha }`$. Adding these two parts we find the complete order $`\alpha `$ corrective factor to the $`\varphi K^+K^{}`$ decay width
$$\begin{array}{c}1+C_f+\beta _f\mathrm{log}\frac{2\mathrm{\Delta }E}{m_{K^+}}1+\frac{\alpha }{\pi }\{\frac{1+v^2}{2v}\pi ^22(1+\mathrm{log}\frac{2\mathrm{\Delta }E}{m_{K^+}})(1+\frac{1+v^2}{2v}\mathrm{log}\frac{1v}{1+v})\hfill \\ \frac{1}{v}\mathrm{log}\frac{1v}{1+v}\frac{1+v^2}{4v}\mathrm{log}\frac{1v}{1+v}\mathrm{log}\frac{1v^2}{4}\frac{1+v^2}{2v}\left[\mathrm{Li}_2\left(\frac{2v}{1+v}\right)\mathrm{Li}_2\left(\frac{2v}{1v}\right)\right]\hfill \\ +\frac{1+v^2}{2v}[\mathrm{Li}_2\left(\frac{1+v}{2}\right)\mathrm{Li}_2\left(\frac{1v}{2}\right)]\frac{1+v^2}{v}[\mathrm{Li}_2(v)\mathrm{Li}_2(v)]\},\hfill \end{array}$$
(6)
where $`v=\sqrt{14m_{K^+}^2/M_\varphi ^2}`$ is the kaon velocity and $`\mathrm{\Delta }E`$ stays for the photon energy resolution. For $`\mathrm{\Delta }E=1`$ MeV the correction (6) amounts to a 4.2% increase. Taking for $`\mathrm{\Delta }E`$ the maximal available photon energy (32.1 MeV, not far from the energy resolution in the KLOE detector at DA$`\mathrm{\Phi }`$NE, which is $``$ 20 MeV) makes no substantial difference as the main contribution comes from the Coulomb term, the first one inside the brackets.
The above discussion ignores the fact that what is actually measured at VEPP-2M and at DA$`\mathrm{\Phi }`$NE is the ratio
$$R_{e^+e^{}}\frac{\sigma (e^+e^{}\varphi K^+K^{})}{\sigma (e^+e^{}\varphi K^0\overline{K}^0)},$$
(7)
and that radiative corrections to $`R`$ correspond to consider the ratio of the radiatively corrected cross-sections which appear at the numerator and denominator of $`R_{e^+e^{}}`$. In addition to consider both initial and final state corrections, a complete treatment also requires to discuss the presence of the $`\varphi `$ resonance and the associated distortion of the cross-sections . At the numerator, radiative corrections include virtual corrections as well as emission of soft unobserved photons, both from the initial and final states, with no interference between initial and final state radiation for an inclusive measurement (i.e. in a measurement that does not distinguish the charges of the kaons) . For the cross-section at the denominator, there are only initial state radiative corrections since the final kaons are neutral. In the absence of final state radiation, the presence of a narrow resonance like the $`\varphi `$ in the intermediate state introduces large double logarithms which can be resummed and factorized in an expression like
$$\left(\frac{\mathrm{\Gamma }_\varphi }{M_\varphi }\right)^{\beta _i}(1+C_i),$$
(8)
where $`\beta _i=\frac{2\alpha }{\pi }\left(\mathrm{log}\frac{s}{m_e^2}1\right)`$ is the initial state radiation factor and $`C_i`$ is the finite part of the initial virtual and soft photon corrections, which survives after the cancellation of the infrared divergence and the exponentiation of the large resonant dependent factors. The same factor for initial state radiation appears both at numerator and denominator, and since there is no interference between initial and final state radiation, the real soft-photon radiative corrections to the initial state cancel out in the ratio (7). In principle, one should also resum the contributions coming from final state radiation but the final state radiative factor $`\beta _f=\frac{2\alpha }{\pi }\left(\frac{1+v^2}{2v}\mathrm{log}\frac{1+v}{1v}1\right)3.9\times 10^4`$ is very small and resummation in this case is irrelevant. One then obtains the following expression for the ratio $`R_{e^+e^{}}`$ as defined in Eq. (7):
$$R_{e^+e^{}}=\frac{\mathrm{\Gamma }(\varphi K^+K^{})}{\mathrm{\Gamma }(\varphi K^0\overline{K}^0)}\frac{1+C_i+C_f+\beta _f\mathrm{log}\frac{2\mathrm{\Delta }E}{m_{K^+}}}{1+C_i}.$$
(9)
Since $`C_i\frac{\alpha }{\pi }\left(\frac{3}{2}\mathrm{log}\frac{s}{m_e^2}+\frac{\pi ^2}{3}2\right)5.6\times 10^2`$ , one can expand the denominator in Eq. (9), canceling the $`C_i`$ term and remaining with the final state correction terms $`C_f`$ and $`\beta _f`$ given explicitly in Eq. (6). We thus conclude that one is justified in using the expressions as above and that the conventional treatment of radiative corrections increases the previous two-$`\sigma `$ discrepancy between experiment and theory for the ratio $`R`$ to the level of three standard deviations.
## 3 $`SU(2)`$-breaking in $`\varphi K\overline{K}`$ vertices
The $`\varphi K^+K^{}`$ and $`\varphi K^0\overline{K}^0`$ vertices are not equal (and thus do not cancel in the ratio $`R`$) once $`SU(2)`$-breaking effects are taken into account. The way $`SU(2)`$-breaking is usually introduced in the effective lagrangians is the same as for $`SU(3)`$-breaking, namely, via quark mass differences. In the latter $`SU(3)`$ case, an improved description of the vector-meson couplings to two pseudoscalar-mesons can easily be achieved as shown, for example, in Refs. . But the situation is by far less convincing when turning to the much smaller $`SU(2)`$-breaking effects . The essential feature —common to most models— is that the dynamics of these flavour symmetry breakings suppress the creation of heavier $`q\overline{q}`$ pairs. In the $`\varphi K^+K^{}`$ and $`\varphi K^0\overline{K}^0`$ vertices, one needs to produce a $`u\overline{u}`$ and a $`d\overline{d}`$ pair, respectively. Since the latter is heavier, the $`\varphi K^0\overline{K}^0`$ decay is further suppressed and then the ratio $`R`$ is further increased. To be somewhat more precise, we will consider two recent and independent models dealing quite explicitly with such kind of effects .
In the $`SU(3)`$-breaking treatment of $`VP_1P_2`$ vertices by Bijnens et al. , these decays proceed through two independent terms containing the relevant vector and pseudoscalar masses ($`M_V`$ and $`m_{1,2}`$) and thus incorporating quark-mass breaking effects. In the notation of Ref. , to which we refer for details, these $`VP_1P_2`$ couplings are then proportional to
$$M_V^2\left(g_V+2\sqrt{2}f_\chi \frac{m_1^2+m_2^2}{M_V^2}\right).$$
(10)
For the $`\varphi K^+K^{}`$ and $`\varphi K^0\overline{K}^0`$ coupling constants, the uncorrected strong coupling constant $`g_0`$ becomes, respectively,
$$\frac{M_\varphi ^2}{2\sqrt{2}g_0f^2}\left(1+4\sqrt{2}\frac{f_\chi }{g_V}\frac{m_{K^+,K^0}^2}{M_\varphi ^2}\right),$$
(11)
with the pion decay constant $`f92`$ MeV. One then obtains the ratio
$$\frac{g_{\varphi K^+K^{}}}{g_{\varphi K^0\overline{K}^0}}1+4\sqrt{2}\frac{f_\chi }{g_V}\frac{m_{K^+}^2m_{K^0}^2|_{m_um_d}}{M_\varphi ^2}1.01,$$
(12)
where we have used $`m_{K^+}^2m_{K^0}^2|_{m_um_d}\mathrm{6\; 10}^3`$ GeV<sup>2</sup> for the non-photonic kaon mass difference and the estimate $`\frac{f_\chi }{g_V}\frac{1}{3}`$ obtained in Ref. when fitting the $`\rho \pi \pi `$ and $`K^{}K\pi `$ decay widths.
Similarly, in the independent treatment of $`SU(3)`$ symmetry breaking , some relevant $`VP_1P_2`$ couplings are given by
$$\begin{array}{ccc}\hfill g_{\rho \pi \pi }& =& \sqrt{2}g,\hfill \\ \hfill g_{\varphi K^+K^{}}& =& g_{\varphi K^0\overline{K}^0}=g(1+2c_V)(1c_A),\hfill \end{array}$$
(13)
with $`c_V0.28`$ and $`c_A0.36`$ (see Ref. for notation and details) mimicking the $`SU(3)`$ mass difference effects discussed in the previous approach . The transition from $`SU(3)`$\- to $`SU(2)`$-breaking offers no difficulties. One now obtains
$$\frac{g_{\varphi K^+K^{}}}{g_{\varphi K^0\overline{K}^0}}1\frac{m_{K^+}^2m_{K^0}^2|_{m_um_d}}{m_K^2m_\pi ^2}c_A1.01.$$
(14)
As in the approach of Ref. , these $`SU(2)`$-breaking corrections work in the undesired direction and the discrepancy between theory and experiment for the ratio $`R`$ increases by an additional 2%.
An independent $`SU(2)`$-breaking effect can arise from $`\rho `$-$`\varphi `$ mixing. This is both isospin and Zweig-rule violating, and should therefore lead to rather tiny corrections. Indeed, in this context one can immediately obtain the following relation among coupling constants<sup>2</sup><sup>2</sup>2 Notice that this isospin relation not only accounts for $`\rho (770)`$$`\varphi `$ mixing effects but also for those between $`\varphi `$ and any other higher mass isovector $`\rho `$-like resonance.: $`g_{\varphi K^+K^{}}g_{\varphi K^0\overline{K}^0}=g_{\varphi \pi ^+\pi ^{}}`$, with a small value for the $`g_{\varphi \pi ^+\pi ^{}}`$ coupling coming from the observed smallness of the $`\varphi \pi ^+\pi ^{}`$ branching ratio ($`𝒪(10^4)`$ ) in spite of its much larger phase space. A more quantitative estimate is now possible thanks to the recent data on $`e^+e^{}\varphi \pi ^+\pi ^{}`$ coming from VEPP-2M . These data describe the pion form factor around the $`\varphi `$ peak, $`F(sM_\varphi ^2)`$, in terms of the complex parameter $`Z`$ by the expression
$$F(s)\left(1\frac{ZM_\varphi \mathrm{\Gamma }_\varphi }{M_\varphi ^2siM_\varphi \mathrm{\Gamma }_\varphi }\right).$$
(15)
This $`Z`$, in turn, can be easily related to $`ϵ_{\varphi \rho }`$, the complex parameter describing the amount of $`\rho `$-like (or $`(u\overline{u}d\overline{d})/\sqrt{2}`$) contamination in the $`\varphi `$ wave function. One finds
$$ϵ_{\varphi \rho }\frac{f_\varphi }{f_\rho }\frac{\mathrm{\Gamma }_\varphi }{M_\varphi }F(s=M_\varphi ^2)Z,$$
(16)
where the first coefficient $`\frac{f_\varphi }{f_\rho }\frac{3}{\sqrt{2}}`$ is the well-known ratio of $`\varphi `$-$`\gamma `$ to $`\rho `$-$`\gamma `$ couplings. One finally obtains
$$\frac{g_{\varphi K^+K^{}}}{g_{\varphi K^0\overline{K}^0}}1\sqrt{2}\mathrm{}(ϵ_{\varphi \rho })1.001,$$
(17)
where an average of the values for $`Z`$ in Ref. and the parametrization of $`F(s=M_\varphi ^2)`$ from Ref. have been used in the final step. This time the correction is tiny and the accuracy of our estimate is rather rough, but again it tends to increase the discrepancy on the ratio $`R`$.
## 4 Further attempts
Since the discrepancy between the theoretical and experimental value for $`R`$ remains (or has even been increased by some additional 2% due to the $`SU(2)`$-breaking effects just discussed), we have tried to improve our analysis in different aspects. First, we have taken into account that the couplings of photons to kaons, rather than being point-like (as assumed in our previous and conventional treatment of radiative corrections), are known to be vector-meson dominated . Accordingly, we have redone the calculation performed in Sec. 2 including the corresponding electromagnetic (vector-meson dominated) kaon form-factors. Now, not only the decay mode $`\varphi K^+K^{}`$ can be affected but also the $`\varphi K^0\overline{K}^0`$ one due to the $`\rho `$, $`\omega `$ and $`\varphi `$ mass differences. For the $`\varphi K^+K^{}`$ case, the contribution of the charged kaon form-factor modifies the point-like result for $`\mathrm{\Gamma }(\varphi K^+K^{})`$ by $`\mathrm{2\; 10}^3`$. For the case of $`\varphi K^0\overline{K}^0`$, a vanishing effect will be obtained in the limit of exact $`SU(3)`$ symmetry, and a fraction of the preceding one if $`SU(3)`$-broken masses are used. In both cases, the effect of kaon form-factors on real-photon emission diagrams is null. So then, the additional net effect of electromagnetic kaon form-factors on the ratio $`R`$ leads to a modification of the point-like radiative corrections result of Sec. 2 by some per mille and is thus fully negligible.
A second and independent possibility consists in adopting a different framework for $`VPP`$ decays. This is usually done in terms of more general effective lagrangians with $`VPP`$ vertices containing two derivatives of the pseudoscalar fields instead of a single one as in our previous discussion. The radiative decay $`\rho \pi ^+\pi ^{}\gamma `$ —quite similar to the processes we are considering— has been quite recently analyzed in this modern context in Ref. . The two relevant coupling constants ($`F_V`$ and $`G_V`$, in the notation of Ref. ) and their relative sign can be fixed to the canonical values $`F_V=2G_V=\sqrt{2}f_\pi `$ thanks to the experimental data for $`\rho \pi ^+\pi ^{}\gamma `$ and other $`\rho `$ meson processes . As discussed in Ref. , a good description of these data is then achieved in terms of an amplitude that coincides with the one previously introduced in Ref. , and which originated from the simple one-derivative $`VPP`$ vertices used by Ref. as well as in our recalculation in Sec. 2. In other words, both types of effective lagrangians lead to exactly the same real-photon emission amplitudes once the coupling constants are properly fixed. This is also true for the other corrections concerning one-loop effects: for the canonical value $`F_V=2G_V`$ one reobtains precisely our previous expression in Eq. (6).
A third attempt includes the effect of final $`K\overline{K}`$ rescattering through scalar resonances. It is well known that the charged kaons emitted in $`\varphi K^+K^{}`$ are always accompanied by soft photons. In the case of single photon emission, the $`K^+K^{}`$ system is found to be in a $`J^{PC}=0^{++}`$ or $`2^{++}`$ state with an invariant mass just below the $`\varphi `$ mass. The presence of the $`J^{PC}=0^{++}`$ scalar resonances $`f_0(980)`$ and $`a_0(980)`$, with masses and decay widths that cover the invariant mass range of interest (from the $`K\overline{K}`$ threshold to the $`\varphi `$ mass) , would suggest that rescattering effects could be important<sup>3</sup><sup>3</sup>3Rescattering effects from $`2^{++}`$ states are suppressed because the nearest tensorial resonances, $`f_2(1270)`$ and $`a_2(1310)`$, are well above the $`\varphi `$ mass .. We have computed these rescattering effects through the exchange of the $`f_0`$ and $`a_0`$ using the charged kaon loop model . In this model, the $`\varphi `$ decays into a $`K^+K^{}`$ system that emits a photon (from the charged kaon internal lines and from the $`\varphi K^+K^{}`$ vertex) before rescattering into a final $`K^+K^{}`$ or $`K^0\overline{K}^0`$ state through the propagation of $`f_0`$ and $`a_0`$ resonances. If the emitted soft photon is unobserved, the process $`\varphi K^+K^{}(\gamma )f_0/a_0(\gamma )K^+K^{}(\gamma )`$ or $`K^0\overline{K}^0(\gamma )`$ contributes to the ratio $`R`$, both at the numerator and denominator. In order to calculate these effects, one needs an estimate of the coupling constant $`g_{SK\overline{K}}`$, where $`S`$ is either the $`f_0`$ or the $`a_0`$. Recent measurements of the $`\varphi f_0\gamma `$ and $`a_0\gamma `$ decay modes at VEPP-2M are consistent with the predictions of the charged kaon loop model for values of the above couplings given by
$$\frac{g_{f_0K\overline{K}}^2}{4\pi }=(1.48\pm 0.32)\mathrm{GeV}^2,\frac{g_{a_0K\overline{K}}^2}{4\pi }=(1.5\pm 0.5)\mathrm{GeV}^2.$$
(18)
We have then found that the contribution of these kaon loops to the $`\mathrm{BR}(\varphi K^+K^{}(\gamma ))`$ is $`𝒪(10^7)`$, while for $`\mathrm{BR}(\varphi K^0\overline{K}^0(\gamma ))`$ is $`𝒪(10^9)`$. For charged kaons in the final state, there is an additional contribution from the interference between the soft-bremsstrahlung and the scalar amplitudes. This contribution is given by
$$\begin{array}{c}\mathrm{\Gamma }_{\mathrm{int}}(\varphi K^+K^{}(\gamma ))=\frac{4}{3}\alpha \frac{g_{\varphi K^+K^{}}^2}{4\pi }\frac{g_{f_0K^+K^{}}^2}{4\pi }\frac{M_\varphi }{2\pi ^2m_{K^+}^2}_0^{\mathrm{\Delta }E}d\omega \omega \times \hfill \\ \mathrm{}\left(I(a,b)\left[\frac{1}{D_{f_0}(m)}+\frac{g_{a_0K^+K^{}}^2}{g_{f_0K^+K^{}}^2}\frac{1}{D_{a_0}(m)}\right]\right)\left(w+\frac{1v^2}{2}\mathrm{log}\frac{1w}{1+w}\right),\hfill \end{array}$$
(19)
where $`I(a,b)`$ is the kaon loop integral defined in Refs. (with $`aM_\varphi ^2/m_{K^+}^2`$ and $`bm^2/m_{K^+}^2`$), $`D_{f_0/a_0}(m)`$ are the scalar propagators and $`w=\sqrt{14m_{K^+}^2/m^2}`$, with $`m=M_\varphi \sqrt{12\omega /M_\varphi }`$ being the invariant mass of the $`K\overline{K}`$ system and $`\omega `$ the photon energy. Using the values in Eq. (18) for the scalar couplings, we find that the interference term, which contributes to $`R`$ only in the numerator, is positive and $`𝒪(10^5)`$, i.e. completely negligible in spite of being the dominant one.
Admittedly, this estimate of the $`K\overline{K}`$ rescattering effects is model dependent and affected by large uncertainties. Before concluding, we would thus like to make a few comments on possible variations on the magnitude of the scalar coupling constants and the expressions for the scalar propagators $`D_{f_0/a_0}(m)`$ which enter into our evaluation in the preceding paragraph. The values of the couplings $`g_{SK\overline{K}}`$ depend on the nature of the scalar mesons, i.e. whether they are two- or four-quark states, or $`K\overline{K}`$ molecules. The results of the $`K\overline{K}`$ molecule model, in addition to the couplings $`g_{SK\overline{K}}`$, depend upon the spatial extension of the scalar $`K\overline{K}`$ bound state, and the predictions for $`\mathrm{BR}(\varphi f_0/a_0\gamma )`$ (for the same $`g_{SK\overline{K}}`$) are always smaller than in the purely point-like case, i.e. the effects on $`R`$ tend to vanish for more extended objects . The two-quark model, irrespectively of the $`s\overline{s}`$ vs. $`(u\overline{u}+d\overline{d})/\sqrt{2}`$ quark content of the $`f_0`$, predicts too small values (see, for example, Refs. ) for the branching ratios $`\mathrm{BR}(\varphi f_0/a_0\gamma )`$ , and is unable anyway to account for the near mass degeneracy of the isoscalar $`f_0`$ and isovector $`a_0`$. On the other hand, such mass degeneracy is well understood in the four-quark model, critically reexamined very recently in Refs. . The four-quark model also predicts values for $`g_{SK\overline{K}}`$ that seem to be in agreement with the available measurements of $`\mathrm{BR}(\varphi f_0/a_0\gamma )`$ . In all cases, the different possibilities are found to modify the previously quoted sizes of the $`K\overline{K}`$ rescattering effects by at most one order of magnitude. Something similar happens with the lack of consensus on the specific form for the scalar propagators to be used in these estimates. Here the uncertainties arise because of the opening of the $`K\overline{K}`$ channels quite close to the nominal scalar masses. This translates into sharp modifications of the conventional Breit-Wigner curves and changes the size of the $`K\overline{K}`$ rescattering effects again by one order of magnitude. Although affected by large uncertainties, the contributions coming from final-state $`K\overline{K}`$ rescattering are thus found to be negligible and their effects on the ratio $`R`$ irrelevant.
## 5 Conclusions
In this letter, we have performed a discussion of the ratio $`R\varphi K^+K^{}/`$ $`K^0\overline{K^0}`$. From the experimental point of view, the value $`R_{\mathrm{exp}}=1.44\pm 0.04`$ seems to be firmly established . However, in our present theoretical analysis of this ratio $`R`$ we have failed to reproduce the value $`R_{\mathrm{exp}}`$ quoted above. In a first and conservative attempt, including isospin symmetry for the strong vertices and the appropriate phase-space factor, one obtains $`R=1.53`$ which is two $`\sigma `$’s above $`R_{\mathrm{exp}}`$. In a second step, we have also included conventional electromagnetic radiative corrections to order $`\alpha `$, thus obtaining $`R=1.59`$ and increasing the previous discrepancy up to three $`\sigma `$’s. This value confirms some existing results and has been checked to be quite independent from the details of the relevant vertices. In a third step, we have tried to correct our predictions for $`R`$ introducing various isospin breaking corrections to the $`\varphi K\overline{K}`$ coupling constants. As a result, the ratio $`R`$ is found to be further increased by some 2%, an estimate affected by rather large errors reflecting our poor knowledge on the $`SU(2)`$-breaking details. In view of all this, we have introduced final-state rescattering effects which should be dominated by the almost on-shell formation of the $`f_0(980)`$ and $`a_0(980)`$ resonances in the $`S`$-wave $`K\overline{K}`$ channel. The controversial nature of these scalar resonances allows for quite disparate estimates of their effects, but one can safely conclude that they are well below those previously mentioned. The disagreement on the ratio $`R`$ persists well above two (experimental) standard deviations. Higher statistics from DA$`\mathrm{\Phi }`$NE are expected in order to settle definitively whether the discrepancy on $`R`$ is a real problem, or final agreement between theory and experimental data can be achieved.
## Acknowledgements
We would like to thank J. Bijnens, M. Block, J. Gasser, E. Oset, M.R. Pennington and E.P. Solodov for helpful comments and clarifying discussions. Work partly supported by the EEC, TMR-CT98-0169, EURODAPHNE network. J.L. Lucio M. acknowledges partial financial support from CONACyT and CONCyTEG.
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# Effect of Exclusion of Double Occupancies in 𝑡-𝐽 Model: Extension of Gutzwiller Approximation
## I Introduction
For understanding the basic physics of high-$`T_\mathrm{c}`$ superconductivity, the two-dimensional $`t`$-$`J`$ model has been extensively studied as one of the most promising and simplified models which describe charge and spin dynamics in the CuO<sub>2</sub> plane. Mean-field theories and numerical calculations have shown that the d$`_{x^2y^2}`$ wave superconductivity (SC) takes place in a reasonable parameter region in the phase diagram ($`J/t=0.3`$ and the doping rate $`\delta <0.3`$). Therefore, as far as the d-wave SC is concerned, the experimentally observed phase diagram of high-$`T_\mathrm{c}`$ superconductors and the results in the $`t`$-$`J`$ model agree quite well.
Recently, however, there are some experiments which indicate the importance of the interplay between the antiferromagnetism and d-wave SC. The most interesting phenomenon is a realization of stripe state in some materials, which has charge ordering as well as incommensurate antiferromagnetic (AF) spin ordering. Other potential problems are the AF state induced around nonmagnetic impurities and a possibility of AF vortex cores.
In order to study these problems, it is necessary to develop a theory in which the d-wave SC and the AF correlation are treated in a reliable way. Usual mean-field theories generally overestimate the AF long-range order, so that they give unphysical results even in the uniform case. In this paper we develop a new type of analytic theory which gives a reliable estimate of the AF correlation as well as the d-wave SC. Our scheme is an extension of the Gutzwiller approximation.
In the $`t`$-$`J`$ model the double occupancy of up- and down-spin electrons at any site is prohibited. To study its ground state, it is natural to use a projected mean-field wave function in which the double occupancies are excluded. As the mean-field wave functions, a SDW mean-field solution at half filling, BCS wave functions, and a coexistent state of the AF and d-wave SC have been used. Variational energies are calculated in the variational Monte Carlo (VMC) simulations, which treat exactly the constraints of no double occupancies. Among the above trial states, the coexistent state has the lowerest variational energy in the doping region $`\delta <0.1`$.
On the other hand, the analytic theories in which the constraint is treated approximately give large discrepancies in the estimation of the AF correlation. For example, in the slave-boson mean-field theory which takes account of the constraint as an average, the AF order is too overestimated and it extends up to unphysical doping rates ($`\delta <0.2`$). Thus in the slave-boson theory, we are unable to discuss the stripe state, for example, because it is stabilized near $`\delta =0.125`$ doping where the slave-boson theory gives the AF state even for the bulk.
The Gutzwiller approximation is more advanced approximation than the simple mean-field theories in treating the constraint. Renormalizations of expectation values are introduced by comparing the statistical weighting factors in the wave functions with and without projection. As a result, the parameters $`t`$ and $`J`$ are renormalized to $`g_tt`$ and $`g_sJ`$. It has been shown that the Gutzwiller approximation gives a fairly reliable estimation of the variational energy for the pure d-wave SC state when it is compared with the VMC results. However it was shown that there is no region in the phase diagram where the AF state is stabilized. This contradicts with the VMC results.
In order to clarify the origin of these discrepancies, we investigated the VMC data and found that the Gutzwiller approximation has to be modified when the AF correlations are present. Based on these observations, we extend the Gutzwiller approximation in this paper and derive an analytic formalism for the renormalization factors which reproduces the VMC results. We show that it is important to take account of the longer-range correlations for the weighting factors, in contrast to the previous approximation where only the site-diagonal expectation values such as $`n_{i\sigma }`$ are considered. As a result, the renormalization factors for $`t`$ and $`J`$ have nonlinear dependences on the expectation values of Cooper pairs and AF moment. We think that this is the essence of the strong correlation in the sense that the renormalization appears solely from the projection operators.
The merit of the present scheme is that it can be easily applied to the inhomogeneous systems, where the VMC simulations are not so feasible. Thus a reliable analytic formulation can be given to the problems, such as the stripe state, vortex cores and states around impurities, where the interplay between the AF and d-wave SC plays an important role. Preliminary results of these problems have been described elsewhere.
The outline of this paper is as follows: In Sec. II we briefly review the formulation of the original Gutzwiller approximation in order to make the present paper self-contained. We also show our final results for the extended Gutzwiller approximation before going into the details of the derivation. In Sec. III, the original Gutzwiller approximation is extended in order to include longer-range correlations. We divide the whole system into cells and introduce the configurations in each cell for evaluating the weighting factors. Using the general formulation developed in Sec. III, we obtain the renormalization factors for the half-filled case in Sec. IV. We approximate the weighting factors in a perturbative way with respect to the nearest-neighbor correlations. In IV.D, the physical origin of the enhancement of the renormalization factor for the $`z`$-component of exchange interaction is discussed in a viewpoint of statistical weights of real-space spin configurations. In Sec. V, the results at half filling are generalized to the less-than-half-filled case. The present method is applied to the projected variational states in Sec. VI. The self-consistency equations are derived and solved numerically to show that they give reasonable estimate of AF long-range order near half filling. Finally section VII is devoted to a summary and discussions on related problems.
## II $`t`$-$`J`$ model and the original Gutzwiller approximation
We consider the Gutzwiller approximation for the two-dimensional $`t`$-$`J`$ model on a square lattice:
$$\widehat{}=t\underset{ij\sigma }{}P_G(c_{i\sigma }^{}c_{j\sigma }+h.c.)P_G+J\underset{ij}{}𝑺_i𝑺_j,$$
(1)
where $`ij`$ represents the sum over the nearest-neighbor sites and $`𝑺_i=c_{i\alpha }^{}(\frac{1}{2}𝝈)_{\alpha \beta }c_{i\beta }`$. The Gutzwiller’s projection operator $`P_G`$ is defined as $`P_G=_j(1\widehat{n}_j\widehat{n}_j)`$.
For this Hamiltonian, we assume the projected BCS-SDW mean-field wave function
$$|\psi =P_G|\psi _0(\mathrm{\Delta }_d^\mathrm{V},\mathrm{\Delta }_{\mathrm{af}}^\mathrm{V},\mu ^\mathrm{V}),$$
(2)
where $`\mathrm{\Delta }_d^\mathrm{V}`$, $`\mathrm{\Delta }_{\mathrm{af}}^\mathrm{V}`$ and $`\mu ^\mathrm{V}`$ are the variational parameters relating to d-wave SC, AF and chemical potential, respectively, and $`|\psi _0(\mathrm{\Delta }_d^\mathrm{V},\mathrm{\Delta }_{\mathrm{af}}^\mathrm{V},\mu ^\mathrm{V})`$ is a Hartree-Fock type wave function with d-wave SC and AF orders. The wave function $`|\psi `$ is a natural generalization of the RVB state proposed by Anderson. It has been shown that the coexistent state of AF and d-wave SC has the best variational energy in the VMC simulations for the doping rate $`\delta <0.1`$.
In evaluating the variational energy, the projection operator makes difficulties for an analytic approach. The Gutzwiller approximation was developed for this purpose. In this method, the effect of the projection is taken into account by renormalizations of expectation values as follows:
$$c_{i\sigma }^{}c_{j\sigma }=g_tc_{i\sigma }^{}c_{j\sigma }_0,𝑺_i𝑺_j=g_s𝑺_i𝑺_j_0,$$
(3)
where $`\mathrm{}_0`$ represents the expectation value in terms of $`|\psi _0=|\psi _0(\mathrm{\Delta }_d^\mathrm{V},\mathrm{\Delta }_{\mathrm{af}}^\mathrm{V},\mu ^\mathrm{V})`$, and $`\mathrm{}`$ represents the normalized expectation value in $`|\psi =P_G|\psi _0`$;
$`\widehat{𝒪}{\displaystyle \frac{\psi |\widehat{𝒪}|\psi }{\psi |\psi }}={\displaystyle \frac{\psi _0|P_G\widehat{𝒪}P_G|\psi _0}{\psi _0|P_GP_G|\psi _0}}.`$ (4)
The coefficients $`g_t`$ and $`g_s`$ are the renormalization factors of the expectation values, which we call as Gutzwiller factors in the following. Using these notations, the variational energy $`E_{\mathrm{var}}=\widehat{}`$ is rewritten as
$$E_{\mathrm{var}}=\widehat{}_{\mathrm{eff}}_0,$$
(5)
where the parameters $`t`$ and $`J`$ in $`\widehat{}`$ are replaced with
$$t_{\mathrm{eff}}=g_tt,J_{\mathrm{eff}}=g_sJ.$$
(6)
A systematic estimation of $`g_t`$ and $`g_s`$ was developed by Ogawa et al. whose method is briefly reviewed in II.A. For the case with AF order parameter, it can be shown that
$$g_t=\frac{2\delta (1\delta )}{1\delta ^2+4m^2},g_s=\frac{4(1\delta )^2}{(1\delta ^2+4m^2)^2},$$
(7)
where $`\delta `$ is the hole concentration, $`\delta =1n`$, and $`m`$ is the expectation value of AF order parameter in terms of $`|\psi _0`$ defined as
$$m=\frac{(1)^j}{2}\left(\widehat{n}_j_0\widehat{n}_j_0\right).$$
(8)
These Gutzwiller factors, however, do not reproduce the VMC results, as mentioned in I. We find that this discrepancy is due to the fact that only the site-diagonal expectation values are taken into account in obtaining Eq. (7). In this paper we extend the method by Ogawa et al. systematically to include the longer range correlations. We show that it is important to include the effects of the nearest-neighbor expectation values, such as
$`\chi `$ $`=c_{i\sigma }^{}c_{j\sigma }_0,`$ (9)
$`\mathrm{\Delta }`$ $`=c_i^{}c_j^{}_0,`$ (10)
in the estimation of $`g_t`$ and $`g_s`$. This corresponds to taking account of the effect of surroundings of the corresponding bond.
Furthermore we find that the Gutzwiller factors $`g_s`$ for the exchange interaction have different values for the $`z`$ component (denoted as $`g_s^Z`$) and $`xy`$ component ($`g_x^{XY}`$); i.e.,
$`S_i^zS_j^z`$ $`=g_s^ZS_i^zS_j^z_0,`$ (11)
$`S_i^+S_j^{}`$ $`=g_s^{XY}S_i^+S_j^{}_0,`$ (12)
instead of (3) in the presence of $`m`$. Actually from the VMC calculations, we have estimated the behaviors of $`g_s^Z`$ and $`g_s^{XY}`$ as a function of $`m`$, using the relations
$$g_s^Z=\frac{S_i^zS_j^z}{S_i^zS_j^z_0},g_s^{XY}=\frac{S_i^+S_j^{}}{S_i^+S_j^{}_0},$$
(13)
where the numerators $`S_i^zS_j^z`$ and $`S_i^+S_j^{}`$ are evaluated numerically in VMC simulations. It was found that $`g_s^Z`$ is enhanced compared with $`g_s^{XY}`$ and that this enhancement is essential for the stabilization of the AF order. In the present extended Gutzwiller approximation, this result is reproduced. Physically the weighting factor of $`g_s^Z`$ for a specific configuration of a bond is enhanced due to the effect of the surrounding AF correlations.
Before going into details of the derivation, we show our final results. They are summarized as follows:
$`g_s^{XY}=\left({\displaystyle \frac{2(1\delta )}{1\delta ^2+4m^2}}\right)^2a^7,`$ (14)
and
$$g_s^Z=g_s^{XY}\frac{1}{4m^2+X_2}\left[X_2+4m^2\left\{1+\frac{6X_2(1\delta )^2}{1\delta ^2+4m^2}a^3\right\}^2\right],$$
(15)
where
$`a`$ $`=1+{\displaystyle \frac{4X}{(1\delta ^2+4m^2)^2}},`$ (16)
$`X`$ $`=2\delta ^2(\mathrm{\Delta }^2\chi ^2)+8m^2(\chi ^2+\mathrm{\Delta }^2)+4(\chi ^2+\mathrm{\Delta }^2)^2,`$ (17)
$`X_2`$ $`=2(\chi ^2+\mathrm{\Delta }^2).`$ (18)
The above expressions sufficiently reproduce the results obtained in the VMC simulations. However we find that the slight difference can be improved by using
$$X=2\delta ^2(\mathrm{\Delta }^2\chi ^2)+8m^2(\chi ^2+\mathrm{\Delta }^2)+2(\chi ^2+\mathrm{\Delta }^2)^2,$$
(19)
instead of Eq. (18). A typical $`m`$-dependence of $`g_s^Z`$ and $`g_s^{XY}`$ are shown in Fig. 1 for the half-filled case, $`\delta =0`$, and with $`\mathrm{\Delta }`$ being fixed at the two values, $`\mathrm{\Delta }=0.02`$ and $`0.18`$. The behaviors in Fig. 1 agree with those obtained in the VMC simulations using the relation in Eq. (13). The meanings of the rather complicated Gutzwiller factors become apparent in the following sections.
The Gutzwiller factor for the hopping term is given by
$`g_t={\displaystyle \frac{2\delta (1\delta )}{1\delta ^2+4m^2}}{\displaystyle \frac{(1+\delta )^24m^22X_2}{(1+\delta )^24m^2}}a.`$ (20)
Figure 2 shows the $`m`$-dependences of $`g_s^{XY},g_s^Z`$ and $`g_t`$ for $`\delta =0.12`$, which can be compared with the VMC results in Ref. . Note here that the slave-boson mean-field theory simply gives $`g_t=b_i^{}b_j=\delta `$ and $`g_s=1`$. The nonlinear dependences of $`g_s^{XY},g_s^Z`$ and $`g_t`$ on $`\mathrm{\Delta },\chi `$ and $`m`$ are beyond the slave boson theory. In the following subsection we review the original formulation of the Gutzwiller approximation which is useful for the generalization in the later sections.
### A Original formulation of Gutzwiller approximation
Let us start with the formulation of the Gutzwiller approximation developed by Ogawa et al when it is applied to the $`t`$-$`J`$ model. In this subsection we do not consider the AF order. The ultimate aim is to evaluate the variational energy
$$E_{\mathrm{var}}=\widehat{}=\frac{\psi |\widehat{}|\psi }{\psi |\psi },$$
(21)
using $`|\psi =P_G|\psi _0`$. First, Ogawa et al rewrote the denominator as ($`P_G^2=P_G`$)
$`\psi |\psi `$ $`=\psi _0|P_GP_G|\psi _0`$ (25)
$`=\psi _0|{\displaystyle \underset{j}{}}(1\widehat{n}_j\widehat{n}_j)|\psi _0`$
$`=\psi _0|{\displaystyle \underset{j}{}}\{\widehat{n}_j(1\widehat{n}_j)+\widehat{n}_j(1\widehat{n}_j)`$
$`+(1\widehat{n}_j)(1\widehat{n}_j)\}|\psi _0.`$
The expansion of the product in (25) leads to various real-space configurations. Therefore $`\psi |\psi `$ can be rewritten as
$`{\displaystyle \underset{\mathrm{config}.}{}}\psi _0`$ $`|{\displaystyle \underset{j𝒜}{}}\widehat{n}_j(1\widehat{n}_j){\displaystyle \underset{j^{}}{}}\widehat{n}_j^{}(1\widehat{n}_j^{})`$ (27)
$`\times {\displaystyle \underset{j^{\prime \prime }}{}}(1\widehat{n}_{j^{\prime \prime }})(1\widehat{n}_{j^{\prime \prime }})|\psi _0,`$
where $`𝒜`$ ($``$) is a subset of the lattice sites which is singly occupied by an up-spin (down-spin) electron and $``$ is a subset which is vacant. The summation is over all the possible configurations, i.e., all the possible combinations of $`\{𝒜,,\}`$.
In principle, Eq. (27) can be calculated using Wick’s theorem since $`|\psi _0`$ is a mean-field wave function, although it is very complicated. The simplest Gutzwiller approximation is to estimate (27) by using only the site-diagonal expectation values, which gives
$`{\displaystyle \underset{\mathrm{config}.}{}}`$ $`{\displaystyle \underset{j𝒜}{}}\widehat{n}_j(1\widehat{n}_j)_0{\displaystyle \underset{j^{}}{}}\widehat{n}_j^{}(1\widehat{n}_j^{})_0`$ (28)
$`\times `$ $`{\displaystyle \underset{j^{\prime \prime }}{}}(1\widehat{n}_{j^{\prime \prime }})(1\widehat{n}_{j^{\prime \prime }})_0,`$ (29)
with $`\mathrm{}_0`$ meaning the expectation value in $`|\psi _0`$.
For example, the contribution from the site $`j𝒜`$ is evaluated as
$$\omega _A\widehat{n}_j(1\widehat{n}_j)_0=\frac{n}{2}\left(1\frac{n}{2}\right),$$
(30)
where $`n`$ is the average electron density, $`n=N_e/N`$, with $`N`$ ($`N_e`$) being the total number of sites (total electron number). We call $`\omega _A`$ as the weight for $`j`$ site belonging to the subset $`𝒜`$. Since the hole density is $`\delta =1n`$, we obtain
$`\omega _A`$ $`=\omega _B={\displaystyle \frac{n}{2}}\left(1{\displaystyle \frac{n}{2}}\right)={\displaystyle \frac{1\delta ^2}{4}},`$ (31)
$`\omega _E`$ $`=\left(1{\displaystyle \frac{n}{2}}\right)^2={\displaystyle \frac{(1+\delta )^2}{4}}.`$ (32)
Then the expectation value (29) is rewritten as
$$\underset{\mathrm{config}.}{}\omega _A^{N_A}\omega _B^{N_B}\omega _E^{N_E}=\frac{N!}{N_A!N_B!N_E!}\omega _A^{N_A}\omega _B^{N_B}\omega _E^{N_E},$$
(33)
where $`N_A`$ is the number of sites in the subset $`𝒜`$, and so on, ($`N_A+N_B+N_E=N`$). In this simple case, $`N_A`$ is equal to the number of up-spin electrons, so that
$$N_A=N_B=\frac{N_e}{2},N_E=NN_e.$$
(34)
A similar procedure can be carried out for the estimation of the numerator in $`E_{\mathrm{var}}`$. For the exchange term, we have
$`\psi _0`$ $`|P_GS_{\mathrm{}}^+S_m^{}P_G|\psi _0`$ (39)
$`=\psi _0|S_{\mathrm{}}^+S_m^{}{\displaystyle \underset{j\mathrm{},m}{}}\{\widehat{n}_j(1\widehat{n}_j)+\widehat{n}_j(1\widehat{n}_j)`$
$`+(1\widehat{n}_j)(1\widehat{n}_j)\}|\psi _0`$
$`={\displaystyle \underset{\mathrm{config}.}{}}\psi _0|S_{\mathrm{}}^+S_m^{}{\displaystyle \underset{j𝒜^{}}{}}\widehat{n}_j(1\widehat{n}_j){\displaystyle \underset{j^{}^{}}{}}\widehat{n}_j^{}(1\widehat{n}_j^{})`$
$`\times {\displaystyle \underset{j^{\prime \prime }^{}}{}}(1\widehat{n}_{j^{\prime \prime }})(1\widehat{n}_{j^{\prime \prime }})|\psi _0.`$
Here the summation is over all the possible configurations $`\{𝒜^{},^{},^{}\}`$ in which the two sites $`\mathrm{}`$ and $`m`$ are excluded. The evaluation of (39) with the site-diagonal expectation values for $`j\mathrm{},m`$ leads to
$$\frac{(N2)!}{N_A^{}!N_B^{}!N_E^{}!}\omega _A^{N_A^{}}\omega _B^{N_B^{}}\omega _E^{N_E^{}}S_{\mathrm{}}^+S_m^{}_0.$$
(40)
Since the sites $`\mathrm{}`$ and $`m`$ contain one up-spin electron and one down-spin electron, we have $`N_A^{}=\frac{N_e}{2}1`$, $`N_B^{}=\frac{N_e}{2}1`$, and $`N_E^{}=N_E=NN_e`$.
By combining the estimation of the denominator in Eq. (33), we obtain the original Gutzwiller approximation
$`S_{\mathrm{}}^+S_m^{}`$ $`={\displaystyle \frac{\psi _0|P_GS_{\mathrm{}}^+S_m^{}P_G|\psi _0}{\psi _0|P_GP_G|\psi _0}}`$ (43)
$`={\displaystyle \frac{\frac{N_e}{2}\frac{N_e}{2}}{N(N1)}}\omega _A^1\omega _B^1S_{\mathrm{}}^+S_m^{}_0`$
$`={\displaystyle \frac{4}{(1+\delta )^2}}S_{\mathrm{}}^+S_m^{}_0.`$
This gives $`g_s^{XY}=4/(1+\delta )^2`$.
### B Antiferromagnetic case
Ogawa et al also considered the above formalism in the case with AF long-range order for the Hubbard model. The application to the $`t`$-$`J`$ model was carried out by Zhang et al. In this subsection we follow their methods to obtain the Gutzwiller approximation in the $`t`$-$`J`$ model.
In the AF case, the sublattices 1 and 2 are distinguished and their magnetizations are defined as
$$\frac{1}{2}\left(\widehat{n}_j_0\widehat{n}_j_0\right)=m,$$
(44)
for the sublattice 1 and $`m`$ for the sublattice 2. Thus we denote
$`\widehat{n}_j_0`$ $`={\displaystyle \frac{n}{2}}+mr,`$ (45)
$`\widehat{n}_j_0`$ $`={\displaystyle \frac{n}{2}}mw,`$ (46)
for the sublattice 1, where $`r`$ ($`w`$) means the average electron density with the right (wrong) spin direction on the sublattice. $`r`$ and $`w`$ are exchanged for the sublattice 2.
Using these notations, the denominator of $`E_{\mathrm{var}}`$ is evaluated as
$`\psi _0|P_GP_G|\psi _0=`$ $`{\displaystyle \underset{\mathrm{config}.}{}}\omega _{A1}^{N_{A1}}\omega _{B1}^{N_{B1}}\omega _{E1}^{N_{E1}}\omega _{A2}^{N_{A2}}\omega _{B2}^{N_{B2}}\omega _{E2}^{N_{E2}}`$ (47)
$`=`$ $`{\displaystyle \underset{N_{A1},N_{B1}}{}}{\displaystyle \frac{\left(\frac{N}{2}\right)!}{N_{A1}!N_{B1}!N_{E1}!}}{\displaystyle \frac{\left(\frac{N}{2}\right)!}{N_{A2}!N_{B2}!N_{E2}!}}`$ (49)
$`\times \omega _{A1}^{N_{A1}}\omega _{B1}^{N_{B1}}\omega _{E1}^{N_{E1}}\omega _{A2}^{N_{A2}}\omega _{B2}^{N_{B2}}\omega _{E2}^{N_{E2}},`$
where the configurations in the sublattice 1 are specified by $`(N_{A1},N_{B1},N_{E1})`$, and so on. Contrary to the previous subsection, Eq. (49) has a summation over possible values of $`N_{A1}`$ and $`N_{B1}`$ because there are only four constraints between six numbers, $`(N_{A1},N_{B1},N_{E1},N_{A2},N_{B2},N_{E2})`$:
$`N_{A1}+N_{B1}+N_{E1}`$ $`={\displaystyle \frac{N}{2}},N_{A2}+N_{B2}+N_{E2}={\displaystyle \frac{N}{2}},`$ (50)
$`N_{A1}+N_{A2}`$ $`={\displaystyle \frac{N_e}{2}},N_{B1}+N_{B2}={\displaystyle \frac{N_e}{2}}.`$ (51)
The weights in the site-diagonal expectation values are
$`\omega _{A1}`$ $`=\widehat{n}_j_01\widehat{n}_j_0=r(1w),`$ (53)
$`\omega _{B1}`$ $`=\widehat{n}_j_01\widehat{n}_j_0=w(1r),`$ (54)
$`\omega _{E1}`$ $`=1\widehat{n}_j_01\widehat{n}_j_0=(1r)(1w),`$ (55)
for the sublattice 1, and
$`\omega _{A2}`$ $`=\omega _{B1},`$ (56)
$`\omega _{B2}`$ $`=\omega _{A1},`$ (57)
$`\omega _{E2}`$ $`=\omega _{E1},`$ (58)
for the sublattice 2.
Since $`N_{A1},N_{B1}`$ and $`N`$ are large numbers, Eq. (49) is further approximated by the largest term in the summation. By taking the partial derivative of the logarithm of each term in (49) with respect to $`N_{A1}`$, we obtain
$$\mathrm{ln}\overline{N_{A1}}+\mathrm{ln}(\frac{N_e}{2}\overline{N_{A1}})+\mathrm{ln}\omega _{A1}\mathrm{ln}\omega _{A2}=0,$$
(59)
with $`\overline{N_{A1}}`$ being the value of $`N_{A1}`$ which gives the largest term in the summation. By solving (59), $`\overline{N_{A1}}`$ is given by
$$\overline{N_{A1}}=\frac{r(1w)}{r(1w)+w(1r)}\frac{N_e}{2}=\frac{r(1w)}{n2rw}\frac{N_e}{2}.$$
(60)
Similarly we obtain
$`\overline{N_{B1}}={\displaystyle \frac{w(1r)}{n2rw}}{\displaystyle \frac{N_e}{2}}.`$ (61)
The summation in (49) is approximated by the term with $`\overline{N_{A1}}`$ and $`\overline{N_{B1}}`$.
Again the numerator such as $`\psi _0|P_GS_{\mathrm{}}^+S_m^{}P_G|\psi _0`$ is approximated in the similar way. As a result, the Gutzwiller approximation in the presence of AF order, $`m`$, becomes
$`S_{\mathrm{}}^+S_m^{}`$ $`={\displaystyle \frac{\psi _0|P_GS_{\mathrm{}}^+S_m^{}P_G|\psi _0}{\psi _0|P_GP_G|\psi _0}}`$ (65)
$`={\displaystyle \frac{\overline{N_{A1}}\overline{N_{B1}}}{\frac{N}{2}\frac{N}{2}}}\omega _{A1}^1\omega _{B1}^1S_{\mathrm{}}^+S_m^{}_0`$
$`={\displaystyle \frac{n^2}{(n2rw)^2}}S_{\mathrm{}}^+S_m^{}_0`$
$`={\displaystyle \frac{4(1\delta )^2}{(1\delta ^2+4m^2)^2}}S_{\mathrm{}}^+S_m^{}_0,`$
where a useful relation
$$\frac{n}{n2rw}=\frac{2(1\delta )}{1\delta ^2+4m^2}$$
(66)
has been used. Equation (65) gives
$$g_s^{XY}=\frac{4(1\delta )^2}{(1\delta ^2+4m^2)^2},$$
(67)
which is simply Eq. (7).
It is straightforward to see
$$g_s^Z=g_s^{XY},$$
(68)
and
$$g_t=\frac{2\delta (1\delta )}{1\delta ^2+4m^2}.$$
(69)
## III Extension of the Gutzwiller approximation: Formulation
The approximation in the preceding section is restricted to the site-diagonal (on-site) expectation values. It was shown that this approximation is not good enough to reproduce the AF order obtained in the VMC simulation. A straightforward extension of the previous Gutzwiller approximation is to consider longer-range expectation values. The nearest-neighbor expectation was considered before. In this section, we develop a general formalism to take account of the longer-range effects systematically.
In evaluating the denominator of $`E_{\mathrm{var}}`$, i.e., $`\psi _0|P_GP_G|\psi _0`$, various configurations have been classified depending on the state ($`,`$ or a hole) on each site as $`\{𝒜,,\}`$ in the previous section. Instead of this we divide the whole system into cells consisting of $`N_c`$ sites. (Later we take $`N_c`$ as a large enough number.) For each cell, there are $`3^{N_c}`$ configurations because each site has $`,`$ or a hole. We denote these configurations as states of a cell. Using these states, the denominator of $`E_{\mathrm{var}}`$ can be approximated as
$`\psi _0|P_GP_G|\psi _0`$ $`={\displaystyle \underset{\mathrm{all}\mathrm{the}\mathrm{possible}\mathrm{states}}{}}{\displaystyle \underset{i=1}{\overset{K}{}}}\omega _i^{N_i}`$ (71)
$`={\displaystyle \underset{\{N_i\}}{}}{\displaystyle \frac{\left(\frac{N}{N_c}\right)!}{_{i=1}^K(N_i)!}}{\displaystyle \underset{i=1}{\overset{K}{}}}\omega _i^{N_i},`$
with $`K=3^{N_c}`$. Here $`N_i`$ is the number of cells in the $`i`$-th state, and $`\omega _i`$ is the weight of the $`i`$-th state cell in the analogy with Eq. (49). Explicitly the weight $`\omega _i`$ is defined as
$`\omega _i=`$ $`\psi _0|{\displaystyle \underset{j}{}}\widehat{n}_j(1\widehat{n}_j){\displaystyle \underset{j^{}}{}}\widehat{n}_j^{}(1\widehat{n}_j^{})`$ (73)
$`\times {\displaystyle \underset{j^{\prime \prime }}{}}(1\widehat{n}_{j^{\prime \prime }})(1\widehat{n}_{j^{\prime \prime }})|\psi _0,`$
with $`j,j^{}`$ and $`j^{\prime \prime }`$ being the sites in the cell. To calculate $`\omega _i`$ is the most important part of the present theory which we carry out later.
The last summation in (71) is over all the possible values of $`N_i`$ under the following constraints:
$`{\displaystyle \underset{i=1}{\overset{K}{}}}n_iN_i={\displaystyle \frac{N_e}{2}},{\displaystyle \underset{i=1}{\overset{K}{}}}n_iN_i={\displaystyle \frac{N_e}{2}},`$ (74)
$`{\displaystyle \underset{i=1}{\overset{K}{}}}n_{\mathrm{h}i}N_i=NN_e,`$ (75)
where $`n_i`$ ($`n_i`$) represents the number of up-spins (down-spins) in the $`i`$-th state cell, and $`n_{\mathrm{h}i}`$ represents the number of holes, respectively. Since a cell has $`N_c`$ sites, a relation
$`n_i+n_i+n_{\mathrm{h}i}=N_c,`$ (76)
holds.
Under the above constraints, we look for the largest term in Eq. (71) as was done in the original Gutzwiller approximation (see Eq. (59)). Since the constraints (75) are rather complicated, it is useful to introduce the Lagrange multipliers ($`\mu _{},\mu _{},\lambda `$). Taking the partial derivative of
$`\mathrm{ln}\left[{\displaystyle \frac{\left(\frac{N}{N_c}\right)!}{_{i=1}^K(N_i)!}}{\displaystyle \underset{i=1}{\overset{K}{}}}\omega _i^{N_i}\right]\mu _{}\left({\displaystyle \underset{i=1}{\overset{K}{}}}n_iN_i{\displaystyle \frac{N_e}{2}}\right)`$ (77)
$`\mu _{}\left({\displaystyle \underset{i=1}{\overset{K}{}}}n_iN_i{\displaystyle \frac{N_e}{2}}\right)\lambda \left({\displaystyle \underset{i=1}{\overset{K}{}}}n_{\mathrm{h}i}N_iN+N_e\right),`$ (78)
we obtain
$`\overline{N_i}=\omega _i\mathrm{exp}(\mu _{}n_i\mu _{}n_i\lambda n_{\mathrm{h}i}).`$ (79)
From symmetry we find $`\mu _{}=\mu _{}=\mu `$, and using the relation (76) we get
$`\overline{N_i}=\omega _i\mathrm{e}^{\mu N_c}\mathrm{e}^{(\mu \lambda )n_{\mathrm{h}i}}.`$ (80)
Furthermore we introduce new variables $`W`$ and $`p`$ as
$`\mathrm{e}^{\mu N_c}`$ $`{\displaystyle \frac{N}{N_c}}{\displaystyle \frac{1}{W}},`$ (81)
$`\mathrm{e}^{\mu \lambda }`$ $`p.`$ (82)
Using these variables, $`\overline{N_i}`$ is rewritten as
$`\overline{N_i}={\displaystyle \frac{N}{N_c}}{\displaystyle \frac{\omega _i}{W}}p^{n_{\mathrm{h}i}}.`$ (83)
The variables $`W`$ and $`p`$ are to be determined from the constraints (75) which are equivalent to
$`{\displaystyle \underset{i=1}{\overset{K}{}}}\overline{N_i}={\displaystyle \frac{N}{N_c}},{\displaystyle \underset{i=1}{\overset{K}{}}}n_{\mathrm{h}i}\overline{N_i}=NN_e.`$ (84)
These conditions become
$`{\displaystyle \underset{i=1}{\overset{K}{}}}{\displaystyle \frac{\omega _i}{W}}p^{n_{\mathrm{h}i}}`$ $`=1,`$ (85)
$`{\displaystyle \underset{i=1}{\overset{K}{}}}n_{\mathrm{h}i}{\displaystyle \frac{\omega _i}{W}}p^{n_{\mathrm{h}i}}`$ $`=\delta N_c.`$ (86)
It is convenient to classify all the possible states of a cell into subgroups which contain $`j`$ holes ($`j=0,1,2,\mathrm{},N_c`$). (We call these subgroups as $`j`$-hole sectors.) Then the quantity
$`W_j={\displaystyle \underset{i\mathrm{with}j\mathrm{holes}}{}}\omega _i,`$ (87)
represents the total weight of the states in this subgroup. Using $`W_j`$, we can rewrite the constraints as
$`{\displaystyle \underset{j=0}{\overset{N_c}{}}}{\displaystyle \frac{W_j}{W}}p^j`$ $`=1,`$ (88)
$`{\displaystyle \underset{j=0}{\overset{N_c}{}}}j{\displaystyle \frac{W_j}{W}}p^j`$ $`=\delta N_c.`$ (89)
The first equality is rewritten as $`W=_jW_jp^j`$ so that $`W`$ represents the total weight. In the following sections, we evaluate $`W_j`$ and then determine $`W`$ and $`p`$ from Eq. (89).
The numerator in $`E_{\mathrm{var}}`$ is evaluated in a similar way. Let us consider the expectation value of an operator $`\widehat{𝒪}`$ (such as $`S_{\mathrm{}}^+S_m^{}`$) which operates on a certain part inside a cell of the $`i_0`$-th state. (We call this cell as the central cell.) Then all the configurations in the other cells are classified by $`\{N_i^{}\}`$ which is the number of cells of the $`i`$-th state. Then the expectation value becomes
$$\psi _0|P_G\widehat{𝒪}P_G|\psi _0=\underset{i_0}{}\underset{\{N_i^{}\}}{}\frac{\left(\frac{N}{N_c}1\right)!}{_{i=1}^K(N_i^{})!}\underset{i=1}{\overset{K}{}}\omega _i^{N_i^{}}\widehat{𝒪}_{i_0},$$
(90)
in the analogy with Eq. (40). Here $`\widehat{𝒪}_{i_0}`$ indicates the expectation value of $`\widehat{𝒪}`$ together with terms such as $`_i\widehat{n}_i(1\widehat{n}_i)`$ inside the central cell of the $`i_0`$-th state.
Determining the largest terms on the right-hand side and taking the ratio to the denominator, we finally obtain
$`\widehat{𝒪}{\displaystyle \frac{\psi _0|P_G\widehat{𝒪}P_G|\psi _0}{\psi _0|P_GP_G|\psi _0}}`$ $`={\displaystyle \underset{i_0}{}}{\displaystyle \frac{1}{\left(\frac{N}{N_c}\right)}}N_{i_0}\omega _{i_0}^1\widehat{𝒪}_{i_0}`$ (92)
$`={\displaystyle \underset{i_0}{}}{\displaystyle \frac{p^{n_{\mathrm{h}i_0}}}{W}}\widehat{𝒪}_{i_0}.`$
The evaluation of the largest term and the derivation of the above result is shown in Appendix A.
If we consider only the site-diagonal expectation values to estimate $`\omega _i`$, Eq. (92) reproduces the original Gutzwiller approximation shown in the previous section. This is summarized in Appendix B.
## IV Half-Filled Case
To obtain the Gutzwiller factors formulated in the last section, it is necessary to evaluate the weights $`\omega _i`$ for each state. Although these weights are generally complicated, we estimate them assuming that the corrections to the original Gutzwiller approximation (reproduced in Appendix B) are small. To be more specific, we evaluate $`W_j`$ in the lowest orders with respect to the intersite correlations $`\chi `$ and $`\mathrm{\Delta }`$ defined in Eq. (10). Typical values for $`\chi `$ and $`\mathrm{\Delta }`$ are less than $`0.2`$, so that the perturbation with respect to them will be justified.
In this section we calculate the Gutzwiller approximation for the half filling, since it is less complicated than the doped case. All the cells do not contain holes so that only the weight, $`W_0`$, is to be calculated. We call this subset of states as 0-hole sector. After the half-filled case, it is rather straightforward to extend the results to the doped case.
### A Evaluation of $`W_0`$
First we calculate the weight $`\omega _i`$ of the $`i`$-th state which has $`n_i`$ up-spin electrons, $`n_i`$ down-spin electrons:
$$\omega _i=\underset{j𝒜}{}\widehat{n}_j(1\widehat{n}_j)\underset{j^{}}{}\widehat{n}_j^{}(1\widehat{n}_j^{})_0.$$
(93)
In the lowest order of $`\chi `$ and $`\mathrm{\Delta }`$ or in the zeroth order, we have to take only the site-diagonal expectation values. This gives
$`\omega _i^{(0)}=[r(1w)]^{n_{\mathrm{right}}}[w(1r)]^{n_{\mathrm{wrong}}},`$ (94)
where $`n_{\mathrm{right}}`$ ($`n_{\mathrm{wrong}}`$) means the number of sites on which the right (wrong) spices of spin direction is located depending on sublattice 1 and 2. The summation over all the possible configurations gives
$`W_0^{(0)}=[r(1w)+w(1r)]^{N_c}=(n2rw)^{N_c}.`$ (95)
In the next order with respect to $`\chi `$ and $`\mathrm{\Delta }`$, we have to consider contributions from expectation values of bonds in the cell as shown in Fig. 3. For example, consider a bond connecting the sites $`i`$ and $`j`$ (see Fig. 3(a)) and a state in which two up-spin electrons occupy both $`i`$ and $`j`$ sites. The expectation value in this bond gives
$`P_{}`$ $`=(1\widehat{n}_i)\widehat{n}_i\widehat{n}_j(1\widehat{n}_j)_0`$ (98)
$`=rw(1r)(1w)rw\chi ^2(1r)(1w)\chi ^2`$
$`r(1r)\mathrm{\Delta }^2w(1w)\mathrm{\Delta }^2+(\chi ^2+\mathrm{\Delta }^2)^2,`$
where we have used the Wick’s theorem. Generally $`\chi `$ and $`\mathrm{\Delta }`$ can be complex numbers and $`\chi ^2,\mathrm{\Delta }^2`$ in the above expression mean $`|\chi |^2,|\mathrm{\Delta }|^2`$ implicitly. In the similar way we calculate $`P_{},P_{},P_{}`$. Writing that the left-spin in the subscript indicates the spin on the sublattice 1, they are calculated as
$`P_{}`$ $`=P_{},`$ (99)
$`P_{}`$ $`=r^2(1w)^2+2r(1w)\chi ^2`$ (101)
$`+r^2\mathrm{\Delta }^2+(1w)^2\mathrm{\Delta }^2+(\chi ^2+\mathrm{\Delta }^2)^2,`$
$`P_{}`$ $`=w^2(1r)^2+2w(1r)\chi ^2`$ (103)
$`+w^2\mathrm{\Delta }^2+(1r)^2\mathrm{\Delta }^2+(\chi ^2+\mathrm{\Delta }^2)^2.`$
In the estimation of $`W_0^{(1)}`$ we need their summation
$`P_{}+P_{}+P_{}+P_{}=(n2rw)^2+X,`$ (104)
with $`X`$ defined in Eq. (18), where we have used the relations $`r=n/2+m`$ and $`w=n/2m`$. The first term $`(n2rw)^2`$ on the r.h.s. is the zeroth order contribution which has been included in $`W_0^{(0)}`$. In this sense, $`X`$ is a kind of connected (or cummulant) contribution coming from the real-space diagram in Fig. 3(a). The actual values of $`X`$ is roughly $`1/20`$ so that our perturbation scheme can be justified. Denoting the number of bonds in a cell as $`N_b`$, we have a contribution to $`W_0`$ as
$`W_0^{(1)}=N_bX(n2rw)^{N_c2},`$ (105)
where the factor $`(n2rw)^{N_c2}`$ comes from the contributions of the other $`N_c2`$ sites in the cell except for $`i`$ and $`j`$.
In the similar way, we evaluate the higher order contributions to $`W_0`$. By considering two bonds in the cell we approximate their contribution as
$`W_0^{(2)}=_{N_b}C_2X^2(n2rw)^{N_c4},`$ (106)
where $`{}_{N_b}{}^{}C_{2}^{}=N_b!/(N_b2)!2!`$ is the number of choices for the positions of the two bonds in a cell. The summation of these series leads to
$`W_0`$ $`={\displaystyle \underset{j=0}{\overset{N_b}{}}}{}_{N_b}{}^{}C_{j}^{}X^j(n2rw)^{N_c2j}`$ (108)
$`=(n2rw)^{N_c}\left(1+{\displaystyle \frac{X}{(n2rw)^2}}\right)^{N_b}.`$
Let us discuss the neglected contributions to $`W_0`$. Figures 3(b) and 3(c) show real-space diagrams whose connected expectation values can contribute to $`W_0`$. Although their contributions are not included, it can be shown that they are smaller than the terms in (108). For example, the contribution from Fig. 3(b) is in the order of $`N_b\delta ^2\chi ^4`$ and $`N_bm^2\chi ^4`$ etc., so that it is smaller than $`W_0^{(1)}`$. The contribution from Fig. 3(c) is in the order of $`X^2N_b`$ which is smaller than (106) by a factor $`1/N_b`$.
There is another effect neglected in $`W_0`$: The higher order terms in (108) become less and less correct because the number of the bonds giving $`X`$ is not exactly $`{}_{N_b}{}^{}C_{j}^{}`$ due to their overlapping. However the error is again in the order smaller than $`{}_{N_b}{}^{}C_{j}^{}`$. In this sense, Eq. (108) is a kind of a summation of most dominant terms in the order of $`X^jN_b^j`$.
### B Evaluation of $`g_s^{XY}`$
In order to obtain the Gutzwiller factor for $`S_{\mathrm{}}^+S_m^{}`$ at half filling we need to calculate
$`S_{\mathrm{}}^+S_m^{}={\displaystyle \underset{i_0}{}}{\displaystyle \frac{1}{W_0}}S_{\mathrm{}}^+S_m^{}_{i_0},`$ (109)
according to the general formulation (92). Note that $`W=W_0`$ and the central cell (which includes sites $`\mathrm{}`$ and $`m`$) does not contain holes, i.e., $`n_{\mathrm{h}i_0}=0`$. The average in a central cell of the $`i_0`$-th state, $`S_{\mathrm{}}^+S_m^{}_{i_0}`$, is evaluated in a similar perturbation scheme as in $`W_0`$. In the lowest order of $`\chi `$ and $`\mathrm{\Delta }`$, we have (Fig. 4(a))
$`S_{\mathrm{}}^+S_m^{}_0(n2rw)^{N_c2},`$ (110)
where the factor $`(n2rw)^{N_c2}`$ comes from the contributions from the sites in the central cell other than the sites $`\mathrm{}`$ and $`m`$.
In the next order of $`\chi `$ and $`\mathrm{\Delta }`$, there are contributions of $`X`$ from the bonds in the cell. This leads to
$`\stackrel{~}{N}_bXS_{\mathrm{}}^+S_m^{}_0(n2rw)^{N_c4},`$ (111)
where $`\stackrel{~}{N}_b`$ is the number of bonds which are not connected to the sites $`\mathrm{}`$ and $`m`$ directly (Fig. 4(b)). In general $`\stackrel{~}{N}_b`$ depends on the shape of the cell. However, if we choose the large enough cell, we have
$`\stackrel{~}{N}_b=N_b7,`$ (112)
as is evident from Fig. 4(b).
Although the dotted bonds in Fig. 4(b) do not give $`X`$, they may still give different contributions to $`S_{\mathrm{}}^+S_m^{}_{i_0}`$, such as
$`S_{\mathrm{}}^+S_m^{}\widehat{n}_m^{}(1\widehat{n}_m^{})_c,`$ (113)
where $`m^{}`$ is a nearest neighbor site of $`m`$ and $`\mathrm{}_c`$ means a connected expectation value excluding the disconnected terms such as $`S_{\mathrm{}}^+S_m^{}_0\widehat{n}_m^{}(1\widehat{n}_m^{})_0.`$ (The disconnected expectation values have been already taken into account in Eq. (110).) By calculating (113), however, we find that there are no connected conributions: Actually the two electron operators in $`S_{\mathrm{}}^+=c_{\mathrm{}}^{}c_{\mathrm{}}`$ have to make contractions with the electron operators on the site $`m`$ in the Wick’s expansion, which gives just the disconnected expectation value.
By summing up the series of corrections such as (111) we obtain
$`{\displaystyle \underset{i_0}{}}S_{\mathrm{}}^+S_m^{}_{i_0}`$ (114)
$`=S_{\mathrm{}}^+S_m^{}_0\{(n2rw)^{N_c2}`$ (115)
$`+\stackrel{~}{N}_bX(n2rw)^{N_c4}+\mathrm{}\}`$ (116)
$`=S_{\mathrm{}}^+S_m^{}_0(n2rw)^{N_c2}\left(1+{\displaystyle \frac{X}{(n2rw)^2}}\right)^{\stackrel{~}{N}_b}.`$ (117)
Combining with $`W_0`$, the final result for the extended Gutzwiller approximation for $`S_{\mathrm{}}^+S_m^{}`$ becomes
$`S_{\mathrm{}}^+S_m^{}={\displaystyle \frac{S_{\mathrm{}}^+S_m^{}_0}{(n2rw)^2}}\left(1+{\displaystyle \frac{X}{(n2rw)^2}}\right)^{(N_b\stackrel{~}{N}_b)}.`$ (118)
Since the Gutzwiller factor is defined as the ratio between $`S_{\mathrm{}}^+S_m^{}`$ and $`S_{\mathrm{}}^+S_m^{}_0`$, we have
$`g_s^{XY}`$ $`={\displaystyle \frac{1}{(n2rw)^2}}\left(1+{\displaystyle \frac{X}{(n2rw)^2}}\right)^{(N_b\stackrel{~}{N}_b)}`$ (120)
$`={\displaystyle \frac{a^{(N_b\stackrel{~}{N}_b)}}{(n2rw)^2}},`$
where we have put
$`a1+{\displaystyle \frac{X}{(n2rw)^2}},`$ (121)
which appears frequently in the following discussion. If we choose a large enough cell like $`N_c>>2`$, we obtain
$`g_s^{XY}={\displaystyle \frac{a^7}{(n2rw)^2}},`$ (122)
according to Eq. (112). Physically the factor $`a^7`$ with $`a>1`$ represents the exclusion effect of the bonds as shown in Fig. 4(b).
### C Evaluation of $`g_s^Z`$
The Gutzwiller factor for the $`z`$-component of $`𝑺_{\mathrm{}}𝑺_m`$ is calculated in the similar way. The lowest orders with respect to $`\chi `$ and $`\mathrm{\Delta }`$ give the similar results to (110) and (111), in which $`S_{\mathrm{}}^+S_m^{}_0`$ is replaced with $`S_{\mathrm{}}^zS_m^z_0`$. However, for $`S_{\mathrm{}}^zS_m^z_{i_0}`$, an additional contribution appears from the diagram in Fig. 5(a) which did not give contributions in $`S_{\mathrm{}}^+S_m^{}_{i_0}`$. This is the reason why the anisotropy for the Gutzwiller factor between $`xy`$ component and $`z`$ component appears, which is the most important feature of our extended Gutzwiller approximation.
Assuming that the site $`\mathrm{}`$ and $`m^{}`$ in Fig. 5(a) are in the sublattice 1, we have a contribution
$`S_{\mathrm{}}^zS_m^z\left\{\widehat{n}_m^{}(1\widehat{n}_m^{})+\widehat{n}_m^{}(1\widehat{n}_m^{})\right\}_c`$ (123)
$`=S_{\mathrm{}}^z_0S_m^z\widehat{n}_m^{}(1\widehat{n}_m^{})+S_m^z\widehat{n}_m^{}(1\widehat{n}_m^{})_c`$ (124)
$`={\displaystyle \frac{m}{2}}(P_{}P_{}+P_{}P_{})`$ (125)
$`=m^2X_2,`$ (126)
with $`X_2`$ defined in Eq. (18). Since the number of bonds connected to the site $`m`$ is
$`N_2{\displaystyle \frac{N_b\stackrel{~}{N}_b1}{2}},`$ (127)
the contribution to $`_{i_0}S_{\mathrm{}}^zS_m^z_{i_0}`$ turns out to be
$`2N_2m^2X_2(n2rw)^{N_c3},`$ (128)
where the contributions from the bonds connected to the site $`\mathrm{}`$ are also included.
Similarly the diagram in Fig. 5(b) gives a contribution
$`S_{\mathrm{}}^z\widehat{n}_{\mathrm{}^{}}(1\widehat{n}_{\mathrm{}^{}})+S_{\mathrm{}}^z\widehat{n}_{\mathrm{}^{}}(1\widehat{n}_{\mathrm{}^{}})_c`$ (129)
$`\times S_m^z\widehat{n}_m^{}(1\widehat{n}_m^{})+S_m^z\widehat{n}_m^{}(1\widehat{n}_m^{})_c`$ (130)
$`={\displaystyle \frac{1}{4}}(P_{}P_{}+P_{}P_{})^2`$ (131)
$`=m^2X_2^2.`$ (132)
By counting the number of possible combination of the bonds, we have
$`N_2^2m^2X_2^2(n2rw)^{N_c4}.`$ (133)
The higher order terms of $`X^j`$ are calculated as before, which become
$`S_{\mathrm{}}^zS_m^z={\displaystyle \frac{1}{W_0}}\{`$ $`S_{\mathrm{}}^zS_m^z_0(n2rw)^{N_c2}a^{\stackrel{~}{N}_b}`$ (134)
$``$ $`2N_2m^2X_2(n2rw)^{N_c3}a^{\stackrel{~}{N}_{b}^{}{}_{}{}^{}}`$ (135)
$``$ $`N_2^2m^2X_2^2(n2rw)^{N_c4}a^{\stackrel{~}{N}_{b}^{}{}_{}{}^{\prime \prime }}\},`$ (136)
with $`\stackrel{~}{N}_{b}^{}{}_{}{}^{}`$ and $`\stackrel{~}{N}_{b}^{}{}_{}{}^{\prime \prime }`$ being the numbers of the bonds which are not connected directly to the diagrams in Figs. 5(a) and 5(b), respectively. They are shown in Figs. 6(a) and (b). In the large enough cell, we have
$`\stackrel{~}{N}_{b}^{}{}_{}{}^{}`$ $`=N_b10,`$ (137)
$`\stackrel{~}{N}_{b}^{}{}_{}{}^{\prime \prime }`$ $`=N_b13.`$ (138)
Combining with $`W_0`$, we obtain
$`S_{\mathrm{}}^zS_m^z=`$ $`{\displaystyle \frac{a^{(N_b\stackrel{~}{N}_b)}}{(n2rw)^2}}\{S_{\mathrm{}}^zS_m^z_0`$ (142)
$`{\displaystyle \frac{2N_2}{n2rw}}m^2X_2a^{(\stackrel{~}{N}_b\stackrel{~}{N}_{b}^{}{}_{}{}^{})}`$
$`{\displaystyle \frac{N_2^2}{(n2rw)^2}}m^2X_2^2a^{(\stackrel{~}{N}_b\stackrel{~}{N}_{b}^{}{}_{}{}^{\prime \prime })}\}.`$
In this case, $`S_{\mathrm{}}^zS_m^z`$ is not proportional to $`S_{\mathrm{}}^zS_m^z_0`$. Since the Gutzwiller factor is the ratio between the two, we have some nontrivial contribution from the second and third terms on the r.h.s of (142). Using
$`S_{\mathrm{}}^zS_m^z_0=m^2{\displaystyle \frac{X_2}{4}},`$ (143)
the final expression for the Gutzwiller factor turns out to be
$`g_s^Z`$ $`=g_s^{XY}{\displaystyle \frac{1}{4m^2+X_2}}`$ (145)
$`\times \left[X_2+4m^2\left\{1+{\displaystyle \frac{N_2X_2}{n2rw}}a^{(\stackrel{~}{N}_b\stackrel{~}{N}_{b}^{}{}_{}{}^{})}\right\}^2\right],`$
where we have assumed
$`\stackrel{~}{N}_b\stackrel{~}{N}_{b}^{}{}_{}{}^{\prime \prime }=2(\stackrel{~}{N}_b\stackrel{~}{N}_{b}^{}{}_{}{}^{}).`$ (146)
For a large enough cell we have $`N_2=3`$ and $`\stackrel{~}{N}_b\stackrel{~}{N}_{b}^{}{}_{}{}^{}=3`$, so that
$$g_s^Z=g_s^{XY}\frac{1}{4m^2+X_2}\left[X_2+4m^2\left\{1+\frac{3X_2}{n2rw}a^3\right\}^2\right].$$
(147)
The above formula is one of the most important results in this paper. Typical $`m`$-dependences have been shown in Fig. 1. Let us here check some limiting cases. When $`m=0`$, we have
$$g_s^{XY}=g_s^Z=4a^7,$$
(148)
which reproduces the original Gutzwiller approximation. For small $`m`$, we obtain
$$g_s^Z=g_s^{XY}\left\{1+4m^2\left(\frac{6}{n2rw}a^3+\frac{9X_2}{(n2rw)^2}a^6\right)\right\}.$$
(149)
It is apparent that $`g_s^Z`$ has an enhancement as a function of $`m`$, which gives the reasonable AF long-range order at half filling.
### D Physical meaning of the enhancement of $`g_s^Z`$
In the usual interpretation of the Gutzwiller approximation, a comparison is made between the probabilities of spin configurations in the wave functions with and without the projection operators. When we consider only two sites $`\mathrm{}`$ and $`m`$ for $`𝑺_{\mathrm{}}𝑺_m`$, the summation of the probabilities of spin configurations $`|,|,|`$ and $`|`$ in the wave function without the projection is given by $`P_{}+P_{}+P_{}+P_{}`$ which was calculated as $`(n2rw)^2+X`$ in Eq. (104). On the other hand, in the wave function with the projection, we have always one of the above four states at half filling, and thus the probability is equal to 1. The ratio of these probabilities gives the Gutzwiller factor
$`g_s^{XY}=g_s^Z={\displaystyle \frac{1}{(n2rw)^2+X}}.`$ (150)
However this Gutzwiller approximation does not give a reasonable answer compared with the VMC results as mentioned before.
For the physical understanding of the results in the previous subsections, it is necessary to go beyond the two sites, $`\mathrm{}`$ and $`m`$, because the enhancement of $`g_s^Z`$ comes from the diagrams in Fig. 5. For this purpose we consider the spin configurations on three sites as shown in Fig. 7. To calculate the probabilities of these configurations is straightforward. For example for Fig. 7(c) we have
$`\left(\begin{array}{ccc}& & \\ & & \end{array}\right)=2r(1w)P_{}2r^3(1w)^3,`$ (151)
where the second term on the r.h.s. appears to avoid the double counting of the zeroth order with respect to $`\chi `$ and $`\mathrm{\Delta }`$. We can see that the presence of the third site $`m^{}`$ enhances the probability of this configuration, $`\left(\begin{array}{ccc}& & \\ & & \end{array}\right)`$, because $`P_{}>P_{}=P_{}>P_{}`$ in the presence of AF order. As a result, the weight of the configuration $``$ on $`\mathrm{}`$ and $`m`$ sites is increased in the wave function. This causes the enhancement of $`g_s^Z`$.
To understand this effect more quantitatively, we calculate $`S_{\mathrm{}}^zS_m^z`$ directly from the configurations in Fig. 7 as follows:
$`4S_{\mathrm{}}^zS_m^z`$ (152)
$`={\displaystyle \frac{(\mathrm{a})+(\mathrm{b})(\mathrm{c})(\mathrm{d})(\mathrm{e})(\mathrm{f})+(\mathrm{g})+(\mathrm{h})}{(\mathrm{a})+(\mathrm{b})+(\mathrm{c})+(\mathrm{d})+(\mathrm{e})+(\mathrm{f})+(\mathrm{g})+(\mathrm{h})}}.`$ (153)
The denominator becomes
$`(n2rw)\{(n2rw)^2+2X\},`$ (154)
while the numerator becomes
$`(n2rw)(P_{}P_{}P_{}+P_{})`$ (155)
$`+2m(P_{}P_{}+P_{}P_{})+4m^2(n2rw)`$ (156)
$`=4(n2rw)S_{\mathrm{}}^zS_m^z_04m^2X_2.`$ (157)
Although the denominator is larger than the two-site case in Eq. (150), the second term in the numerator enhances $`S_{\mathrm{}}^zS_m^z`$, which is mainly from the contribution of the configuration 7(c). Using $`4S_{\mathrm{}}^zS_m^z_0=(4m^2+X_2)`$ and taking the ratio, we obtain
$`{\displaystyle \frac{S_{\mathrm{}}^zS_m^z}{S_{\mathrm{}}^zS_m^z_0}}`$ $`={\displaystyle \frac{1}{(n2rw)^2+2X}}{\displaystyle \frac{1}{4m^2+X_2}}`$ (159)
$`\times \left\{X_2+4m^2(1+{\displaystyle \frac{X_2}{n2rw}})\right\}.`$
This is essentially the result obtained in the previous subsection for $`g_s^Z`$. Since we have considered only the site $`m^{}`$ as the third term, $`N_2`$ in Eq. (145) is replaced with $`1`$ in this case. The first factor $`1/\{(n2rw)^2+2X\}`$ corresponds the exclusion effect which was represented by $`a^{(N_b\stackrel{~}{N}_b)}`$ in $`g_s^{XY}`$. From this analysis it is apparent that the enhancement of $`g_s^Z`$ as a function of $`m`$ is due to the increase of the probability of $``$ in the presence of the AF circumstances.
## V Less-than-half-filled case
### A Evaluation of the total weight $`W`$
In the presence of holes, the Gutzwiller factors are calculated similarly as in the half-filled case. In this subsection we obtain the total weight, $`W`$. Since the weight for the zero-hole sector $`W_0`$ has been calculated, we have to evaluate
$`W_1={\displaystyle \underset{i\mathrm{with}1\mathrm{hole}}{}}\omega _i,`$ (160)
and $`W_2`$ and so on.
In the lowest order of $`\chi `$ and $`\mathrm{\Delta }`$, we have
$`W_1^{(0)}=N_c(1r)(1w)(n2rw)^{N_c1},`$ (161)
where the factor $`N_c`$ comes from the possible position of the hole in a cell and the factor $`(1r)(1w)`$ comes from the expectation value of the hole position, $`(1\widehat{n}_i)(1\widehat{n}_i)_0`$. In the next order, by counting the number of bonds which contribute $`X`$, we obtain
$`N_c(1r)(1w)N_{1b}X(n2rw)^{N_c3},`$ (162)
where $`N_{1b}`$ means the number of bonds which are not connected the hole site as shown in Fig. 8(a). In the large enough cell we have
$`N_{1b}=N_b4.`$ (163)
Although the dotted bonds in Fig. 8(a) do not give $`X`$, they give alternative contributions to $`W_1`$ in the same order. This contribution is calculated as
$`(1\widehat{n}_i)(1\widehat{n}_i)\{\widehat{n}_j(1\widehat{n}_j)+\widehat{n}_j(1\widehat{n}_j)\}_0`$ (164)
$`=(n2rw)(1r)(1w)+Y,`$ (165)
with
$$Y=\delta (1+\delta )(\chi ^2\mathrm{\Delta }^2)4m^2(\chi ^2+\mathrm{\Delta }^2)2(\chi ^2+\mathrm{\Delta }^2)^2.$$
(166)
Here $`i`$ represents the site of the hole and $`j`$ is one of the nearest neighbor sites. Since the number of the dotted bonds in Fig. 8(a) is $`(N_bN_{1b})`$, the additional contribution to $`W_1`$ becomes
$`N_c(N_bN_{1b})Y(n2rw)^{N_c2}.`$ (167)
The higher order terms with respect to $`X`$ can be taken into account in the similar way as before and we obtain
$`W_1`$ $`=N_c(1r)(1w)(n2rw)^{N_c1}a^{N_{1b}}`$ (169)
$`+N_c(N_bN_{1b})Y(n2rw)^{N_c2}a^{\stackrel{~}{N}_b},`$
where $`\stackrel{~}{N}_b`$ appears because the diagram in Fig. 8(a) excludes the same number of bonds as in Fig. 5(a). We assume that
$`\stackrel{~}{N}_b=N_{1b}3,`$ (170)
which leads to
$`W_1`$ $`=N_c\left\{(1r)(1w)+{\displaystyle \frac{N_bN_{1b}}{n2rw}}Ya^3\right\}`$ (172)
$`\times (n2rw)^{N_c1}a^{N_{1b}}.`$
Strictly speaking, when the hole is located on the boundary of the cell, the weight should be different. However we make an approximation of large enough cell.
Continuing the similar arguments, we obtain
$`W_2`$ $`=_{N_c}C_2\left\{(1r)(1w)+{\displaystyle \frac{N_bN_{1b}}{n2rw}}Ya^3\right\}^2`$ (174)
$`\times (n2rw)^{N_c2}a^{N_{2b}},`$
for the two-hole sector, where $`N_{2b}`$ is the number of bonds which are not connected to the two sites of holes (Fig. 8(b)). We have neglected the case in which two holes come to the nearest neighbor sites. From Eq. (174) it is apparent that the $`j`$-hole sector, $`W_j`$, can be generally approximated as
$`W_j=_{N_c}C_jz^j(n2rw)^{N_cj}a^{N_{j,b}},`$ (175)
with
$`z=(1r)(1w)+{\displaystyle \frac{N_bN_{1b}}{n2rw}}Ya^3.`$ (176)
Finally approximating that $`N_{j,b}=N_b4j`$, the constraints which determine $`W`$ and $`p`$ (Eq. (89)) become
$`{\displaystyle \underset{j=0}{\overset{N_c}{}}}{\displaystyle \frac{W_j}{W}}p^j`$ $`={\displaystyle \frac{1}{W}}(n2rw+pza^4)^{N_c}a^{N_b}=1,`$ (177)
$`{\displaystyle \underset{j=0}{\overset{N_c}{}}}j{\displaystyle \frac{W_j}{W}}p^j`$ $`={\displaystyle \frac{N_c}{W}}pza^4(n2rw+pza^4)^{N_c1}a^{N_b}`$ (179)
$`=\delta N_c.`$
From these constraints we obtain
$`p`$ $`={\displaystyle \frac{\delta (n2rw)}{nz}}a^4,`$ (180)
$`W`$ $`=\left({\displaystyle \frac{n2rw}{n}}\right)^{N_c}a^{N_b},`$ (181)
for less-than-half-filling.
### B Evaluation of $`g_s^{XY}`$ and $`g_s^Z`$
In the similar way we estimate $`g_s^{XY}`$. In the case when there are $`j`$ holes in the central cell, we have
$`{\displaystyle \underset{i_0\mathrm{with}j\mathrm{holes}}{}}S_{\mathrm{}}^+S_m^{}_{i_0}`$ (182)
$`=_{N_c2}C_jz^j(n2rw)^{N_c2j}a^{\stackrel{~}{N}_b4j}S_{\mathrm{}}^+S_m^{}_0.`$ (183)
By summing up all the sectors with different number of holes, we have (according to the general formula (92)),
$`S_{\mathrm{}}^+S_m^{}`$ $`={\displaystyle \underset{j=0}{\overset{N_c2}{}}}{\displaystyle \frac{p^j}{W}}{\displaystyle \underset{i_0\mathrm{with}j\mathrm{holes}}{}}S_{\mathrm{}}^+S_m^{}_{i_0}`$ (186)
$`={\displaystyle \frac{1}{W}}(n2rw+pza^4)^{N_c2}a^{\stackrel{~}{N}_b}S_{\mathrm{}}^+S_m^{}_0`$
$`=\left({\displaystyle \frac{n}{n2rw}}\right)^2a^{(N_b\stackrel{~}{N}_b)}S_{\mathrm{}}^+S_m^{}_0.`$
From this we obtain the Gutzwiller factor
$`g_s^{XY}=\left({\displaystyle \frac{n}{n2rw}}\right)^2a^7.`$ (187)
This is a simple generalization of the results at half filling obtained in Eq. (122).
Although the estimation of $`g_s^Z`$ is essentially the same as in the half-filled case, there are some additional diagrams to be taken into account which are shown in Fig. 9. The detailed calculations are summarized in Appendix C. The final result is
$`g_s^Z`$ $`=g_s^{XY}{\displaystyle \frac{1}{4m^2+X_2}}`$ (189)
$`\times \left[X_2+4m^2\left\{1+{\displaystyle \frac{N_2X_2n}{n2rw}}\left(1{\displaystyle \frac{p}{2}}\right)a^3\right\}^2\right].`$
The quantity $`p`$ is given in Eq. (181). For $`g_s^Z`$, we use the lowest order approximation
$$p\frac{\delta (n2rw)}{nz}2\delta ,$$
(190)
since $`n2rw=1/2`$ and $`z=1/4`$ in the lowest order with respect to $`\delta ,\chi `$ and $`\mathrm{\Delta }`$. This gives our final expression for $`g_s^Z`$ given in Eq. (15).
### C Evaluation of $`g_t`$
As for the Gutzwiller factor for the hopping term, $`g_t`$, we have to calculate
$`(1\widehat{n}_i)c_i^{}c_j(1\widehat{n}_j)_{i_0}.`$ (191)
Since the hopping term needs at least one hole, the lowest sector is the one-hole sector. For this sector, we have
$`{\displaystyle \underset{i_0\mathrm{with}1\mathrm{hole}}{}}(1\widehat{n}_i)c_i^{}c_j(1\widehat{n}_j)_{i_0}`$ (192)
$`=(1\widehat{n}_i)c_i^{}c_j(1\widehat{n}_j)_0a^{\stackrel{~}{N}_b}(n2rw)^{N_c2}`$ (193)
$`=Ta^{\stackrel{~}{N}_b}(n2rw)^{N_c2},`$ (194)
where the expectation value of the correlated hopping, $`T`$, is calculated as
$`T`$ $`(1\widehat{n}_i)c_i^{}c_j(1\widehat{n}_j)_0`$ (197)
$`=(1r)(1w)\chi \chi ^3\chi \mathrm{\Delta }^2`$
$`=\chi \{(1r)(1w){\displaystyle \frac{X_2}{2}}\}.`$
It is easy to show that for the $`j+1`$-hole sector
$`{}_{N_c2}{}^{}C_{j}^{}Ta^{\stackrel{~}{N}_{j,b}}(n2rw)^{N_c2j}z^j.`$ (198)
Therefore summing up all the sectors, we obtain
$`c_i^{}c_j`$ $`={\displaystyle \frac{1}{W}}{\displaystyle \underset{j=0}{\overset{N_c2}{}}}{}_{N_c2}{}^{}C_{j}^{}Ta^{\stackrel{~}{N}_{j,b}}(n2rw)^{N_c2j}z^jp^{j+1}`$ (200)
$`=\left({\displaystyle \frac{n}{n2rw}}\right)^2Ta^{(N_b\stackrel{~}{N}_b)}p.`$
Since $`c_i^{}c_j_0=\chi `$, the Gutzwiller factor $`g_t=c_i^{}c_j/c_i^{}c_j_0`$ becomes
$`g_t`$ $`=\left({\displaystyle \frac{n}{n2rw}}\right)^2\{(1r)(1w){\displaystyle \frac{X_2}{2}}\}a^{(N_b\stackrel{~}{N}_b)}p`$ (203)
$`={\displaystyle \frac{n}{n2rw}}{\displaystyle \frac{\delta \{(1r)(1w)\frac{X_2}{2}\}}{z}}a^{(N_b\stackrel{~}{N}_b)+4},`$
where we have substituted the value of $`p`$ obtained in Eq. (181). The diagrams in Fig. 5 give contributions of the order of $`\delta ^2\chi ^2`$ or $`\delta m^2\chi ^2`$ so that they are neglected.
Finally comparing the definitions of $`X`$ and $`Y`$, we use an approximation
$`Y{\displaystyle \frac{X}{2}},`$ (204)
for the small $`\delta `$ case. Using this approximation, the quantity $`z`$ can be rewritten as
$`z`$ $`=(1r)(1w)\left\{1+{\displaystyle \frac{N_bN_{1b}}{(1r)(1w)(n2rw)}}Ya^3\right\}`$ (208)
$`=(1r)(1w)\left\{14(N_bN_{1b})X\right\}`$
$`=(1r)(1w)a^{(N_bN_{1b})}`$
$`=(1r)(1w)a^4.`$
Here we are considering $`\delta ,m`$ and $`\chi `$ as small quantities. Substituting this approximate $`z`$ into $`g_t`$ we have
$`g_t`$ $`={\displaystyle \frac{\delta n}{n2rw}}{\displaystyle \frac{(1r)(1w)\frac{X_2}{2}}{(1r)(1w)}}a^{(N_b\stackrel{~}{N}_b)+8}`$ (210)
$`={\displaystyle \frac{2\delta (1\delta )}{1\delta ^2+4m^2}}{\displaystyle \frac{(1+\delta )^24m^22X_2}{(1+\delta )^24m^2}}a.`$
This is our final expression for $`g_t`$ given in Eq. (20).
## VI Optimized variational state
In this section we calculate the energy of the variational state using the Gutzwiller factors obtained in the previous sections. We will show that the variational energies and the magnitudes of order parameters agree with the results in VMC simulations.
In the variational state $`|\psi =P_G|\psi _0(\mathrm{\Delta }_d^\mathrm{V},\mathrm{\Delta }_{\mathrm{af}}^\mathrm{V},\mu ^\mathrm{V})`$ in Eq. (2), $`|\psi _0(\mathrm{\Delta }_d^\mathrm{V},\mathrm{\Delta }_{\mathrm{af}}^\mathrm{V},\mu ^\mathrm{V})`$ is a Hartree-Fock type wave function with the $`d`$-wave SC and AF orders. It is expressed as
$`|\psi _0(\mathrm{\Delta }_d^\mathrm{V},\mathrm{\Delta }_{\mathrm{af}}^\mathrm{V},\mu ^\mathrm{V})={\displaystyle \underset{k,s(=\pm )}{}}(u_k^{(s)}+v_k^{(s)}d_k^{(s)}d_k^{(s)})|0`$ (211)
$`={\displaystyle \underset{k,s}{}}u_k^{(s)}\mathrm{exp}\left[{\displaystyle \underset{k,s}{}}{\displaystyle \frac{v_k^{(s)}}{u_k^{(s)}}}d_k^{(s)}d_k^{(s)}\right]|0,`$ (212)
where
$`{\displaystyle \frac{v_k^{(\pm )}}{u_k^{(\pm )}}}`$ $`=`$ $`{\displaystyle \frac{\pm \mathrm{\Delta }_d^\mathrm{V}\eta _k}{\left(\pm E_k\mu ^\mathrm{V}\right)+\sqrt{\left(\pm E_k\mu ^\mathrm{V}\right)^2+(\mathrm{\Delta }_d^\mathrm{V}\eta _k)^2}}},`$ (213)
$`E_k`$ $`=`$ $`\sqrt{ϵ_k^2+\mathrm{\Delta }_{\mathrm{af}}^{\mathrm{V2}}},`$ (214)
$`ϵ_k=t\gamma _k`$, $`\gamma _k=2(\mathrm{cos}k_x+\mathrm{cos}k_y)`$ and $`\eta _k=2(\mathrm{cos}k_x\mathrm{cos}k_y)`$. The annihilation operators $`d_{k\sigma }^{(s)}`$ are related to the electron operators through the following unitary transformation,
$$\left(\begin{array}{c}d_{k\sigma }^{(+)}\\ d_{k\sigma }^{()}\end{array}\right)=\left(\begin{array}{cc}\alpha _{k\sigma }& \beta _{k\sigma }\\ \beta _{k\sigma }& \alpha _{k\sigma }\end{array}\right)\left(\begin{array}{c}c_{Ak\sigma }\\ c_{Bk\sigma }\end{array}\right),$$
(215)
with
$$\{\begin{array}{c}\alpha _{k\sigma }=\sqrt{\frac{1}{2}\left(1\frac{\sigma \mathrm{\Delta }_{\mathrm{af}}^\mathrm{V}}{E_k}\right)}\\ \beta _{k\sigma }=\sqrt{\frac{1}{2}\left(1+\frac{\sigma \mathrm{\Delta }_{\mathrm{af}}^\mathrm{V}}{E_k}\right)}\end{array}.$$
(216)
Here $`c_{Ak\sigma }(c_{Bk\sigma })`$ are annihilation operators of an electron on the A(B)-sublattice and $`\sigma `$ represent $``$(+1) and $``$(-1). The wave vector $`𝒌`$ is limited to half of the Brillouin zone where $`ϵ_k<0`$. We can confirm that $`|\psi _0`$ is a vacuum of the annihilation operators which diagonalize
$`{\displaystyle \underset{k}{}}`$ $`[{\displaystyle \underset{\sigma }{}}\{ϵ_k(c_{Ak\sigma }^{}c_{Bk\sigma }+h.c.)`$ (219)
$`(\mu ^\mathrm{V}+\sigma \mathrm{\Delta }_{\mathrm{af}}^\mathrm{V})c_{Ak\sigma }^{}c_{Ak\sigma }(\mu ^\mathrm{V}\sigma \mathrm{\Delta }_{\mathrm{af}}^\mathrm{V})c_{Bk\sigma }^{}c_{Bk\sigma }\}`$
$`\mathrm{\Delta }_d^\mathrm{V}\eta _k(c_{Ak}c_{Bk}+c_{Bk}c_{Ak}+h.c.)].`$
In order to clarify the correspondence to the mean-field theory, let us consider the effective Hamiltonian, $`\widehat{}_{\mathrm{eff}}`$ in Eq. (5). In $`\widehat{}_{\mathrm{eff}}`$ the parameter $`t`$ in $``$ is replaced with
$$t_{\mathrm{eff}}=g_tt,$$
(220)
and the exchange term is replaced with
$`J𝑺_i𝑺_j=g_s^{XY}J\left(S_i^xS_j^x+S_i^yS_j^y\right)+g_s^ZJS_i^zS_j^z,`$ (221)
where $`g_t`$, $`g_s^{XY}`$ and $`g_s^Z`$ are the Gutzwiller factors obtained in the previous sections. When we apply the mean-field theory to $`\widehat{}_{\mathrm{eff}}`$, we obtain the similar Hamiltonian as in Eq. (219) but with the replacements
$`t`$ $`t_{\mathrm{eff}}+J_{\mathrm{eff}}\chi ^\mathrm{V},`$ (222)
$`\mathrm{\Delta }_d^\mathrm{V}`$ $`J_{\mathrm{eff}}\mathrm{\Delta }^\mathrm{V},`$ (223)
$`\mathrm{\Delta }_{\mathrm{af}}^\mathrm{V}`$ $`2J_{\mathrm{eff}}^Zm^\mathrm{V},`$ (224)
with
$$J_{\mathrm{eff}}=\frac{1}{2}g_s^{XY}J+\frac{1}{4}g_s^ZJ,$$
(225)
$$J_{\mathrm{eff}}^Z=g_s^ZJ.$$
(226)
Note that in the usual mean-field theory the self-consistency equations give $`\chi ^\mathrm{V}=\chi ,\mathrm{\Delta }^\mathrm{V}=\mathrm{\Delta }`$ and $`m^\mathrm{V}=m`$, where $`\chi ,\mathrm{\Delta }`$ and $`m`$ are the expectation values $`\mathrm{\Delta }=c_i^{}c_j^{}_0`$, $`\chi =c_{i\sigma }^{}c_{j\sigma }_0`$, and $`m=\frac{1}{2}(1)^i\left(\widehat{n_i}_0\widehat{n_i}_0\right)`$ used in the previous sections. We will see shortly that $`\chi ^\mathrm{V}`$ and $`\chi `$ etc. are slightly different due to the dependence of the Gutzwiller factors on $`\chi ,\mathrm{\Delta }`$ and $`m`$.
Using the wave function $`|\psi _0`$, the expectation values become
$`\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{1}{8N}}{\displaystyle \underset{k,\pm }{}}{\displaystyle \frac{F_k\eta _k}{\sqrt{(\pm E_k\mu ^\mathrm{V})^2+F_k^2}}},`$ (227)
$`\chi `$ $`=`$ $`{\displaystyle \frac{1}{8N}}{\displaystyle \underset{k,\pm }{}}{\displaystyle \frac{(t_{\mathrm{eff}}+J_{\mathrm{eff}}\chi ^\mathrm{V})\gamma _k^2(\pm E_k\mu ^\mathrm{V})}{\pm E_k\sqrt{(\pm E_k\mu ^\mathrm{V})^2+F_k^2}}},`$ (228)
$`m`$ $`=`$ $`{\displaystyle \frac{1}{2N}}{\displaystyle \underset{k,\pm }{}}{\displaystyle \frac{2J_{\mathrm{eff}}^Zm^\mathrm{V}(\pm E_k\mu ^\mathrm{V})}{\pm E_k\sqrt{(\pm E_k\mu ^\mathrm{V})^2+F_k^2}}},`$ (229)
$`n`$ $`=`$ $`1{\displaystyle \frac{1}{N}}{\displaystyle \underset{k,\pm }{}}{\displaystyle \frac{(\pm E_k\mu ^\mathrm{V})}{\sqrt{(\pm E_k\mu ^\mathrm{V})^2+F_k^2}}},`$ (230)
with
$`F_k`$ $`=`$ $`J_{\mathrm{eff}}\mathrm{\Delta }^\mathrm{V}\eta _k,`$ (231)
$`E_k`$ $`=`$ $`\sqrt{(t_{\mathrm{eff}}+J_{\mathrm{eff}}\chi ^\mathrm{V})^2\gamma _k^2+(2J_{\mathrm{eff}}^Zm^\mathrm{V})^2},`$ (232)
where $`N`$ is the number of sites and the summation over $`𝒌`$ is limited to half of the Brillouin zone.
Using these expectation values we obtain
$`E_{\mathrm{var}}`$ $`=`$ $`\widehat{}_{\mathrm{eff}}_0`$ (233)
$`=`$ $`8Nt_{\mathrm{eff}}\chi _{}4NJ_{\mathrm{eff}}\left(\mathrm{\Delta }^2+\chi ^2\right)`$ (235)
$`2NJ_{\mathrm{eff}}^Zm^2\mu (NnN_\mathrm{e}).`$
By taking the derivatives with respect to the variational parameters, $`\mathrm{\Delta }^\mathrm{V},\chi ^\mathrm{V}`$, and $`m^\mathrm{V}`$, we find the following self-consistency equations:
$`\mathrm{\Delta }^\mathrm{V}`$ $`=`$ $`\mathrm{\Delta }{\displaystyle \frac{1}{8NJ_{\mathrm{eff}}}}{\displaystyle \frac{H_{\mathrm{eff}}}{\mathrm{\Delta }}}_0,`$ (236)
$`\chi ^\mathrm{V}`$ $`=`$ $`\chi {\displaystyle \frac{1}{8NJ_{\mathrm{eff}}}}{\displaystyle \frac{H_{\mathrm{eff}}}{\chi }}_0,`$ (237)
$`m^\mathrm{V}`$ $`=`$ $`m{\displaystyle \frac{1}{4NJ_{\mathrm{eff}}^Z}}{\displaystyle \frac{H_{\mathrm{eff}}}{m}}_0,`$ (238)
$`\mu ^\mathrm{V}`$ $`=`$ $`\mu {\displaystyle \frac{1}{N}}{\displaystyle \frac{H_{\mathrm{eff}}}{n}}_0,`$ (239)
where the partial derivative of $`H_{\mathrm{eff}}`$ is applied to the Gutzwiller factors, $`g_s^{XY},g_s^Z`$ and $`g_t`$.
Figure 10 shows the self-consistent parameters $`\mathrm{\Delta },\chi `$ and $`m`$ satisfying (239) as a function of the doping $`\delta =1n`$ for $`J/t=0.3`$. Note that these are the expectation values in the wave function $`|\psi _0`$ without the projection. The dashed line in Fig. 10 represents the results when the AF order is suppressed, i.e., $`m`$ is fixed to zero. It is apparent that the presence of the AF order has a small effect on the expectation values $`\mathrm{\Delta }`$ and $`\chi `$. The coexistent state between the d-wave SC and AF order parameters is stabilize up to the doping rate $`\delta =0.1`$. This is consistent with the VMC simulations.
The actual expectation values in the wave function $`P_G|\psi _0`$ with projection are different from $`\mathrm{\Delta },\chi `$ and $`m`$. For example, the expectation value of $`c_{i\sigma }^{}c_{j\sigma }`$ is
$$\chi _{\mathrm{exp}}c_{i\sigma }^{}c_{j\sigma }=g_tc_{i\sigma }^{}c_{j\sigma }_0=g_t\chi .$$
(240)
In the similar way, the actual expectation values in $`P_G|\psi _0`$ are obtained using corresponding Gutzwiller factors. For $`m_{\mathrm{exp}}`$ we repeat the similar arguments as for $`g_s^Z`$ to obtain
$`m_{\mathrm{exp}}{\displaystyle \frac{1}{2}}n_in_i=g_mm,`$ (241)
$`g_m={\displaystyle \frac{n}{n2rw}}a^4\left[1+{\displaystyle \frac{(N_bN_{1b})X_2n}{n2rw}}\left(1{\displaystyle \frac{p}{2}}\right)a^3\right].`$ (242)
(243)
For the expectation value of $`c_i^{}c_j^{}`$, we have to take care of the fact that the number of holes changes in the cell. Here we use the average number of holes as an approximation. In the same level of approximations as for $`g_t`$ we obtain
$`\mathrm{\Delta }_{\mathrm{exp}}`$ $`c_i^{}c_j^{}=g_\mathrm{\Delta }\mathrm{\Delta },`$ (244)
$`g_\mathrm{\Delta }`$ $`={\displaystyle \frac{\delta n}{n2rw}}{\displaystyle \frac{(1r)^2+(1w)^2+X_2}{2(1r)(1w)}}a^{(N_b\stackrel{~}{N}_b)+8}`$ (246)
$`={\displaystyle \frac{2\delta (1\delta )}{1\delta ^2+4m^2}}{\displaystyle \frac{(1+\delta )^2+4m^2+2X_2}{(1+\delta )^24m^2}}a.`$
Figure 11 shows the actual expectation values, $`\mathrm{\Delta }_{\mathrm{exp}}`$ and $`m_{\mathrm{exp}}`$ as a function of the doping $`\delta =1n`$ for $`J/t=0.3`$.
## VII Summary and Discussion
In this paper we have developed a new approach for studying the effect of strong correlation or the Gutzwiller’s projection in the two-dimensional $`t`$-$`J`$ model. It is based on the extended Gutzwiller approximation, in which the effects of longer-range correlations are taken into account. These correlations play important roles for the interplay between the AF and d-wave SC. Let us summarize our main results and discuss related problems.
(1) Generally the expectation values with respect to the projected wave function are strongly renormalized due to the exclusion of double occupancies. In the slave-boson mean-field theory, this effect is taken into account by assuming the replacement $`t\delta t`$ and $`JJ`$. However we have shown that the renormalization is not so simple. First of all the renormalization factor for the exchange interaction is anisotropic, i.e., $`J_{\mathrm{eff}}^{XY}=g_s^{XY}JJ_{\mathrm{eff}}^Z=g_s^ZJ`$ in the presence of AF moment. Furthermore $`g_t,g_s^{XY}`$ and $`g_s^Z`$ have nonlinear dependences on the expectation values $`\mathrm{\Delta },\chi `$ and $`m`$. In this sense the extended Gutzwiller approximation is beyond the simple slave-boson mean-field theory. We think that essence of the strong correlation is contained in these Gutzwiller factors, because they stem solely from the projection.
The physical meanings of the Gutzwiller factors are clarified:
(i) The factor $`a^7`$ appearing in $`g_s^{XY}`$ and $`g_s^Z`$ represents the exclusion effect shown in Fig. 4(b). When we calculate $`𝑺_i𝑺_j`$, the surrounding six bonds cannot make contribution of nearest neighbor correlation $`X`$. As is evident in Fig. 1, the factor $`a^7`$ reduces the value of $`J_{\mathrm{eff}}`$ even for the case $`m=0`$.
(ii) The enhancement of $`g_s^Z`$ as a function of $`m`$ is the most important feature. This was observed in the VMC simulations. Here we have identified its origin as the effect of surrounding AF correlations by discussing the probabilities of spin configurations in Fig. 7. The enhancement of $`g_s^Z`$ is due to the increase of $``$ configuration caused by the AF circumstances.
(2) We have studied the projected variational states as an application of the new Gutzwiller approximation. It is shown that our scheme reproduces the results in the VMC simulations. The enhancement of $`g_s^Z`$ as a function of $`m`$ is essential for reproducing the coexistent state of the AF and d-wave SC orders found in VMC.
(3) At half filling ($`\delta =0`$), there are several interesting features. In the original Gutzwiller approximation, the AF state at half filling was not obtained which was unphysical. However in the present Gutzwiller approximation, the AF state with d-wave SC order parameter becomes the most optimized state as obtained in VMC. Note that the expectation value of the SC order parameter, $`\mathrm{\Delta }_{\mathrm{exp}}`$, is zero because of the Gutzwiller factor $`g_\mathrm{\Delta }=0`$ at half filling, although $`\mathrm{\Delta }`$ and the variational parameter $`\mathrm{\Delta }^V`$ are nonzero. The self-consistent values which we obtain are (Fig. 10)
$`\mu ^V`$ $`=0,`$ (247)
$`\mathrm{\Delta }^V`$ $`=\chi ^V=0.080,`$ (248)
$`\mathrm{\Delta }`$ $`=\chi =0.163,`$ (249)
$`m^V`$ $`=0.017,`$ (250)
$`m`$ $`=0.147.`$ (251)
The actual expectation value with the projection, $`m_{\mathrm{exp}}`$, is (see Fig. 11)
$$m_{\mathrm{exp}}=g_mm=0.373.$$
(252)
This is a reasonable magnitude and close to the Monte Carlo result for the Heisenberg model ($`m_{\mathrm{exp}}=0.31\pm 0.02`$).
The relations $`\mathrm{\Delta }^V=\chi ^V`$ and $`\mathrm{\Delta }=\chi `$ are the manifestation of the SU(2) symmetry at half filling. Actually the self-consistency equations for $`\chi `$ and $`\mathrm{\Delta }`$ become the same at $`\delta =0`$. It is worth while noting here that, due to the SU(2) symmetry, the coexistent state of the AF and d-wave SC is equivalent to the $`\pi `$-flux state with AF long range order discussed by Hsu.
The self-consistent value, $`m=0.147`$ in Eq. (251), is smaller than the result $`m=1/2`$ in the simple mean-field theory. Instead it is slightly larger than the value at which the Gutzwiller factor $`g_s^Z`$ has a maximum, as shown in Fig. 1. If $`m`$ is increased from the self-consistent value, $`g_s^Z`$ decreases. This indicates that $`m`$ is determined so as to minimize the exchange energy $`2NJ_{\mathrm{eff}}^Zm^2=2Ng_s^ZJm^2`$ by optimizing the energy gain due to the long-range order and the energy loss due to the reduction of $`g_s^Z`$. This situation is, in some sense, similar to the quantum fluctuations of the Heisenberg spin system discussed in the spin-wave theory: The spin fluctuation, which reduces $`m`$ from the mean-field value $`1/2`$, leads to the enhancement of the Gutzwiller’s renormalization factors to gain the energy.
(4) For less-than-half-filled case, the difference between the variational parameter $`\mathrm{\Delta }^V`$ and the expectation value $`\mathrm{\Delta }_{\mathrm{exp}}`$ is remarkable: $`\mathrm{\Delta }^V`$ is finite and moreover it increases as the doping rate $`\delta `$ decreases, while $`\mathrm{\Delta }_{\mathrm{exp}}`$ is proportional to $`\delta `$ near half filling. This is because the Gutzwiller factor, $`g_\mathrm{\Delta }`$, is proportional to $`\delta `$ (see section VI). We interpret that $`\mathrm{\Delta }^V`$ is the BCS-type energy gap observed in scanning tunnel spectroscopy or in break junctions, because it is the parameter embedded in the wave function even at half filling. The excitation spectra will have a large energy gap corresponding to $`\mathrm{\Delta }^V`$, while the true long-range order, $`\mathrm{\Delta }_{\mathrm{exp}}`$, is reduced due to the projection or the strong correlation. The increase of $`\mathrm{\Delta }^V`$ in decreasing $`\delta `$ is consistent with the dependence observed experimentally.
(5) We discuss here the relation to the SO(5) theory. In our formulation, the combination $`m^2+\mathrm{\Delta }^2`$ appears frequently. If the SO(5) symmetry is exact in the $`t`$-$`J`$ model, the free energy will have a systematic dependence such as
$$F(\mathrm{\Delta }^2+m^2)\mathrm{or}F(\mathrm{\Delta }_{\mathrm{exp}}^2+m_{\mathrm{exp}}^2).$$
(253)
For example, the numerator in $`a`$ has a combination
$$16(m^2+\mathrm{\Delta }^2+\chi ^2).$$
(254)
Also the factor giving the enhancement of $`g_s^Z`$ contains
$$4m^2+X_2=4m^2+2\mathrm{\Delta }^2+2\chi ^2.$$
(255)
Although these combinations remind us of the SO(5) symmetry, our Gutzwiller approximation does not show exact symmetry.
Nevertheless a tendency similar to the SO(5) prediction can be seen from the Gutzwiller factors in Fig. 1. Comparing $`g_s^{XY}`$ and $`g_s^Z`$ for the cases with $`\mathrm{\Delta }=0.02`$ and $`0.18`$, we find that $`g_s^{XY}`$ and $`g_s^Z`$ are larger for smaller $`\mathrm{\Delta }`$. This is mainly from the exclusion effect of $`a`$, because even at $`m=0`$ this effect is observed. Therefore if we consider a situation where the d-wave SC order parameter is suppressed, then the Gutzwiller factor is enhanced causing the increase of AF moment. This gives the similar phenomena predicted in the SO(5) theory.
(6) One advantage of the present theory is that it is easily applied to the inhomogeneous cases, such as the stripe state, vortex cores, and magnetic states around nonmagnetic or magnetic impurities, where the interplay between the AF and d-wave SC correlations plays an important role. Since our Gutzwiller approximation gives a reasonable estimate of the variational energies in the presence of AF correlations, a reliable analytic formulation can be given.
The simplest way of applying the present scheme to inhomogeneous problems is to assume that $`g_s^{XY},g_s^Z`$ and $`g_t`$ for each bond $`i,j`$ are determined locally from the expectation values $`\mathrm{\Delta }_{ij},\chi _{ij}`$ and $`m_i,m_j`$ for the bond. In this case, if the d-wave order parameter is reduced around vortex cores, impurities or stripes, then the Gutzwiller factors $`g_s^{XY}`$ and $`g_s^Z`$ are enhanced locally as expected from Fig. 1. This effect causes the local development of AF correlations, which can be observed experimentally. From these viewpoint, preliminary calculations for the vortex cores and stripe states have been published elsewhere.
Of course the VMC simulations give more accurate evaluation of the variational energies, if they are used for the inhomogeneous systems like stripe states. However there are some difficulties in applying the VMC. Firstly we have to treat a fairly large unit cell to study the slowly varying order parameters especially near half filling. For example, the incommensurability in the stripe state is close to $`(\pi ,\pi )`$ so that the period of the stripe pattern becomes fairly long. In this case the VMC simulations become difficult. Furthermore the choice of the functional form of the trial state in VMC is restricted, since only the small number of variational parameters can be used practically. On the contrary, our scheme based on the mean-field-type Gutzwiller approximation can be used in a fairly large system sizes. Moreover the order parameters on all the bonds, $`\mathrm{\Delta }_{ij},\chi _{ij}`$ and $`m_i`$ can be optimized in the similar sense to unrestricted Hartree-Fock theory. Therefore we can search for the microscopically optimized variational states in our scheme.
The authors wish to thank T. M. Rice, H. Fukuyama, Y. Tanaka, H. Tsuchiura and M. Sigrist for useful discussions. They also thank C.-M. Ho for critical reading of the manuscript. This work is supported in part by a Grant-in-Aid of of the Ministry of Education, Science, Sports and Culture.
## A Derivation of the expectation value of an operator
The evaluation of Eq. (90) is slightly complicated. The constraints for $`\{N_i^{}\}`$ are
$`{\displaystyle \underset{i=1}{\overset{K}{}}}n_iN_i^{}`$ $`={\displaystyle \frac{N_e}{2}}n_{i_0},{\displaystyle \underset{i=1}{\overset{K}{}}}n_iN_i^{}={\displaystyle \frac{N_e}{2}}n_{i_0},`$ (A1)
$`{\displaystyle \underset{i=1}{\overset{K}{}}}n_{\mathrm{h}i}N_i^{}`$ $`=NN_en_{\mathrm{h}i_0},`$ (A2)
instead of (75) because $`N_i^{}`$ represents the number of cells of the $`i`$-th state except for the central cell.
In the same way as in the denominator, the largest term is given by
$`\overline{N_i^{}}={\displaystyle \frac{N}{N_c}}{\displaystyle \frac{\omega _i}{W^{}}}(p^{})^{n_{\mathrm{h}i}},`$ (A3)
with slightly different values of $`W^{}`$ and $`p^{}`$ because of the difference of the constraints. From the constraints (A1), we can see that $`\overline{N_i^{}}`$ satisfies the relations
$$\underset{i=1}{\overset{K}{}}\overline{N_i^{}}=\frac{N}{N_c}1,\underset{i=1}{\overset{K}{}}n_{\mathrm{h}i}\overline{N_i^{}}=NN_en_{\mathrm{h}i_0}.$$
(A4)
Therefore if we define the difference $`\mathrm{\Delta }\overline{N_i}=\overline{N_i^{}}\overline{N_i}`$, we have important relations
$`{\displaystyle \underset{i=1}{\overset{K}{}}}\mathrm{\Delta }\overline{N_i}=1,{\displaystyle \underset{i=1}{\overset{K}{}}}n_{\mathrm{h}i}\mathrm{\Delta }\overline{N_i}=n_{\mathrm{h}i_0}.`$ (A5)
By use of these, the ratio between the numerator and the denominator in Eq. (92) is calculated as follows:
$`\widehat{𝒪}`$ $`={\displaystyle \frac{\psi _0|P_G\widehat{𝒪}P_G|\psi _0}{\psi _0|P_GP_G|\psi _0}}`$ (A11)
$`={\displaystyle \underset{i_0}{}}{\displaystyle \frac{1}{\left(\frac{N}{N_c}\right)}}{\displaystyle \underset{i=1}{\overset{K}{}}}{\displaystyle \frac{\overline{N_i}!}{\overline{N_i^{}}!}}{\displaystyle \underset{i=1}{\overset{K}{}}}\omega _i^{\overline{N_i^{}}\overline{N_i}}\widehat{𝒪}_{i_0}`$
$`={\displaystyle \underset{i_0}{}}{\displaystyle \frac{N_c}{N}}{\displaystyle \underset{i=1}{\overset{K}{}}}\left({\displaystyle \frac{\overline{N_i}}{\omega _i}}\right)^{\mathrm{\Delta }\overline{N_i}}\widehat{𝒪}_{i_0}`$
$`={\displaystyle \underset{i_0}{}}{\displaystyle \frac{N_c}{N}}{\displaystyle \underset{i=1}{\overset{K}{}}}\left({\displaystyle \frac{N}{N_c}}{\displaystyle \frac{p^{n_{\mathrm{h}i}}}{W}}\right)^{\mathrm{\Delta }\overline{N_i}}\widehat{𝒪}_{i_0}`$
$`={\displaystyle \underset{i_0}{}}{\displaystyle \frac{N_c}{N}}\left({\displaystyle \frac{N}{N_c}}{\displaystyle \frac{1}{W}}\right)^{{\scriptscriptstyle \mathrm{\Delta }\overline{N_i}}}\times p^{{\scriptscriptstyle n_{\mathrm{h}i}\mathrm{\Delta }\overline{N_i}}}\widehat{𝒪}_{i_0}`$
$`={\displaystyle \underset{i_0}{}}{\displaystyle \frac{p^{n_{\mathrm{h}i_0}}}{W}}\widehat{𝒪}_{i_0}.`$
This gives Eq. (92).
## B Reproduction of the original Gutzwiller approximation
In this appendix we show that the original Gutzwiller approximation described in II is reproduced if we assume only the site-diagonal expectation values in the generalized formulation in section III.
First we calculate the weight $`\omega _i`$ of the $`i`$-th state which has $`n_i`$ up-spin electrons, $`n_i`$ down-spin electrons and $`n_{\mathrm{h}i}`$ holes. If we take only the site-diagonal expectation values, we have
$$\omega _i=[r(1w)]^{n_{\mathrm{right}}}[w(1r)]^{n_{\mathrm{wrong}}}[(1r)(1w)]^{n_{\mathrm{h}i}},$$
(B1)
where $`n_{\mathrm{right}}`$ ($`n_{\mathrm{wrong}}`$) means the number of sites where the right (wrong) spices of spin direction is located depending on the sublattice 1 and 2. Then the total weight in the subgroup with $`j`$ holes (Eq. (87)) can be calculated exactly as
$`W_j`$ $`={\displaystyle \underset{i\mathrm{with}j\mathrm{holes}}{}}\omega _i`$ (B4)
$`=_{N_c}C_j[r(1w)+w(1r)]^{N_cj}\left[(1r)(1w)\right]^j`$
$`=_{N_c}C_j(n2rw)^{N_cj}(1r)^j(1w)^j,`$
where $`{}_{N_c}{}^{}C_{j}^{}`$ is the number of choices of the positions of $`j`$ holes.
Using these $`W_j`$, the constraints (89) have simple forms as
$`{\displaystyle \underset{j=0}{\overset{N_c}{}}}{\displaystyle \frac{W_j}{W}}p^j={\displaystyle \frac{1}{W}}\{n2rw+(1r)(1w)p\}^{N_c}=1,`$ (B5)
$`{\displaystyle \underset{j=0}{\overset{N_c}{}}}j{\displaystyle \frac{W_j}{W}}p^j={\displaystyle \frac{N_c}{W}}(1r)(1w)p`$ (B6)
$`\times \{n2rw+(1r)(1w)p\}^{N_c1}=\delta N_c.`$ (B7)
These can be solved easily to give
$`p={\displaystyle \frac{\delta (n2rw)}{n(1r)(1w)}},`$ (B8)
$`W=\left({\displaystyle \frac{n2rw}{n}}\right)^{N_c}.`$ (B9)
Finally the expectation value for $`S_{\mathrm{}}^+S_m^{}`$, for example, becomes
$`{\displaystyle \frac{\psi _0|P_GS_{\mathrm{}}^+S_m^{}P_G|\psi _0}{\psi _0|P_GP_G|\psi _0}}`$ (B10)
$`={\displaystyle \underset{i_0}{}}{\displaystyle \frac{p^{n_{\mathrm{h}i_0}}}{W}}S_{\mathrm{}}^+S_m^{}_{i_0}`$ (B11)
$`={\displaystyle \underset{j=0}{\overset{N_c2}{}}}{\displaystyle \frac{p^j}{W}}_{N_c2}C_j(n2rw)^{N_c2j}(1r)^j(1w)^jS_{\mathrm{}}^+S_m^{}_0`$ (B12)
$`={\displaystyle \frac{1}{W}}\{n2rw+(1r)(1w)p\}^{N_c2}S_{\mathrm{}}^+S_m^{}_0`$ (B13)
$`={\displaystyle \frac{n^2}{(n2rw)^2}}S_{\mathrm{}}^+S_m^{}_0,`$ (B14)
which is exactly the same as the results in the original Gutzwiller approximation in section II.B, Eq. (65).
## C Evaluation of $`g_s^Z`$ for the less-than-half-filled case
In the similar way to $`g_s^{XY}`$, we evaluate $`g_s^Z`$. In the zero-hole sector, we have
$`{\displaystyle \underset{i_0\mathrm{with}0\mathrm{holes}}{}}S_{\mathrm{}}^zS_m^z_{i_0}`$ $`=S_{\mathrm{}}^zS_m^z_0(n2rw)^{N_c2}a^{\stackrel{~}{N}_b}`$ (C1)
$``$ $`2N_2m^2X_2(n2rw)^{N_c3}a^{\stackrel{~}{N}_b3}`$ (C2)
$``$ $`N_2^2m^2X_2^2(n2rw)^{N_c4}a^{\stackrel{~}{N}_b6}.`$ (C4)
For the one-hole sector, there are five contributions from the diagrams in Fig. 9. The first three diagrams are simple extensions of the zero-hole sector. However Figs. 9(d) and 9(e) are new-type contributions due to the presence of a hole. For the diagrams in Figs. 9(d) and 9(e) we calculate
$`\left(\begin{array}{ccc}& & \\ & & \end{array}\right)`$ $`=S_{\mathrm{}}^zS_m^z(1\widehat{n}_m^{})(1\widehat{n}_m^{})_c`$ (C6)
$`={\displaystyle \frac{1}{2}}m^2X_2,`$
$`\left(\begin{array}{ccc}& & \\ & & \\ & & \end{array}\right)`$ $`=\{\widehat{n}_{\mathrm{}^{}}(1\widehat{n}_{\mathrm{}^{}})+\widehat{n}_{\mathrm{}^{}}(1\widehat{n}_{\mathrm{}^{}})\}S_{\mathrm{}}^zS_m^z`$ (C9)
$`\times (1\widehat{n}_m^{})(1\widehat{n}_m^{})_c`$
$`={\displaystyle \frac{1}{2}}m^2X_2^2.`$
Using these expectation values, we obtain
$`{\displaystyle \underset{i_0\mathrm{with}1\mathrm{holes}}{}}S_{\mathrm{}}^zS_m^z_{i_0}`$ (C10)
$`=_{N_c2}C_1S_{\mathrm{}}^zS_m^z_0z(n2rw)^{N_c3}a^{\stackrel{~}{N}_{1b}}`$ (C11)
$`_{N_c3}C_12N_2m^2X_2z(n2rw)^{N_c4}a^{\stackrel{~}{N}_{1b}3}`$ (C12)
$`_{N_c4}C_1N_2^2m^2X_2^2z(n2rw)^{N_c5}a^{\stackrel{~}{N}_{1b}6}`$ (C13)
$`+N_2m^2X_2(n2rw)^{N_c3}a^{\stackrel{~}{N}_b3}`$ (C14)
$`+N_2^2m^2X_2^2(n2rw)^{N_c4}a^{\stackrel{~}{N}_b6}.`$ (C15)
In the two hole sector, we need to calculate a diagram in Fig. 9(f) which is $`\frac{1}{4}m^2X_2^2`$. Taking account of these diagrams and counting the higher order terms with respect to $`X`$, we approximate as
$`{\displaystyle \underset{i_0\mathrm{with}2\mathrm{holes}}{}}S_{\mathrm{}}^zS_m^z_{i_0}`$ (C16)
$`=_{N_c2}C_2S_{\mathrm{}}^zS_m^z_0z^2(n2rw)^{N_c4}a^{\stackrel{~}{N}_{2b}}`$ (C17)
$`_{N_c3}C_22N_2m^2X_2z^2(n2rw)^{N_c5}a^{\stackrel{~}{N}_{2b}3}`$ (C18)
$`_{N_c4}C_2N_2^2m^2X_2^2z^2(n2rw)^{N_c6}a^{\stackrel{~}{N}_{2b}6}`$ (C19)
$`+_{N_c3}C_1N_2m^2X_2z(n2rw)^{N_c4}a^{\stackrel{~}{N}_{1b}3}`$ (C20)
$`+_{N_c4}C_1N_2^2m^2X_2^2z(n2rw)^{N_c5}a^{\stackrel{~}{N}_{1b}6}`$ (C21)
$`{\displaystyle \frac{N_2^2}{4}}m^2X_2^2(n2rw)^{N_c4}a^{\stackrel{~}{N}_b6}.`$ (C22)
Finally generalization to higher order contributions leads to
$`S_{\mathrm{}}^zS_m^z={\displaystyle \frac{1}{W}}{\displaystyle \underset{j=0}{\overset{N_c}{}}}p^j{\displaystyle \underset{i_0\mathrm{with}j\mathrm{holes}}{}}S_{\mathrm{}}^zS_m^z_{i_0}`$ (C23)
$`={\displaystyle \frac{1}{W}}{\displaystyle \underset{j=0}{\overset{N_c2}{}}}p_{N_c2}^jC_jS_{\mathrm{}}^zS_m^z_0z^j(n2rw)^{N_c2j}a^{\stackrel{~}{N}_{j,b}}`$ (C24)
$`{\displaystyle \frac{1}{W}}{\displaystyle \underset{j=0}{\overset{N_c3}{}}}p_{N_c3}^jC_j2N_2m^2X_2z^j(n2rw)^{N_c3j}a^{\stackrel{~}{N}_{j,b}3}`$ (C25)
$`{\displaystyle \frac{1}{W}}{\displaystyle \underset{j=0}{\overset{N_c4}{}}}p_{N_c4}^jC_jN_2^2m^2X_2^2z^j(n2rw)^{N_c4j}a^{\stackrel{~}{N}_{j,b}6}`$ (C26)
$`+{\displaystyle \frac{1}{W}}{\displaystyle \underset{j=0}{\overset{N_c3}{}}}p_{N_c3}^{j+1}C_jN_2m^2X_2z^j(n2rw)^{N_c3j}a^{\stackrel{~}{N}_{j,b}3}`$ (C27)
$`+{\displaystyle \frac{1}{W}}{\displaystyle \underset{j=0}{\overset{N_c4}{}}}p_{N_c4}^{j+1}C_jN_2^2m^2X_2^2z^j(n2rw)^{N_c4j}a^{\stackrel{~}{N}_{j,b}6}`$ (C28)
$`{\displaystyle \frac{1}{W}}{\displaystyle \underset{j=0}{\overset{N_c4}{}}}p_{N_c4}^{j+2}C_j{\displaystyle \frac{N_2^2}{4}}m^2X_2^2z^j(n2rw)^{N_c4j}a^{\stackrel{~}{N}_{j,b}6}`$ (C29)
$`=\left({\displaystyle \frac{n}{n2rw}}\right)^2S_{\mathrm{}}^zS_m^z_0a^{(N_b\stackrel{~}{N}_b)}`$ (C30)
$`\left({\displaystyle \frac{n}{n2rw}}\right)^32N_2m^2X_2\left(1{\displaystyle \frac{p}{2}}\right)a^{(N_b\stackrel{~}{N}_{b}^{}{}_{}{}^{})}`$ (C31)
$`\left({\displaystyle \frac{n}{n2rw}}\right)^4N_2^2m^2X_2^2\left(1{\displaystyle \frac{p}{2}}\right)^2a^{(N_b\stackrel{~}{N}_{b}^{}{}_{}{}^{\prime \prime })}`$ (C32)
$`=\left({\displaystyle \frac{n}{n2rw}}\right)^2a^{(N_b\stackrel{~}{N}_b)}`$ (C33)
$`\times \left[{\displaystyle \frac{X_2}{4}}m^2\left\{1+{\displaystyle \frac{N_2X_2n}{n2rw}}\left(1{\displaystyle \frac{p}{2}}\right)a^3\right\}^2\right].`$ (C34)
Thus the Gutzwiller factor is
$$g_s^Z=g_s^{XY}\frac{1}{4m^2+X_2}[X_2+4m^2\{1+\frac{N_2X_2n}{n2rw}(1\frac{p}{2})a^3)^2].$$
(C35)
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# Spin transport in interacting quantum wires and carbon nanotubes
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## Abstract
We present a general formulation of spin-dependent transport through a clean one-dimensional interacting quantum wire or carbon nanotube, connected to non-collinear ferromagnets via tunnel junctions. We show that the low energy description of each junction is given by a conformally-invariant boundary condition representing exchange coupling, in addition to a pair of electron tunneling operators. The effects of the exchange coupling are strongly enhanced by interactions, leading to a dramatic suppression of spin accumulation. This is a direct signature of spin-charge separation in a Luttinger liquid. Furthermore, we demonstrate that magnetic polarization can lead to oscillations in the non-linear current-voltage relation. This phenomena is a surprising purely nonequilibrium effect due to backscattering interactions, which are thus dangerously marginally irrelevant in the repulsively-interacting Luttinger liquid.
\]
Recent studies on metal-ferromagnet hybrid systems have revealed new and interesting physics due to the interplay between the electronic charge and spin , e.g., the giant magnetoresistance effect . Following the initial spin-injection proposal , the work of Johnson and Silsbee and subsequent advances have opened the way for the field of spintronics, where the electron spin is the central element for information storage and transport. Spin-dependent transport plays an important role in quantum computation proposals, and has already led to new technological applications. In this context, a detailed understanding of transport through ferromagnetic-normal-ferromagnetic devices is both of fundamental and technological interest. In such structures, the current-voltage relation is predicted to sensitively depend on the relative angle $`\theta `$ between the magnetization directions of the ferromagnets (FMs).
Current theoretical models are based on Fermi liquid theory, thereby ignoring the effect of interactions in the metal. As the inevitable miniaturization of spin-dependent devices proceeds, however, at least the interconnects must ultimately reach the one-dimensional (1D) quantum limit, in which Fermi liquid theory breaks down. This theoretically expected change from Fermi liquid to Luttinger liquid (LL) behavior drastically alters transport phenomena, as has recently been verified in experiments on charge conduction in carbon nanotubes, which are nearly ideal 1D quantum wires (QWs). Despite these developments, spin injection into a LL has received surprisingly little attention . In this paper, we present a general low-energy theory for spin transport in a LL, which directly applies to nanotubes and semiconductor QWs . We assume, as expected theoretically and recently observed experimentally for carbon nanotubes, that spin-orbit coupling in the LL is negligible. Its inclusion is, however, straightforward.
Our analysis shows that spin transport in LLs is qualitatively different both from charge transport in LLs and from Fermi liquid spin transport. We focus for concrete results on the case of an end-contacted quantum wire, and assume that the distance $`L`$ between contacts is sufficiently long, $`\mathrm{max}(V,T)v/L`$, to ensure an incoherent stepwise transport mechanism through the tunnel barriers between each FM and the QW. (Here $`v`$ is the Fermi velocity, and we put $`e=k_B=\mathrm{}=1`$.) A complete analytic solution of this problem is contained in Eqs. (2) and (15) to (17). In contrast to charge transport, we find that spin conduction occurs not only through electron transfer but also exchange. This exchange effectively gives rise to a modification of the boundary conditions at the end of the LL, e.g., for the left contact,
$$\stackrel{}{J}_R=(\mathrm{\Theta })\stackrel{}{J}_L+\stackrel{}{J}_{\mathrm{tunnel}},$$
(1)
where $`\stackrel{}{J}_{L/R}`$ is the left/right moving spin current into/out of the contact, and $`\stackrel{}{J}_{\mathrm{tunnel}}`$ represents the effect of electron transfer, see Eq. (15). The effect of exchange coupling is given by the one-parameter $`SO(3)`$ matrix $`(\mathrm{\Theta })=\mathrm{exp}(\mathrm{\Theta }\mathrm{\Gamma })`$, where $`\mathrm{\Gamma }_{\mu \nu }=_\lambda \widehat{m}_\lambda ϵ_{\lambda \mu \nu }`$, and $`\widehat{m}`$ is a unit vector in the direction of magnetization of the FM. Physically, $`\mathrm{\Theta }`$ represents the angle an incident spin in the LL precesses due to exchange interaction with the FM. Due to spin-charge separation in the LL, the exchange contribution is not suppressed by the orthogonality catastrophe affecting the tunneling current, and therefore dominates the physics in many situations. This enhancement of the exchange current does not occur in a Fermi liquid, and its observation would provide a direct experimental signature of electron fractionalization. In addition to the novel physics arising at the contact, we find that a long ballistic QW exhibits a bulk precession of the magnetization,
$$v_x\stackrel{}{M}+_t\stackrel{}{J}=b\stackrel{}{M}\times \stackrel{}{J},$$
(2)
where $`\stackrel{}{M}=\stackrel{}{J}_R+\stackrel{}{J}_L`$ and $`\stackrel{}{J}=\stackrel{}{J}_R\stackrel{}{J}_L`$ are the local magnetization and current in the QW, respectively. Eq. (2) leads to oscillations in the nonlinear current-voltage relation. Remarkably, the latter is a purely non-equilibrium effect that arises from a marginally irrelevant backscattering interaction in the LL. The detailed character of these oscillations is also influenced by interactions.
We now turn to the derivation of these results. In the incoherent limit, we may consider each contact separately, as an initial system composed of two decoupled pieces, $`H_0=H_{FM}+H_{QW}`$. The FM $`(x<0)`$, described by $`H_{FM}`$, is polarized along direction $`\widehat{m}`$, while the $`SU(2)`$ invariant QW ($`x>0`$) is described by $`H_{QW}`$. The $`SU(2)`$ invariance guarantees the existence of a continuity equation for spin density and current. At time $`t\mathrm{}`$, each half is assumed at equilibrium at its own chemical potential, $`\mu _{FM}`$ and $`\mu _{QW}`$, and with a spin chemical potential $`\stackrel{}{h}`$ (see below) in the QW. We are interested in the steady state achieved at $`t=0`$, long after the tunneling perturbation has been adiabatically turned on, $`H(t)=H_0+e^{\delta t}H^{}`$ ($`\delta 0^+`$). The calculation is non-trivial primarily due to its non-equilibrium nature: the system evolves according to $`H(t)`$ while the initial states are distributed according to $`\mathrm{exp}(\beta H_0)`$. Consider
$$H^{}=F^{}W\mathrm{\Psi }^{}+\mathrm{\Psi }^{}W^{}F^{},$$
(3)
where $`F`$ and $`\mathrm{\Psi }`$ are spin-$`1/2`$ Fermion annihilation operators at $`x=0`$ for the FM and the QW, respectively. Employing the projection operators $`\widehat{u}_s=(1\pm \widehat{m}\stackrel{}{\sigma })/2`$, the $`2\times 2`$ tunneling matrix $`W`$ takes the form $`W=_st_s\widehat{u}_s`$, with spin-dependent transmission coefficients $`t_s`$ for spin quantization axis parallel to $`\widehat{m}`$. The junction is then characterized by the conductance $`G=G_{}+G_{}`$ and the polarization $`P=(G_{}G_{})/G`$, where $`G_,=(e^2/h)|t_,|^2`$ are the spin conductances .
From Eq. (3) and the spin continuity equation, the tunneling spin current across the junction is $`\stackrel{}{J}=i\left(F^{}W\stackrel{}{\sigma }\mathrm{\Psi }^{}\mathrm{\Psi }^{}\stackrel{}{\sigma }W^{}F^{}\right)/2`$. By defining $`\stackrel{~}{H}_0=H_0+\mu _{FM}N_{FM}+\mu _{QW}N_{QW}+\stackrel{}{h}\stackrel{}{S}_{QW}`$, the standard perturbative result can be rewritten as
$`\stackrel{}{J}`$ $`=`$ $`\mathrm{Re}{\displaystyle \underset{\alpha \beta \gamma \lambda }{}}{\displaystyle _{\mathrm{}}^0}𝑑te^{\delta t}\left(W\stackrel{}{\sigma }\right)_{\alpha \beta }\left(U^{}(t)W^{}\right)_{\gamma \lambda }`$ (5)
$`\times [F_\alpha ^{}(0)\mathrm{\Psi }_\beta ^{}(0),\mathrm{\Psi }_\gamma ^{}(t)F_\lambda ^{}(t)]_{\stackrel{~}{H}_0}.`$
Thereby an intrinsically non-equilibrium expectation value is expressed in terms of an equilibrium average using the shifted Hamiltonian $`\stackrel{~}{H}_0`$, where the non-equilibrium nature of the problem is fully encoded in the time dependent unitary matrix $`U(t)=\mathrm{exp}[i(V+\stackrel{}{h}\stackrel{}{\sigma }/2)t]`$, with $`V=\mu _{QW}\mu _{FM}`$. A formula similar to Eq. (5) can easily be written down for the charge current, $`I=i(F^{}W\mathrm{\Psi }^{}\mathrm{\Psi }^{}W^{}F^{})`$. Thus both charge and spin current can be calculated using equilibrium correlation functions.
To proceed, we specify the Hamiltonians $`H_{FM}`$ and $`H_{QW}`$. For energies well below the electronic bandwidth $`D`$, the $`F`$ and $`\mathrm{\Psi }`$ equilibrium correlators are identical for $`H_0`$ and $`\stackrel{~}{H}_0`$, and moreover a non-interacting Fermi liquid model with constant density of states (DOS) applies to the leads. Because the lead couples to the QW only at $`x=0`$, the difference in DOS for majority and minority spin carriers can be absorbed in a spatial rescaling of the Fermi fields of the FM and a suitable redefinition of the transmission coefficients ($`t_s`$). Then
$$H_{FM}=i_{\mathrm{}}^0𝑑xf^{}\tau ^z_xf^{},$$
(6)
where the spinor $`f=f_{b\beta }`$ is indexed by $`b=(R,L)`$, describing right- and left-moving modes of the FM, and by $`\beta =(,)`$ for the spin, with the boundary condition $`f_R(0)=f_L(0)`$. Here the Pauli matrix $`\tau ^z`$ acts in the $`R/L`$ space. Putting $`F=f(0)`$, we see that $`F`$ has $`SU(2)`$ invariant correlation functions. The low-energy description of the QW is an interacting LL model,
$$H_{QW}=_0^{\mathrm{}}𝑑x\left\{i\psi ^{}v\tau ^z_x\psi ^{}+u\left(\psi ^{}\psi ^{}\right)^2\right\},$$
(7)
where $`\psi _R(0)=\psi _L(0)`$ and $`\mathrm{\Psi }=\psi (0)`$. Only the forward-scattering interaction $`u`$ is kept in Eq. (7). Alternatively, the exponent $`\alpha >0`$ for tunneling into the end of the LL ($`x=0`$) serves to measure the interaction strength .
The $`SU(2)`$ invariance of Eqs. (6) and (7) implies
$`[F_\alpha ^{}(0)\mathrm{\Psi }_\beta ^{}(0),\mathrm{\Psi }_\gamma ^{}(t)F_\lambda ^{}(t)]_{\stackrel{~}{H}_0}\theta (t)=\delta _{\alpha \lambda }\delta _{\beta \gamma }iC(t),`$
where $`C(t)`$ is the retarded Green’s function of the operator $`F_{}^{}\mathrm{\Psi }_{}`$ (the choice of spin components is arbitrary). From Eq. (5) and the corresponding expression for $`I`$, it is then straightforward to obtain
$`\stackrel{}{J}`$ $`=`$ $`{\displaystyle \frac{G}{2}}{\displaystyle \underset{s}{}}[(P\widehat{m}+s\widehat{h})\mathrm{Im}\stackrel{~}{C}(V+hs/2+i\delta )`$ (9)
$`Ps\widehat{m}\times \widehat{h}\mathrm{Re}\stackrel{~}{C}(V+hs/2+i\delta )],`$
$`I`$ $`=`$ $`G{\displaystyle \underset{s}{}}(1+Ps\widehat{m}\widehat{h})\mathrm{Im}\stackrel{~}{C}(V+hs/2+i\delta ),`$ (10)
where $`\stackrel{~}{C}(V)=𝑑tC(t)e^{iVt}`$. The terms involving
$`_\alpha (V,T)G\mathrm{Im}\stackrel{~}{C}(V+i\delta )`$ (11)
$`=GT(T/D)^\alpha \mathrm{sinh}(V/2T)\left|\mathrm{\Gamma }\left(1+{\displaystyle \frac{\alpha }{2}}+i{\displaystyle \frac{V}{2\pi T}}\right)\right|^2`$ (12)
have a simple interpretation in terms of tunneling via Fermi’s golden rule, as can be seen from the spectral representation of $`\stackrel{~}{C}`$. However, the appearance of $`\mathrm{Re}\stackrel{~}{C}`$ in Eq. (9) indicates the presence of a physical process other than tunneling. It can be shown that it corresponds to a virtual process in which an electron near the Fermi energy in the QW hops into a state of the FM (which could be far from the Fermi energy) and hops back, thereby generating an exchange coupling. We thus include it from the start, $`HH_0+H^{}+H^{\prime \prime }`$, with
$$H^{\prime \prime }=K\widehat{m}\mathrm{\Psi }^{}\stackrel{}{\sigma }\mathrm{\Psi }^{}/2.$$
(13)
Since the FM possesses a non-vanishing average magnetization, the spin operator in the FM may be replaced by this average to leading approximation.
It is helpful to view both $`H^{}`$ and $`H^{\prime \prime }`$ in a renormalization group (RG) framework, as perturbations to a decoupled fixed point described by $`H_0`$. Standard arguments give the scaling dimension of both $`t_s`$ and $`K`$, $`\mathrm{\Delta }_{t_s}=1+\alpha /2`$ and $`\mathrm{\Delta }_K=1`$. The scaling dimension $`\mathrm{\Delta }_K`$ is not renormalized due to spin-charge separation in the QW. A simple calculation gives the RG scaling equations
$`_{\mathrm{}}|t_s|^2(\mathrm{})=\alpha |t_s|^2,_{\mathrm{}}K(\mathrm{})=c\left(|t_{}|^2|t_{}|^2\right),`$
where $`\mathrm{}=\mathrm{ln}(D/E)`$, and $`c`$ is a non-universal constant. Following the RG flow from the ultraviolet cutoff $`D`$ down to energy $`E\mathrm{max}(T,V)D`$, we find $`|t_s|^2(E)=|t_s|^2(E/D)^\alpha `$, and
$`K(E)=K+\alpha ^1cGP[1(E/D)^\alpha ]|t_s|^2(E).`$
Therefore the tunneling spin current is much smaller than the exchange contribution. Neglecting the tunneling contribution completely, one still obtains a $`T`$-independent exchange spin current as $`T0`$.
This fact can be understood from a simple analogy to the Andreev current through a ballistic superconductor-normal-superconductor (SNS) junction . Let us consider a LL connected to two insulating FMs at $`x=0`$ and $`x=L`$, with $`\widehat{m}\times \widehat{m}^{}0`$. This is an equilibrium situation, which can be modeled using Eq. (7) for the LL and two copies of Eq. (13) for the contacts to the FMs. The exchange interaction operates entirely within the spin sector of the LL due to spin-charge separation. Since the charge sector is decoupled, we are free to consider it at the non-interacting point, $`u=0`$. Then the resulting fictitious charge boson and the physical spin boson can be combined and refermionized into a spin-ful Dirac fermion $`\eta `$. Choosing arbitrary quantization axes $`\widehat{m}=\widehat{x}`$ and $`\widehat{m}^{}=\mathrm{cos}(\theta )\widehat{x}+\mathrm{sin}(\theta )\widehat{y}`$, it is instructive to perform the particle-hole transformation $`\eta _{}\eta _{}^{}`$. This yields the Hamiltonian
$$H_\eta =iv_0^L𝑑x\eta ^{}\tau ^z_x\eta ^{}\mathrm{Re}[K\mathrm{\Delta }(0)+K^{}\mathrm{\Delta }(L)e^{i\theta }],$$
(14)
where $`\mathrm{\Delta }(x)=\eta _{}(x)\eta _{}(x)`$. Equation (14) describes a ballistic SNS junction, and for phase twist $`0<\theta <2\pi `$ between the superconductors, supports an equilibrium current due to Andreev reflection. Since the Andreev current is $`v\eta ^{}\tau ^z\eta ^{}`$, the original FM-LL-FM device indeed has a non-zero spin current $`J_z`$. The analogy to a SNS junction also demonstrates that this current does not rely upon the incoherence of the two contacts.
A more general perspective on the exchange coupling can be gained by viewing the low-energy physics entirely in terms of boundary operators and boundary conditions . For that purpose, we may make an arbitrary choice of short-scale physics, and let the exchange coupling act on right-movers slightly away from the junction. In this case, using the boundary condition $`\psi _L(0)=\psi _R(0)`$, the equations of motion for the spin currents can be integrated over the junction region to give the steady-state relation $`\stackrel{}{J}_R(0^+)=(\mathrm{\Theta })\stackrel{}{J}_L(0^+)`$ (the brackets denoting expectation values are omitted henceforth). The parameter $`\mathrm{\Theta }K/v`$ ultimately defines the “exchange coupling constant” of the low-energy theory. In principle, since the boundary exchange operator is exactly marginal, $`\mathrm{\Theta }`$ need not be small. Then the “bulk” spin currents are $`\stackrel{}{J}_L=\stackrel{}{J}_L(0^+)`$ and $`\stackrel{}{J}_R=\stackrel{}{J}_R(0^+)+\stackrel{}{J}_{\mathrm{tunnel}}`$, and we obtain Eq. (1), with the tunneling spin current
$$\stackrel{}{J}_{\mathrm{tunnel}}=\frac{1}{2}\underset{s}{}(P\widehat{m}+s\widehat{h})_\alpha (V+hs/2,T).$$
(15)
The term proportional to $`\mathrm{Re}\stackrel{~}{C}`$ in Eq. (9) has been dropped, as its physical effects are included via the $`SO(3)`$ rotation $``$. Since the magnetization far from the contact is $`\stackrel{}{M}=\chi \stackrel{}{h}`$ with the LL spin susceptibility $`\chi `$, one can then compute the spin current $`\stackrel{}{J}`$ for arbitrary exchange coupling $`\mathrm{\Theta }`$. We arrive at the general result
$$\stackrel{}{J}=S\chi \stackrel{}{h}+(1S)\stackrel{}{J}_{\mathrm{tunnel}},$$
(16)
where $`S=(1)/(+1)`$ is a real antisymmetric matrix. Similarly, the charge current is
$$I=\underset{s}{}(1+s\widehat{m}\widehat{h})_\alpha (V+hs/2,T).$$
(17)
From Eqs. (16) and (17), by exploiting spin and charge current conservation in order to obtain $`\mu _{QW}`$ and $`\stackrel{}{h}`$, one can then compute all transport properties in a given circuit for arbitrary parameters .
We now specialize to a FM-LL-FM device with identical contacts at $`T=0`$ and applied voltage $`V`$ within $`v/LVD`$. For algebraic simplicity, we require $`P^2(1+\alpha )^2`$. For a tunneling contact, one expects $`\mathrm{\Theta }1`$ , and Eq. (16) then yields
$`\stackrel{}{J}=(\mathrm{\Theta }\chi /2)\widehat{m}\times \stackrel{}{h}+\stackrel{}{J}_{\mathrm{tunnel}}.`$
Under these conditions, it is straightforward to obtain the $`\theta `$-dependent FM-LL-FM current-voltage relation,
$$I(\theta )=\frac{GV}{2}(V/D)^\alpha \left(1P^2\frac{\mathrm{tan}^2(\theta /2)}{\mathrm{tan}^2(\theta /2)+Y_\alpha (V)}\right),$$
(18)
where $`Y_\alpha (V)=1+(\mathrm{\Theta }\chi /2)^2(1+\alpha )^2(V/D)^{2\alpha }`$. For $`\alpha 0`$, this reproduces the result of Ref. . Notably, unless the magnetizations of the FMs are anti-parallel ($`\theta =\pi `$) or the exchange coupling vanishes ($`\mathrm{\Theta }=0`$), the spin accumulation effect, in which the current is reduced due to pile-up of spin in the QW, is strongly suppressed by the voltage dependence of $`Y_\alpha `$. Physically, this suppression is caused by the exchange coupling which is much more efficient in relaxing the injected spin polarization compared to the tunneling current.
Finally, we turn to backscattering electron-electron interactions of the form
$$H_b=b_0^L𝑑x\stackrel{}{J}_L\stackrel{}{J}_R.$$
(19)
For a carbon nanotube, with the lattice spacing $`a`$ and the tube radius $`R`$, one may estimate $`bae^2/R`$ . Since $`H_b`$ is marginally irrelevant in a single-channel QW, it is usually neglected in the LL model (7). Nevertheless, as is shown here, dynamical effects caused by (19) can be important. The equations of motion away from the contacts give Eq. (2) and $`v_x\stackrel{}{J}+_t\stackrel{}{M}=0`$. In the steady state, we have conserved spin current $`\stackrel{}{J}`$, and a bulk precession equation for $`\stackrel{}{M}`$. Since $`\stackrel{}{M}=\chi \stackrel{}{h}`$, the vector $`\widehat{h}`$ must then precess around $`\stackrel{}{J}`$.
To get sizeable consequences, detailed analysis shows that it is essential to have small exchange couplings. For simplicity, we put $`\mathrm{\Theta }=0`$ below. For the FM-LL-FM device,
$$\widehat{h}(0)\widehat{h}(L)=\mathrm{cos}(\mathrm{\Delta }\phi ),\mathrm{\Delta }\phi =bJL/v.$$
(20)
Since $`\mathrm{\Delta }\phi L`$, precession will then always be significant for a sufficiently long QW. The computation of $`\stackrel{}{h}(0)`$ and $`\stackrel{}{h}(L)`$ leads to the self-consistency equation
$$(1x^2)\mathrm{cos}\left(\frac{\pi x}{\mathrm{cos}(\theta /2)}\frac{V}{\mathrm{\Delta }V}(V/D)^\alpha \right)=\mathrm{sin}^2(\theta /2),$$
(21)
where $`\mathrm{\Delta }V=8\pi v/[GPbL\mathrm{cos}(\theta /2)]`$. Here solutions $`x=x_n`$ ($`n0`$) in the interval $`0x\mathrm{cos}(\theta /2)`$ have to be found. The current through the device is then
$$I_n(V)=\frac{GV}{2}(V/D)^\alpha (1P^2[1x_n^2(V)]).$$
(22)
The solution $`I_n`$ corresponding to $`n`$ full precession periods exists only for voltages $`V>V_nD[\mathrm{\Delta }V/D]^{1/(1+\alpha )}n^{1/(1+\alpha )}`$. Using typical parameters appropriate for a $`1\mu `$m long single-wall nanotube, one finds $`V_10.1`$ to 1 V. In general, the current-voltage relation could then be multi-valued, where in the regime $`V_n<V<V_{n+1}`$, the solution $`I_n(V)`$ is expected to be realized. Hence the current-voltage relation becomes oscillatory, with sawtooth-like oscillations. The observation of several periods could provide a direct and accurate measurement of $`\alpha `$ via the $`V_n`$.
To conclude, we have presented a general formalism for spin-dependent transport through interacting 1D conductors. An experimental check of the theory should be possible by measuring the current-voltage relation for a ferromagnet-nanotube-ferromagnet device. The approach is general enough to apply to numerous other problems. Several interesting extensions currently under investigation are the description of bulk-contacted wires, inclusion of the subband degree of freedom of single-wall nanotubes, LL to LL contacts, and the ballistic–diffusive crossover potentially relevant for multi-wall nanotubes.
We thank Gerrit Bauer for helpful discussions. Financial support was provided by the NSF CAREER program under Grant NSF-DMR-9985255, the NSF grant PHY-94-07194, and by the DFG under the Gerhard-Hess and the Heisenberg program.
$``$ Present address: Fakultät für Physik, Universität Freiburg, D-79104 Freiburg, Germany
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# CERN-TH/2000-092March 2000 RADIATIVE HIGGS-SECTOR CP VIOLATION IN THE MSSM
## Abstract
We briefly review the phenomenological implications of the minimal supersymmetric standard model (MSSM) with explicit radiative breaking of CP invariance in the Higgs sector for the LEP2 and Tevatron colliders.
It has recently been shown that the tree-level CP invariance of the MSSM Higgs potential can sizeably be broken at the one-loop level by large soft CP-violating trilinear couplings of the Higgs bosons to stop and sbottom squarks. Several recent studies have been devoted to analyze in more detail the effective Higgs potential of the MSSM with explicit radiative breaking of CP invariance. We shall briefly review the main phenomenological implications of this rather rich and very predictive theoretical framework of the MSSM for the LEP2 and Tevatron colliders.
In the $`\overline{\mathrm{MS}}`$ scheme, the one-loop CP-violating effective potential of the MSSM is given by
$$_V=_V^0+\frac{3}{32\pi ^2}\underset{q=t,b}{}\left[\underset{i=1,2}{}\stackrel{~}{m}_{q_i}^4\left(\mathrm{ln}\frac{\stackrel{~}{m}_{q_i}^2}{Q^2}\frac{3}{2}\right)\mathrm{\hspace{0.17em}2}\overline{m}_q^4\left(\mathrm{ln}\frac{\overline{m}_q^2}{Q^2}\frac{3}{2}\right)\right],$$
(1)
where $`_V^0`$ is the tree-level Lagrangian of the MSSM Higgs potential, and $`\overline{m}_q`$ and $`\stackrel{~}{m}_{q_i}`$ are the field-dependent quark and squark masses of the third generation. Here, we adopt the notation of Refs. . The minimization of the CP-violating effective potential differs from that in the CP-conserving case by the presence of a non-trivial CP-odd tadpole condition of the would-be CP-odd scalar $`a`$. In the $`\overline{\mathrm{MS}}`$ scheme, the CP-odd tadpole condition reads
$`T_A=v\mathrm{Im}(m_{12}^2e^{i\xi })`$ $`=`$ $`{\displaystyle \frac{3}{16\pi ^2}}{\displaystyle \underset{q=t,b}{}}{\displaystyle \frac{s_{2q}}{s_\beta }}\mathrm{Im}h_1^q\mathrm{\Delta }m_{\stackrel{~}{q}}^2B_0^{\mathrm{fin}}(0,m_{\stackrel{~}{q}_1}^2,m_{\stackrel{~}{q}_2}^2),`$ (2)
where $`t_\beta =s_\beta /c_\beta =v_2/v_1`$, $`s_{2q}=2\mathrm{sin}\theta _q\mathrm{cos}\theta _q`$, $`\mathrm{\Delta }m_{\stackrel{~}{q}}^2=m_{\stackrel{~}{q}_2}^2m_{\stackrel{~}{q}_1}^2`$, $`h_1^t=m_t\mu ^{}e^{i\delta _t}/(s_\beta v)`$, $`h_1^b=m_bA_b^{}e^{i\delta _b}/(c_\beta v)`$, $`\delta _{t(b)}=\mathrm{arg}[A_{t(b)}\mu ^{}t_\beta (1/t_\beta )]`$ and
$$B_0^{\mathrm{fin}}(0,m_1^2,m_2^2)=\mathrm{ln}\left(\frac{m_1m_2}{Q^2}\right)+\mathrm{\hspace{0.17em}1}+\frac{m_1^2+m_2^2}{m_1^2m_2^2}\mathrm{ln}\left(\frac{m_2}{m_1}\right).$$
(3)
Furthermore, $`\mu `$ and $`m_{12}^2`$ are respectively the SUSY and soft SUSY-breaking Higgs-mixing terms, $`A_{t,b}`$ are the soft SUSY-breaking Yukawa couplings, and $`\theta _q`$ is the mixing angle between the mass and weak squark eigenstates.
To one-loop order, CP violation is introduced into the MSSM Higgs potential through the complex parameters $`\mu `$ and $`A_{t,b}`$. The CP-violating terms are proportional to the rephasing invariant combination: $`\mathrm{Im}(m_{12}^2A_{t,b}\mu )`$. From this last expression, it is clear that only the relative phase between $`\mu A_{t,b}`$ and $`m_{12}^2`$ plays a role. Therefore, a good phase convention is to define $`m_{12}^2`$ to be real. This can always be achieved by a global U(1) rotation. As a result, the relative phase $`\xi `$ between the two Higgs vacuum expectation values $`v_1`$ and $`v_2`$ vanishes at the tree level. The phase choice $`\xi =0`$ can be preserved order by order in perturbation theory by a corresponding choice of the counter-term of $`\mathrm{Im}m_{12}^2`$, exactly as is given in Eq. (2). In fact, this perturbative resetting of the phase $`\xi `$ to zero is equivalent to the general requirement that within the effective-potential formalism, $`v_1`$, $`\mathrm{Re}v_2=|v_2|\mathrm{cos}\xi `$ and $`\mathrm{Im}v_2=|v_2|\mathrm{sin}\xi `$ do not receive finite radiative shifts in higher orders other than those due to Higgs wave-function renormalization. The latter approach is also consistent with the one followed in the CP-conserving studies.
It is known that the MSSM faces the difficulty of explaining naturally the apparent absence of electric dipole moments (EDMs) of the neutron and electron. Several suggestions have been made to suppress the SUSY contributions to electron and neutron EDMs, at a level just below their present experimental upper limits. Apart from the obvious choice of suppressing the new CP-violating phases of the theory to the $`10^3`$ level, a more phenomenologically appealing possibility is to make the first two generations of scalar fermions as heavy as few TeV, but keep the soft-breaking mass parameters of the third generation relatively small, e.g. 0.5–0.7 TeV. An interesting alternative is to arrange for partial cancellations among the different EDM contributions either at the short-distance level or the non-perturbative long-distance one.
In addition to the one-loop EDM contributions of the first two generations, one may have to worry that third-generation squarks do not induce observable effects on the electron and neutron EDMs through the three-gluon operator , through the effective coupling of the ‘CP-odd’ Higgs boson to the gauge bosons , and through two-loop gaugino/higgsino-mediated EDM graphs . For low-$`t_\beta `$ scenarios, the two-loop EDM contributions are found to be of the order of the experimental upper bounds. Therefore, it not very difficult to arrange the different two-loop EDM terms to partially cancel one another, and so reduce significantly their total size.
An immediate consequence of CP violation in the Higgs potential of the MSSM is the presence of mixing-mass terms between the CP-even and CP-odd Higgs fields. In the weak basis $`(\varphi _1,\varphi _2,a)`$, the neutral Higgs-boson mass matrix $`_N^2`$ takes on the form
$$_N^2=\left[\begin{array}{cc}_S^2& _{SP}^2\\ (_{SP}^2)^T& _P^2\end{array}\right],$$
(4)
where $`_S^2`$ and $`_P^2`$ describe the CP-conserving transitions between scalar and pseudoscalar particles, respectively, whereas $`_{SP}^2`$ describes CP-violating scalar-pseudoscalar transitions. The characteristic size of these CP-violating off-diagonal terms in the Higgs-boson mass matrix was found to be
$`M_{SP}^2`$ $``$ $`𝒪\left({\displaystyle \frac{m_t^4}{v^2}}{\displaystyle \frac{|\mu ||A_t|}{32\pi ^2M_{\mathrm{SUSY}}^2}}\right)\mathrm{sin}\varphi _{\mathrm{CP}}`$ (5)
$`\times (6,{\displaystyle \frac{|A_t|^2}{M_{\mathrm{SUSY}}^2}},{\displaystyle \frac{|\mu |^2}{\mathrm{tan}\beta M_{\mathrm{SUSY}}^2}},{\displaystyle \frac{\mathrm{sin}2\varphi _{\mathrm{CP}}}{\mathrm{sin}\varphi _{\mathrm{CP}}}}{\displaystyle \frac{|\mu ||A_t|}{M_{\mathrm{SUSY}}^2}}),`$
where the last bracket summarizes the relative sizes of the different contributions, and $`\varphi _{\mathrm{CP}}=\mathrm{arg}(A_t\mu )`$. As can be seen from Eq. (5), the CP-violating effects can become substantial if $`|\mu |`$ and $`|A_t|`$ are larger than the average of the stop masses, denoted as $`M_{\mathrm{SUSY}}`$. For example, the off-diagonal terms of the neutral Higgs-mass matrix may be of order $`(100\mathrm{GeV})^2`$, for $`|\mu ||A_t|\stackrel{<}{_{}}3M_{\mathrm{SUSY}}`$, and $`\varphi _{\mathrm{CP}}90^{}`$.
The main effect of Higgs-sector CP violation is the modification of the couplings of the Higgs bosons to fermions and the $`W`$ and $`Z`$ bosons, i.e. $`ffH_i`$, $`WWH_i`$, $`ZZH_i`$ and $`ZH_iH_j`$. The modified effective Lagrangians are given by
$`_{H\overline{f}f}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}H_i[{\displaystyle \frac{g_wm_d}{2M_Wc_\beta }}\overline{d}(O_{1i}is_\beta O_{3i}\gamma _5)d`$ (6)
$`+{\displaystyle \frac{g_wm_u}{2M_Ws_\beta }}\overline{u}(O_{2i}ic_\beta O_{3i}\gamma _5)u],`$
$`_{HVV}`$ $`=`$ $`g_wM_W{\displaystyle \underset{i=1}{\overset{3}{}}}(c_\beta O_{1i}+s_\beta O_{2i})(H_iW_\mu ^+W^{,\mu }`$ (7)
$`+{\displaystyle \frac{1}{2c_w^2}}H_iZ_\mu Z^\mu ),`$
$`_{HHZ}`$ $`=`$ $`{\displaystyle \frac{g_w}{4c_w}}{\displaystyle \underset{i,j=1}{\overset{3}{}}}\left[O_{3i}(c_\beta O_{2j}s_\beta O_{1j})O_{3j}(c_\beta O_{2i}s_\beta O_{1i})\right]`$ (8)
$`\times Z^\mu (H_i\underset{\mu }{\overset{}{}}H_j),`$
where $`c_w=M_W/M_Z`$, $`\underset{\mu }{\overset{}{}}\underset{\mu }{\overset{}{}}\underset{\mu }{\overset{}{}}`$, and $`O`$ is the orthogonal transformation matrix relating the weak with the mass Higgs-boson eigenstates.
Let us now discuss a representative example demonstrating the phenomenological consequences of Higgs-sector CP violation on the LEP2 and Tevatron colliders. We consider an intermediate value for $`\mathrm{tan}\beta =4`$, and a relatively light charged Higgs boson $`M_{H^+}=150`$ GeV, with $`M_{\mathrm{SUSY}}=0.5`$ TeV, $`A_t=A_b=1`$ TeV and $`\mu =2`$ TeV. In Fig. 1, we then find regions for which the lightest Higgs-boson mass $`M_{H_1}`$ is as small as 60–70 GeV for $`\mathrm{arg}(A_t)90^{}`$, and the $`H_1ZZ`$ coupling $`g_{H_1ZZ}`$, which is normalized to its SM value, is small enough for the $`H_1`$ boson to escape detection at the latest LEP2 run with $`\sqrt{s}=202`$ GeV. Moreover, the $`H_2`$ boson is too heavy to be detected through the $`H_2ZZ`$ channel. In addition, either the coupling $`H_1H_2Z`$, $`g_{H_1H_2Z}=g_{H_3ZZ}`$, is too small or $`H_2`$ is too heavy to allow Higgs detection in the $`H_1H_2Z`$ channel. An upgraded Tevatron machine has the potential capabilities to close most of such experimentally open windows. The results of the recent complete RG analysis of Ref. are in good qualitative agreement with earlier studies.
In conclusion, the MSSM with explicit radiative breaking of CP invariance in the Higgs sector constitutes a very rich and predictive theoretical framework, with interesting consequences on collider experiments, CP asymmetries in $`B`$-meson decays, and electroweak baryogenesis.
I wish to thank Marcela Carena, Darwin Chang, John Ellis, Wai-Yee Keung and Carlos Wagner for collaboration.
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# VanVleck Response Of A Two-Level System And Mesoscopic Orbital Magnetism Of Small Metals
## I Introduction
The significance of the few-level physics for small metal particles has long been realized,. At temperatures smaller than the mean level spacing at the Fermi level, $`T\mathrm{\Delta }`$, the spin susceptibility, for instance, is determined by the (spin-flip) transitions to the first unoccupied state for the particles that contain even number of electrons and is typically much smaller than the Pauli susceptibility. For the odd-electron particles, it is given by the Curie susceptibility. This is in contrast to the $`T\mathrm{\Delta }`$ (and the bulk) case where the leading term is given by the Pauli susceptibility. Moreover, due to the exponential activation factors associated with transitions to unoccupied states, one should expect extremely broad distribution of susceptibility values and, therefore, large variations from particle to particle.
The orbital magnetic response of small metal particles, on the other hand, does not easily reduce to a few-level problem. While it does turn out to be essentially a Fermi level property for $`T\mathrm{\Delta }`$ and in the bulk systems-, a priori it appears as a property of the entire Fermi sea since all the levels are perturbed by the magnetic field. The two competing contributions to orbital magnetism are the precession diamagnetism and the polarization (van Vleck) paramagnetism. While no rigorous description of the orbital magnetism exists for $`T\mathrm{\Delta }`$, it is nonetheless expected to be a Fermi level property as well due to the large Fermi sea cancellations between the diamagnetic and paramagnetic contributions. However, at the Fermi level the van Vleck part of the orbital response should be dominant and also the one that is very sensitive to variations of the energy level spacings. Consequently, we study a model wherein it is assumed that among all the ”virtual transitions” between the ground and the excited states of the Fermi sea that are responsible for the van Vleck response, the one that determines the particle susceptibility is a single-electron ”transition” from the last occupied to the first unoccupied state.
Imry has conjectured that, since the electrons in a metal particle should be considered as a canonical ensemble, one can describe the particle response in terms of the difference between the magnetic energies of the canonical and grand canonical ensemble $`F_H\mathrm{\Omega }_H`$, the latter being the case for a bulk system and relatively well understood. In the limiting case of $`T\mathrm{\Delta }`$, one should expect the difference to be due to a single electron level crossing the position of the chemical potential in the equivalent grand canonical ensemble. However, the formalism based on this approach works only in the limit of $`T\mathrm{\Delta }`$, as explained in Refs.,, in which case $`F\mathrm{\Omega }=\mathrm{\Delta }\delta N^2/2`$ where the particle number fluctuation in the equivalent grand canonical ensemble is actually much larger than one, $`\delta N^2T/\mathrm{\Delta }1`$. In this limit, the perturbation theory can be used. An improvement on this approach involves using the exact, non-perturbative level correlation function and, in fact, yields a saturation value of the magnetic energy as $`T`$ approaches $`\mathrm{\Delta }`$ from above.
In what follows, we first give an argument in support of the large Fermi sea cancellation between the diamagnetic and paramagnetic contributions to the magnetic response. We then evaluate the two-level van Vleck response for the Fermi level and first unoccupied level. This involves averaging of the inverse energy spacing using the Gaussian Orthogonal Ensemble (GOE) statistics for energy eigenvalues and calculation of the magnetic dipole moment. The latter can be estimated using a semiclassical argument but also rigorously evaluated by considering the magnetic dipole absorption. We compare our result with the expression for the orbital magnetic response obtained in the limit $`T>\mathrm{\Delta }`$. We restrict our analysis to 2D diffusive particles (disks, narrow strips, and rings) and use the units where $`c=\mathrm{}=1`$. We also neglect any sample-specific effects (fluctuations), except when specifically mentioned, and all the quantities in consideration are presumed disorder-averaged.
## II Fermi Sea Cancellation Of Van Vleck Paramagnetism And Precession Diamagnetism
The van-Vleck energy of a two-level system in the magnetic field $`𝐇=H_0\widehat{𝐳}`$ is
$$\delta ϵ_{vV}^{\left(1\right)}=\frac{\left|i\left|\widehat{M}_z\right|f\right|^2H_0^2}{\epsilon _i\epsilon _f}\frac{\left|\widehat{M}_{if}\right|^2H_0^2}{\epsilon _i\epsilon _f}$$
(1)
and the total van Vleck energy is given by
$$ϵ_{vV}^{\left(tot\right)}=\frac{1}{2}\underset{i,f}{}\left|\widehat{M}_{if}\right|^2H_0^2\frac{n_in_f}{\epsilon _i\epsilon _f}$$
(2)
where $`n_i=\theta \left(\epsilon _i\right)`$. In the case of a continuous spectrum, we can rewrite eq. (2), using $`dn_i/d\epsilon _i=\delta \left(\epsilon _i\right)`$, as
$$ϵ_{vV}^{\left(tot\right)}=\frac{1}{2}\upsilon H_0^2\underset{f}{}\left|\widehat{M}_{0f}\right|^2=\frac{1}{2}\upsilon 0\left|M_z^2\right|0H_0^2=\frac{\upsilon e^2v_F^20\left|r_{}^2\right|0H_0^2}{16}$$
(3)
where $`\upsilon \upsilon \left(0\right)`$ is the mean level density at the Fermi level, $`\upsilon \left(\epsilon \right)`$ is the level density at energy $`\epsilon `$. We have also used the fact that the magnetic moment at the Fermi level can be evaluated semiclassically, that is
$$𝐌=\frac{e}{2}𝐫\times 𝐯_F$$
(4)
where $`𝐫`$ is the classical position vector of an electron and $`𝐯_F`$ is the Fermi velocity. In evaluating the coefficient in eq. (3), we averaged over the angle between $`𝐫`$ and $`𝐯_F`$.
The single-level diamagnetic energy diamagnetic energy is given by
$$\delta ϵ_{diam}^{\left(1\right)}=\frac{e^2i\left|r_{}^2\right|iH_0^2}{8m}$$
(5)
and the total diamagnetic energy of the Fermi sea is given by
$$\delta ϵ_{diam}^{\left(tot\right)}=\underset{i}{}\frac{e^2i\left|r_{}^2\right|iH_0^2}{8m}n_i$$
(6)
It is reasonable to conjecture that the disorder-averaged value $`i\left|r_{}^2\right|i`$ is $`i`$ -independent which yields, upon converting to integration for a continuous spectrum,
$$\delta ϵ_{diam}^{\left(tot\right)}=\frac{\upsilon e^2v_F^20\left|r_{}^2\right|0H_0^2}{16}$$
(7)
and is the same as eq. (3), with the opposite sign. This confirms the Fermi sea cancellation between the van Vleck paramagnetism and precession diamagnetism<sup>*</sup><sup>*</sup>*This argument can be carried over, with minor modifications, to 3D.<sup>,</sup>For a disk of radius $`R`$, $`0\left|r_{}^2\right|0`$ can be evaluated as the area average and equals to $`R^2/2`$.. However, this derivation does not account for the quantum effects beyond the existence of the Fermi sea. Consequently, it is understood that the cancellation is not exact and that a Fermi level contribution is not accounted for in the present approximation. The nature of this contribution should depend on whether the chemical potential or the number of particles is fixed. In the former case, one expects a Landau response, as in bulk systems, while in the latter we make an ansatz of a two-level van Vleck response which involves the last occupied (Fermi) level and the first unoccupied level.
The situation is more complex in a strictly discrete level case where we will give only an order-of-magnitude argument. The first principles estimate of $`\left|\widehat{M}_{if}\right|^2`$ in the diffusive regime is based on the idea first proposed by Shapoval and, later, applied by Gor’kov and Eliashberg to small metal particles. Namely, we use the semiclassical approach to write
$$\left|\widehat{M}_{if}\right|^2=\frac{1}{\pi \upsilon }\overline{_0^{\mathrm{}}𝑑\tau \mathrm{exp}\left[\left(\epsilon _i\epsilon _f\right)\tau \right]M\left(t\right)M\left(t+\tau \right)}$$
(8)
where the bar denotes averaging over all classical trajectories. The correlation time scale $`\tau `$ for $`M\left(t\right)`$ is given by
$$\tau \frac{\mathrm{}^2}{D}$$
(9)
where $`\mathrm{}=v_F\tau `$ is the electron mean-free-path and $`D=v_F^2`$ $`\tau /2`$ is the diffusion coefficient. This is because the directions of $`𝐯_F`$ are uncorrelated after such time. Also, for $`\epsilon _f\epsilon _i>\tau ^1`$ the exponential term becomes oscillatory. The scale of $`𝐫`$ is the relevant sample dimension $`a`$ and we find,
$$\left|\widehat{M}_{if}\right|^2\frac{1}{\pi \upsilon }\overline{_0^{\mathrm{}}𝑑\tau M\left(t\right)M\left(t+\tau \right)}\frac{e^2\mathrm{}^2v_F^2a^2}{\upsilon D}\frac{e^2Da^2}{\upsilon }\mu _B^2\left(\epsilon _F\tau \right)$$
(10)
where $`\mu _B`$ is the Bohr magneton and $`\epsilon _F`$ is the Fermi energy. Consequently, the order of magnitude value of the two-level van Vleck response is obtained from eq. (1) as
$$\delta ϵ_{vV}^{\left(1\right)}\frac{\left|\widehat{M}_{if}\right|^2H_0^2}{\mathrm{\Delta }}\frac{\mu _B^2H_0^2}{\mathrm{\Delta }}\left(\epsilon _F\tau \right)\left|\chi _L\right|H_0^2A\left(\epsilon _F\tau \right)$$
(11)
where $`\chi _L`$ is the Landau susceptibility and $`Aa^2`$ is the sample area. The total van Vleck energy can be estimated by multiplying $`\delta ϵ_{vV}^{\left(1\right)}`$ by $`\left(\tau \mathrm{\Delta }\right)^1`$, which determines the range of the integral in eq. (8), and we find
$$ϵ_{vV}^{\left(tot\right)}\frac{\mu _B^2H_0^2}{\mathrm{\Delta }}\frac{\epsilon _F}{\mathrm{\Delta }}\left|\chi _L\right|H_0^2A\frac{\epsilon _F}{\mathrm{\Delta }}$$
(12)
Since $`0\left|r_{}^2\right|0a^2`$ , this is in qualitative agreement with eq. (3). Notice also that eq. (5) yields
$$\delta ϵ_{diam}^{\left(1\right)}\frac{\mu _B^2H_0^2}{\mathrm{\Delta }}\left|\chi _L\right|H_0^2A$$
(13)
for a single-level contribution to the precession diamagnetism.
## III Van Vleck Paramagnetism In The Two-Level Model
Turning again to eq. (1), where now $`\epsilon _i`$ is the energy of the last occupied state (Fermi level) and $`\epsilon _f`$ is the energy of the first unoccupied state, we find, averaging with the Wigner-Dyson distribution,
$$\delta ϵ=s\left|\widehat{M}_{if}\right|^2H_0^2\frac{\pi }{2\mathrm{\Delta }}_0^{\mathrm{}}\mathrm{exp}\left(\frac{\pi x^2}{4}\right)𝑑x=\frac{\pi \upsilon }{2}\left|\widehat{M}_{if}\right|^2H_0^2$$
(14)
Here $`s`$ is the level degeneracy ($`s=2`$, on the account of spin) and $`\upsilon =s\mathrm{\Delta }^1`$. We have already estimated the matrix element using the semiclassical argument. However, a precise derivation of $`\left|\widehat{M}_{if}\right|^2`$ in the diffusive regime can be done by means of evaluation of the low-frequency magneto-dipole absorption in the field $`𝐇=H_0\mathrm{exp}\left(i\omega t\right)\widehat{𝐳}`$. In the case of the continuous energy spectrum (for instance, when the level broadening $`\gamma `$ is larger than $`\mathrm{\Delta }`$ ), the quantum-mechanical expression for the absorption should yield, up to small corrections, the classical value. This is in complete analogy with the electric-dipole absorption (barring screening effects for the latter), which is discussed in detail in Ref.. Since $`\left|\widehat{M}_{if}\right|^2`$ enters into the quantum-mechanical expression (see below), it can be extracted by equating with the classical value of the absorption.
The classical absorption is readily evaluated according to
$$Q_{class}=\frac{1}{2}\omega Im\left\{M_z^{}H_0\right\}$$
(15)
where
$`M_z^{\left(disk\right)}`$ $`=`$ $`{\displaystyle \frac{AH_0}{4\pi }}\left(1{\displaystyle \frac{2}{\varkappa R}}{\displaystyle \frac{J_1\left(\varkappa R\right)}{J_0\left(\varkappa R\right)}}\right){\displaystyle \frac{AH_0}{32\pi }}\left(\varkappa R\right)^2\text{}A=\pi R^2`$ (16)
$`M_z^{\left(strip\right)}`$ $`=`$ $`{\displaystyle \frac{AH_0}{8\pi }}\left(1{\displaystyle \frac{\mathrm{tan}\left(\varkappa L_x/2\right)}{\varkappa L_x/2}}\right){\displaystyle \frac{AH_0}{96\pi }}\left(\varkappa L_x\right)^2\text{}A=L_xL_y`$ (17)
for a disk of radius $`R`$ and for a narrow metal strip, such that $`L_xL_y`$, respectively. Here
$$\varkappa =\frac{1+i}{\delta }\text{}\delta =\frac{1}{\sqrt{2\pi \sigma \omega }}$$
(18)
and $`\sigma `$ is the Boltzmann conductivity. It is assumed that the frequency is such that $`\delta R,L_x`$. Combining eqs. (15)-(18), we obtain
$`Q_{class}^{\left(disk\right)}`$ $`=`$ $`{\displaystyle \frac{\omega ^2H_0^2R^2A\sigma }{16}}`$ (19)
$`Q_{class}^{\left(strip\right)}`$ $`=`$ $`{\displaystyle \frac{\omega ^2H_0^2L_x^2A\sigma }{48}}`$ (20)
for the absorption, respectively, in the disk and the strip.
The quantum absorption, for the continuous spectrum, can be evaluated (in complete analogy with the electric-dipole absorption) by means of the Fermi golden rule and we find
$$Q_{cont}=\frac{\pi }{2}\omega ^2\upsilon ^2H_0^2\left|\widehat{M}_{if}\right|^2\frac{\upsilon \left(0\right)\upsilon \left(\omega \right)}{\upsilon ^2}\frac{\pi }{2}\omega ^2\upsilon ^2H_0^2\left|\widehat{M}_{if}\right|^2$$
(21)
where the small quantum corrections of order $`\mathrm{\Delta }^2/\gamma ^2`$ (or $`\mathrm{\Delta }^2/\omega ^2`$ if $`\omega >\gamma `$) are neglected. For the electric-dipole absorption,, the quantity corresponding to $`\widehat{M}_{if}`$ is the electric dipole matrix element $`\widehat{P}_{if}`$ . The latter can be evaluated from first principles in the diffusive approximation and the classical and quantum results can be evaluated independently and are equal in the considered limit. On physical grounds, we can assume that this is also true for the magnetic-dipole absorption. Consequently, we equate the r.h.s. of eqs. (19), (20) with that of eq. (21) and, with the use of $`\sigma =e^2\upsilon D/A`$, we obtain
$`{\displaystyle \frac{\pi }{2}}\upsilon \left|\widehat{M}_{if}^{\left(disk\right)}\right|^2`$ $`=`$ $`{\displaystyle \frac{e^2DR^2}{16}}`$ (22)
$`{\displaystyle \frac{\pi }{2}}\upsilon \left|\widehat{M}_{if}^{\left(strip\right)}\right|^2`$ $`=`$ $`{\displaystyle \frac{e^2DL_x^2}{48}}`$ (23)
for the disk and the strip respectively.
Substituting thus found value of $`\left|\widehat{M}_{if}\right|^2`$ in eq. (14), we find (for $`s=2`$) the following result for the van Vleck energy:
$`\delta ϵ^{\left(disk\right)}`$ $`=`$ $`{\displaystyle \frac{1}{16}}e^2DR^2H_0^2`$ (24)
$`\delta ϵ^{\left(strip\right)}`$ $`=`$ $`{\displaystyle \frac{1}{48}}e^2DL_x^2H_0^2`$ (25)
for the disk and the strip respectively. Eqs. (24) and (25) should be compared with the magnetic part of the energy obtained for $`T>\mathrm{\Delta }`$ in the so called ”mixed approximation” wherein the exact, non-perturbative level correlation function is used yet the large particle-number fluctuation in the equivalent grand canonical ensemble is assumed also, the latter being true only for $`T\mathrm{\Delta }`$,
$$\delta ϵ_>=\frac{1}{2\pi }\tau _H^1$$
(26)
The procedure for the evaluation of $`\tau _H^1`$ using the gauge where the vector potential is tangential to the surface is described in Ref. and we find
$`\delta ϵ_>^{\left(disk\right)}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}e^2DR^2H_0^2`$ (27)
$`\delta ϵ_>^{\left(strip\right)}`$ $`=`$ $`{\displaystyle \frac{1}{6\pi }}e^2DL_x^2H_0^2`$ (28)
for the disk and the strip respectively.
We also mention the Aharonov-Bohm response of a narrow ring threaded by the flux $`\varphi `$, where the derivation is particularly simple. Namely, the van Vleck energy in this case is given by
$$\delta ϵ_{if}=\frac{\left|\widehat{v}_{if}\right|^2}{\epsilon _i\epsilon _f}\left(\frac{e\varphi }{2\pi R}\right)^2$$
(29)
where $`\left|\widehat{v}_{if}\right|^2`$ is evaluated by considering absorption in the alternating flux $`\varphi \mathrm{exp}\left(i\omega t\right)`$ which generates the electric field according to the Lenz’s law that, in turn, produces the Boltzmann current density proportional to $`v`$. Consequently, we find
$$\left|\widehat{v}_{if}\right|^2=\frac{D}{4\pi \upsilon }$$
(30)
and, averaging over the level spacing,
$$\delta ϵ^{\left(ring\right)}\frac{D}{8}\left(\frac{e\varphi }{\pi R}\right)^2$$
(31)
The latter has the same parametric dependence as
$$\delta ϵ_>^{\left(ring\right)}=\frac{1}{2\pi }\tau _H^1=\frac{D}{2\pi }\left(\frac{e\varphi }{\pi R}\right)^2$$
(32)
which is the limiting value of energy for $`T>\mathrm{\Delta }`$.
## IV Discussion
The main result of this paper is that the parametric dependence of eqs. (24), (25) and (27), (28) is the same (with closely matching numerical coefficients). Consequently, the orbital magnetic response of a level-quantized metal particle can be satisfactorily explained in terms of the two-level van Vleck response. This conclusion also holds for other level distributions characteristic of disordered (chaotic) systems, such as the Gaussian Unitary Ensemble (only the numerical coefficient will be different).
Whereas the two-level picture might be also valid for the Poisson distribution, which is the case for classically integrable systems, the mean response needs to be examined more carefully because level ”bunching,” characteristic of this distribution, implies that for any magnitude of the field there will be a significant number of particles for which the perturbation theory no longer applies.
Another question which we hope to address in a future work is the sample-specific (fluctuation) effects that are expected to be quite large in the limit considered here. We point out that, in addition to the fluctuation of the level spacing and the magnetic moment, the details of the cancellation between the van Vleck paramagnetism and Landau diamagnetism should differ from particle to particle leading to variations in the number of levels contributing to the total response.
## V Acknowledgment
I wish to thank Bernard Goodman for helpful discussions. This work was not supported by any funding agency.
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# Virtual Color Superconductivity–Status and Perspective aafootnote aTalk given at the “Symposium on the Frontiers of Physics at Millennium”, Beijing, China, Oct. 8-10, 1999.
## 1 Introduction
QCD at finite temperature and zero density condition is relatively well understood compared the finite density situation. The technical reasons are well known. Physically, quantum decoherence increases as temperature gets high making the behaviour of the system essentially classical. On the other hand, density effects do not reduce quantum coherence. One has to deal with the quantum effects. A brief review is given here about the current status of our understanding of the role played by the color superconducting phase in the strong interaction ground state close to zero density and in hadron structure.
## 2 High Density Limit
The ground state of QCD at high enough density is most likely to be in a color superconducting phase. This is because quarks on a high Fermi surface carry high momentum and produce large momentum transfer in scattering, so that the two quark interaction kernel in color triplet channel can be approximated well by an attractive one gluon exchange term. The attractive interaction generates the BCS instability that creates a gap of order $`\mathrm{\Delta }_{BCS}e^{c/g}`$ with $`c`$ a constant. The properties of the color superconducting phase are actively studied in the recent literature base on various approaches like the instanton model, the one gluon exchange approximation improved by renormalization group treatment and the contact 4-fermion interaction model.
No matter what models are adopted for the discussion, the fact that color superconducting phase should be the ground state of the high density quark matter is quite clear according to well established BCS theory for superconductivity in non-relativistic condensed matter systems. A gap of order of a few 100 MeV is guaranteed at not so high density since the natural scale of the problem is of order 1 GeV and the coupling between quarks is of order 1. Otherwise, fine tuning is required.
The general consensus on the properties of the color superconducting phase of quark matter at asymptotically high density and zero temperature is that the ground state is in a color flavor locked state involving up, down and strange quarks. As the density is lowered, the color flavor locked state ceases to be the true ground state of the system, the true ground state is in a color superconducting phase involves only the up and down quarks. Such a color superconducting phase is again disfavoured as the density is further lowered to some critical density $`\rho _c`$. The color superconducting phase stops to manifest at a macroscopic scale when $`\rho \rho _c`$. The current estimate of $`\rho _c`$ is so high that the physical effects of the color superconducting phase can only be possible to manifest in a neutron star in the conventional understanding of the problem based mainly on the physical picture gained in the study of non-relativistic many body systems.
Since there are already reviews on this subject$`^\mathrm{?}`$ at high density, I shall mainly concentrate on a brief review of my own efforts in the past 10 years (1989-1999), which actually started the study of vacuum color superconductivity. I however looked at the problem with a different perspective.
## 3 The Possible Role of the Color Superconducting Phase at Low or Zero Density
What is the fate of the color superconducting phase at low density then? One thing is quite clear that the attraction between quark pairs in a color triplet state still exist. It is evidenced by the existence of baryons which are made up of three quarks. One of the reasons besides confinement for the color superconducting phase to disappear at low density is due to the competition from the quark–antiquark pairing. Since as the density is lowered, the antiquark excitation due to interaction becomes less suppressed. A sufficient number of the interaction generated antiquarks begin to pair with the quarks of the system to condense in the ground state, which cause the well know spontaneous breaking of the chiral symmetry.
In order to study the phase and the neighbourhood of the strong interaction vacuum state, the present framework of finite density theory based on non-relativistic space-time is found to be inadequate. Inconsistencies in the logic flow of physical argumentation and unphysical solution occur when one goes beyond the quasi-particle picture. Thus not only the kinematics of individual particles has to be relativistic, but also the covariance related to the relativistic space-time transformation for the whole system at the quantum field level has to be implemented in a way that is consistent to quantum mechanics.
The results of the efforts include: 1) the model building and studying for vector$`^{\mathrm{?},\mathrm{?}}`$ and scalar$`^\mathrm{?}`$ diquark condensation and many of the technical development$`^\mathrm{?}`$ 2) the introduction of an 8–component real representation for fermions$`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$, which is not identical to the conventional 4-component representation at finite density$`^{\mathrm{?},\mathrm{?}}`$. The advantage of using the 8–component theory to study the vacuum properties is discussed. It also enable me to isolate some of the differences near zero density between my work and the later commonly adopted ones$`^\mathrm{?}`$ which contain a continuous mixing between the color superconducting phase and the chiral symmetry breaking phase in the parameter space where no metastable phase is found. The difference originates from a different coupling between quarks in the diquarks. While the “reality condition” adopted here requires that quark and mirror quark<sup>b</sup><sup>b</sup>bIt is identical to a quark at zero density. are correlated to condense, other approaches couple quark and mirror antiquark for the correlated diquarks$`^\mathrm{?}`$, which can naturally be mixed with the chiral symmetry breaking phase. Such a difference disappear at high density. 3) the possibility that the effects of a metastable phase can manifest in local observables and propagators $`^{\mathrm{?},\mathrm{?}}`$ 4) the development of a consistent local finite density theory$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$, which combines all of the above considerations. The effects of statistical blocking in the chiral symmetry breaking phase and the spontaneous CP violation in the color superconducting phase was discovered as a result$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. Also discussed is the dynamics of the statistical gauge field in various phases. It is believed that the statistical blocking effects could provide a mechanism$`^\mathrm{?}`$ to prevent the dissolution of a nucleon inside a nucleus$`^\mathrm{?}`$, which is in direct contradiction to the experimental fact that a nucleon inside a nucleus keeps most of its identity.
Is what are talked about above correspond to reality? To answer it, some of the empirical knowledge about a nucleon, which serves as a carrier of the information about the virtual phase of the vacuum state, are analyzed.
The chiral properties of a nucleon is studied$`^{\mathrm{?},\mathrm{?}}`$, to search for a vector type of virtual color superconducting phase. There are some hints. More detailed works are needed. The recent discovery of a large anapole moment for a nucleon$`^\mathrm{?}`$ that can not be accounted for by conventional electroweak theory$`^\mathrm{?}`$ may also be a hint for a vector type virtual color superconducting phase since the order parameter here breaks parity which can induce an axial type of coupling for a photon to a nucleon$`^\mathrm{?}`$. The electromagnetic properties of a nucleon is studied next base on the mechanism of partial breaking of the electromagnetic gauge symmetry$`^\mathrm{?}`$ in a color superconductor. The Gerasimov-Drell-Hearn sum rule, which is currently unable to be saturated by experimental data, is a possible observable for both scalar and vector type of color superconducting phase$`^\mathrm{?}`$. A collection deep inelastic scattering data for a lepton on a nucleon is analyzed, the results, especially the measured polarized structure function for a nucleon at small Bjorken x, indicates the existence of a virtual superconducting phase$`^\mathrm{?}`$. Other evidences, which is less model independent, is also considered$`^\mathrm{?}`$. In the near future, at least two possibile developements could be foreseen: 1) the nucleon $`\mathrm{\Sigma }`$–term problem, which is made worse by the analysis $`^{\mathrm{?},\mathrm{?}}`$ of most recent data on pion nucleon scattering, could be resolved $`^\mathrm{?}`$ by using the local theory developed here $`^\mathrm{?}`$ 2) type of the CP/T violation data obtained in the KTeV Collaboration should be analyzed$`^\mathrm{?}`$ in more details to search$`^\mathrm{?}`$ for possible anomalies.
## Acknowledgements
This work was support by National Natural Science Foundation of China. I would like to thank Prof. A. W. Thomas and the CSSM of Univeristy of Adelaide for hospitality during which this write up is done.
## References
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# FROM BLACK HOLES TO POMERON: Tensor Glueball and Pomeron Intercept at Strong Coupling11footnote 1 Talk presented by R. C. Brower at ISMD99: QCD and Multiuparticle Production. This work was supported in part by the Department of Energy under Contract No. DE-FG02/91ER40688 and DE-FG02-91ER40676
## 1 Introduction
The Maldacena conjecture and its further extensions allow us to compute quantities in a strongly coupled gauge theory from its dual gravity description. In particular, Witten has pointed out if we compactify the 4-dimensional conformal super Yang Mills (SYM) to 3 dimensions using anti-periodic boundary conditions on the fermions, then we break supersymmetry and conformal invariance and obtain a theory that has interesting mass scales. This approach has been used to calculate a discrete mass spectrum for $`\stackrel{~}{0}^{++}`$ states associated with $`Tr[F^2]`$ at strong coupling by solving the dilaton’s wave equation in the corresponding gravity description. Although the theory at strong coupling is really not pure Yang-Mills, since it has additional fields, some rough agreement was claimed with the pattern of glueball masses.
Here we report on the calculation of the discrete modes for the perturbations of the gravitational metric. A complete description for all discrete fluctuations has also been carried out, both for $`QCD_3`$ and $`QCD_4`$. For simplicity, we shall discuss here mostly $`QCD_3`$. For $`QCD_4`$, from the mass of the $`2^{++}`$ state and the calculated QCD string tension, we obtain a strong coupling expansion for the Pomeron intercept: $`\alpha _P(0)=20(1/g^2N)`$. In this approach, the Pomeron corresponds to a “massive graviton”. Other results will be reported elsewhere.
## 2 AdS/CFT Duality at Finite $`\beta `$
Let us review briefly the proposal for getting a 3-d Yang-Mills theory dual to supergravity. One begins by considering Type IIB supergravity in Euclidean 10-dimensional spacetime with the topology $`M_5\times S^5`$. The Maldacena conjecture asserts that IIB superstring theory on $`AdS^5\times S^5`$ is dual to the $`𝒩=4`$ SYM conformal field theory on the boundary of the $`AdS`$ space. The metric of this spacetime is
$$ds^2/R_{ads}^2=r^2(d\tau ^2+dx_1^2+dx_2^2+dx_3^2)+\frac{dr^2}{r^2}+d\mathrm{\Omega }_5^2,$$
where the radius of the $`AdS`$ spacetime is given through $`R_{AdS}^4=g_sN\alpha ^2`$ ($`g_s`$ is the string coupling and $`l_s`$ is the string length, $`l_s^2=\alpha ^{}`$). The Euclidean time is $`\tau =ix_0`$. To break conformal invariance, following ref. , we place the system at a nonzero temperature described by a periodic Euclidean time $`\tau =\tau +\beta `$, $`\beta =2\pi R_0`$. The metric correspondingly changes, for small enough $`R_0`$, to the non-extremal black hole metric in $`AdS`$ space. For large black hole temperatures, the stable phase of the metric corresponds to a black hole with radius large compared to the $`AdS`$ curvature scale. To see the physics of discrete modes, we may take the limit of going close to the horizon, whereby the metric reduces to that of the black 3-brane. This metric is, (where $`f(r)=r^2\frac{1}{r^2}`$, and we have scaled out all dimensionful quantities),
$$ds^2=fd\tau ^2+f^1dr^2+r^2(dx_1^1+dx_2^2+dx_3^2)+d\mathrm{\Omega }_5^2,$$
(1)
On the gauge theory side, we would have a $`𝒩=4`$ susy theory corresponding to the $`AdS`$ spacetime, but with the $`S^1`$ compactification with antiperiodic boundary conditions for the fermions, supersymmetry is broken and massless scalars are expected to acquire quantum corrections. Consequently from the view point of a 3-d theory, the compactification radius acts as an UV cut-off. Before the compactification the 4-d theory was conformal, and was characterized by a dimensionless effective coupling $`(g_{YM}^{(4)})^2Ng_sN`$. After the compactification the theory is not conformal, and the radius of the compact circle provides a length scale. Let this radius be $`R_0`$. Then a naive dimensional reduction from 4-d Yang-Mills to 3-d Yang-Mills, would give an effective coupling in the 3-d theory equal to $`(g_{YM}^{(3)})^2N=(g_{YM}^{(4)})^2N/(2\pi R_0)`$. The 3-d YM coupling has the units of mass. If the dimensionless coupling $`(g_{YM}^{(4)})^2N`$ is much less than unity, then the length scale associated with this mass is larger than the radius of compactification, and we may expect the 3-d theory to be a dimensionally reduced version of the 4-d theory.
Unfortunately the dual supergravity description only applies at $`(g_{YM}^{(4)})^2N>>1`$, so that the higher Kaluza-Klein modes of the $`S^1`$ compactification have lower energy than the mass scale set by the 3-d coupling. Thus we do not really have a 3-d gauge theory with a finite number of additional fields. One may nevertheless expect that some general properties of the dimensionally reduced theory might survive the strong coupling limit. Moreover, we expect that the pattern of spin splittings might be a good place to look for similarities. In keeping with earlier work, we ignore the Kaluza-Klein modes of the $`S^1`$ and restrict ourselves to modes that are singlets of the $`SO(6)`$, since non-singlets under the $`S^1`$ and the $`SO(6)`$ can have no counterparts in a dimensionally reduced $`QCD_3`$.
## 3 Wave Equations
We wish to consider fluctuations of the metric of the form,
$$g_{\mu \nu }=\overline{g}_{\mu \nu }+h_{\mu \nu }(x),$$
(2)
leading to the linear Einstein equation,
$$h_{\mu \nu ;\lambda }{}_{}{}^{\lambda }+h_\lambda ^\lambda {}_{;\mu \nu }{}^{}h_{\mu \lambda ;\nu }{}_{}{}^{\lambda }h_{\nu \lambda ;\mu }{}_{}{}^{\lambda }8h_{\mu \nu }=0.$$
Our perturbations will have the form
$$h_{\mu \nu }=ϵ_{\mu \nu }(r)e^{mx_3}$$
(3)
where we have chosen to use $`x_3`$ as a Euclidean time direction to define the glueball masses of the 3-d gauge theory. We fix the gauge to $`h_{3\mu }=0`$.
¿From the above ansatz and the metric, we see that we have an $`SO(2)`$ rotational symmetry in the $`x_1x_2`$ space, and we can classify our perturbations with respect spin.
Spin-2: There are two linearly independent perturbations which form the spin-2 representation of $`SO(2)`$: $`h_{12}=h_{21}=q_T(r)e^{mx_3},h_{11}=h_{22}=q_T(r)e^{mx_3}`$ with all other components zero. The Einstein equations give,
$$(r^2\frac{1}{r^2})q_T^{\prime \prime }+(r+\frac{3}{r^3})q_T^{}+(\frac{m^2}{r^2}4\frac{4}{r^4})q_T=0.$$
Defining $`\varphi _T(r)=q_T(r)/r^2`$, this is the same equation as that satisfied by the dilaton (with constant value on the $`S^5`$).
Spin-1: The Einstein equation for the ansatz, $`h_{i\tau }=h_{\tau i}=q_V(r)e^{mx_3},i=1,2`$ gives
$$(r^2\frac{1}{r^2})q_V^{\prime \prime }+(r\frac{1}{r^3})q_V^{}+(\frac{m^2}{r^2}4+\frac{4}{r^4})q_V=0.$$
(4)
Spin-0: Based on the symmetries we choose an ansatz where the nonzero components of the perturbation are
$`h_{11}`$ $`=`$ $`h_{22}=q_1(r)e^{mx_3}`$
$`h_{\tau \tau }`$ $`=`$ $`2q_1(r){\displaystyle \frac{f(r)}{r^2}}e^{mx_3}+q_2(r)e^{mx_3}`$
$`h_{rr}`$ $`=`$ $`q_3(r)e^{mx_3}`$
where $`f(r)`$ is defined above in the metric. The field equation for $`q_3q_S(r)`$, is
$$p_2(r)q_S^{\prime \prime }(r)+p_1(r)q_S^{}(r)+p_0(r)q_S(r)=0,$$
(5)
where $`p_2(r)=r^2(r^41)^2[3(r^41)+m^2r^2]`$, $`p_1(r)=r(r^41)[3(r^41)(5r^4+3)+m^2r^2(7r^4+5)]`$ and $`p_0(r)=9(r^41)^3+2m^2r^2(3+2r^4+3r^8)+m^4r^4(r^41)`$.
## 4 Numerical Solution
To calculate the discrete spectrum for our three equation, one must apply the correct boundary conditions at $`r=1`$ and $`r=\mathrm{}`$. The result is a Sturm-Liouville problem for the propagation of gravitational fluctuations in a “wave guide”.
Table 1. Glueball Excitation Spectrum
Using this shooting method we have computed the the first 10 states given in Table 4. The spin-2 equation is equivalent to the dilaton equation , so the excellent agreement with earlier values validates our method. We used a standard Mathematica routine with boundaries taken to be $`x=r^21=ϵ`$ and $`1/x=ϵ`$ reducing $`ϵ`$ gradually to $`ϵ=10^6`$. Note that since all our eigenfunctions must be even in $`r`$ with nodes spacing in $`x=r^21`$ of $`O(m^2)`$, the variable $`1/x`$ is a natural way to measure the distance to the boundary at infinity. For both boundaries, the values of $`ϵ`$ was varied to demonstrate that they were near enough to $`r=1,`$ and $`\mathrm{}`$ so as not to substantially effect the answer.
As one sees in the accompanying figure, they match very accurately with the leading order WKB approximation. Simple variational forms also lead to very accurate upper bounds for the ground state ($`n=0`$) masses.
## 5 Strong coupling Expansion for Pomeron Intercept
Our current exercise has been extended to 4-d QCD using a scheme involving the finite temperature version of $`AdS^7\times S^4`$. As has been suggested elsewhere, one goal is to find that background metric that has the phenomenologically best strong coupling limit. This should provide an optimal starting point for approaching the continuum weak coupling regime. Here, we shall report briefly the key constraint provided by the Pomeron intercept.
The Pomeron is the leading Regge trajectory passing through the lightest glueball state with $`J^{PC}=2^{++}`$. In a linear approximation, it can be parameterized by
$$\alpha _P(t)=2+\alpha _P^{}(tm_T^2),$$
(6)
where we can use the strong coupling estimate for the lightest tensor mass<sup>2</sup><sup>2</sup>2 We have adopted the normalization in the $`AdS`$-black hole metric to simplify the coefficients, e.g., for $`AdS^7`$, $`\overline{g}_{\tau \tau }=r^2r^4`$. This corresponds to fixing the “thermal-radius” $`R_1=1/3`$ so that $`\beta =2\pi R_1=2\pi /3`$.,
$$m_T[9.86+0(\frac{1}{g^2N})]\beta ^1.$$
(7)
Moreover if we make the standard assumption that the closed string tension is twice that between two static quark sources, we also have a strong coupling expression for the Pomeron slope,
$$\alpha _P^{}[\frac{27}{32\pi g^2N}+0(\frac{1}{g^4N^2})]\beta ^2.$$
(8)
Putting these together, we obtain a strong coupling expansion for the Pomeron intercept,
$$\alpha _P(0)20.66(\frac{4\pi }{g^2N})+0(\frac{1}{g^4N^2}).$$
(9)
Turning this argument around, we can estimate a crossover value between the strong and weak coupling regimes by fixing $`\alpha _P(0)1.2`$ at its phenomenological value. In fact this yields for $`QCD_4`$ at $`N=3`$ a reasonable value for $`\alpha _{strong}=g^2/4\pi =0.176`$ for the crossover. Much more experience with this new approach to strong coupling must be gained before such numerology can be taken seriously. However, similar crude argument have proven to be a useful guide in the crossover regime of lattice QCD. One might even follow the general strategy used in the lattice cut-off formulations. Postpone the difficult question of analytically solving the QCD string to find the true UV fixed point. Instead work at a fixed but physically reasonable cut-off scale (or bare coupling) to calculate the spectrum. If one is near enough to the fixed point, mass ratios should be reliable. After all, the real benefit of a weak/strong duality is to use each method in the domain where it provides the natural language. On the other hand, clearly from a fundamental point of view, finding analytical tools to understand the renormalized trajectory and prove asymptotic scaling within the context of the gauge invariant QCD string would also be a major achievement — an achievement that presumably would include a proof of confinement itself.
Results on these computations will be reported in a future publication.
Acknowledgments: We would like to acknowledge useful conversations with R. Jaffe, A. Jevicki, D. Lowe, J. M. Maldacena, H. Ooguri, and others.
## References
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# Graviton Propagators, Brane Bending and Bending of Light in Theories with Quasi-Localized Gravity
## Abstract:
We derive the graviton propagator on the brane for theories with quasi-localized gravity. In these models the ordinary 4D graviton is replaced by a resonance in the spectrum of massive Kaluza-Klein modes, which can decay into the extra dimension. We find that the effects of the extra polarization in the massive graviton propagator is exactly cancelled by the bending of the brane due to the matter sources, up to small corrections proportional to the width of the resonance. Thus at intermediate scales the classic predictions of Einstein’s gravity are reproduced in these models to arbitrary precision.
preprint: hep-th/0003020
Following the work of Randall and Sundrum (RS) there has been considerable interest in the phenomenon of localization of gravity (for previous relevant work see ). RS found a solution to the five dimensional Einstein equations in a background created by a single positive tension 3-brane and a negative bulk cosmological constant which reproduces the effects of four-dimensional gravity on the brane without the need to compactify the fifth dimension. The reason for this is that the fluctuations of the 4D metric are described by an ordinary quantum mechanical Schrödinger equation, where the shape of the QM potential resembles a volcano. This potential supports exactly one bound state with zero energy, which can be identified with the 4D massless graviton, since the wave functions of the massive continuum Kaluza-Klein (KK) modes are suppressed at the brane due to the tunneling through the potential barrier. Thus the effects of the KK modes are negligible for small energies, and at large distances ordinary 4D gravity is reproduced.
Recently Gregory, Rubakov and Sibiryakov (GRS) found a brane model in which 4D gravity is reproduced only at intermediate scales, since at very small scales there are the same power-law corrections as in the RS model, while at very large scales the gravitational theory on the brane is again modified .<sup>2</sup><sup>2</sup>2A similar proposal has been made in . It has been explained, that the reason behind this phenomenon (dubbed “quasi-localization of gravity”) is that the zero-mode of the QM system of these theories becomes unstable and the exact bound state is replaced by a resonance of zero mass in the continuum spectrum (see also ). As long as the lifetime of this resonance is large, there is a large region of intermediate scales where an effective 4D Newton potential is reproduced. Once we get to large enough distances, the resonance will decay and therefore the gravitational potential will be corrected.
It has been however suggested in , that even though the correct Newton potential is reproduced at intermediate scales, one does not reproduce the results of ordinary 4D Einstein gravity. The argument given in for this is that in these models the resonant mode is a collective effect of very light (but $`m0`$) KK modes. It is, however, well-known that the $`m0`$ limit of a massive graviton propagator does not reproduce the massless graviton propagator, due to the fact that the number of polarizations of the two fields do not match . Based on this discontinuity in the graviton propagator as $`m0`$, it was argued in that the predictions of theories with quasi-localization of gravity would significantly differ from those of ordinary gravity. In particular the bending of light was predicted to be $`\frac{3}{4}`$ of the value in general relativity.
In this letter we show that this is in fact not the case. We find that up to additional small corrections (which vanish in the limit in which the width of the resonance $`\mathrm{\Delta }m0`$) the results of ordinary 4D gravity are reproduced on the brane at intermediate scales as in the RS scenario.
The issue is that the 5D massless graviton propagator is given by
$$G_5(x,x^{})_{\mu \nu \rho \sigma }=\frac{d^5p}{(2\pi )^5}\frac{e^{ip(xx^{})}}{p^2}\left(\frac{1}{2}g_{\mu \rho }g_{\nu \sigma }+\frac{1}{2}g_{\mu \sigma }g_{\nu \rho }\frac{1}{3}g_{\mu \nu }g_{\rho \sigma }\right),$$
(1)
while the 4D massless propagator is
$$G_4(x,x^{})_{\mu \nu \rho \sigma }=\frac{d^4p}{(2\pi )^4}\frac{e^{ip(xx^{})}}{p^2}\left(\frac{1}{2}g_{\mu \rho }g_{\nu \sigma }+\frac{1}{2}g_{\mu \sigma }g_{\nu \rho }\frac{1}{2}g_{\mu \nu }g_{\rho \sigma }\right).$$
(2)
The difference in the tensor structure of these two propagators is caused by the presence of an extra polarization state, a 4D scalar field, contained in the 5D graviton propagator. Thus we see that the situation in theories with localization and quasi-localization of gravity is not very different: In the case of localized gravity the scalar field must be included in the effective 4D theory, while in the case of quasi-localized gravity the scalar is eaten by the massive graviton modes and appears as an additional graviton polarization. One has to explain in both cases why the extra polarization state does not eventually contribute to the propagator on the brane. Intuitively this is the case for the following reason: in brane models the stress tensor for a source on a brane has a vanishing 55 component; thus we might naïvely expect that one degree of freedom, namely the scalar field, decouples from the sources on the brane and therefore the ordinary massless 4D graviton propagator should be reproduced. It has been shown in two beautiful papers by Garriga and Tanaka , and by Giddings, Katz and Randall that this is in fact the case. They have shown that once a source term on the brane is introduced the brane itself is bent in a frame in which gravitational fluctuations are small. The effect of the bending is such that it exactly compensates for the effect of the extra 4D scalar field contained in the 5D graviton propagator, and thus the ordinary 4D massless graviton propagator is reproduced in the RS model. This hints that a similar situation may occur for theories with quasi-localized gravity, since the massive propagator of the KK modes making up the resonance exactly the same tensor structure as the massless graviton plus the scalar. We will show that the effect of the brane bending will again cancel the effects of the extra polarization in the massive propagator, up to corrections depending on the width of the resonance, which can be made arbitrarily small by adjusting the parameters of the theory. Thus there is no discontinuity in these models as $`\mathrm{\Delta }m0`$ and the results of ordinary 4D Einstein theory are reproduced at intermediate scales up to small corrections.
We now turn to the problem of constructing the propagators in a general brane world described by the metric
$$ds^2=dy^2+e^{A(y)}\eta _{\mu \nu }dx^\mu dx^\nu .$$
(3)
We will take $`A(y)`$ to approach the Randall-Sundrum form for $`|y|y_0`$:
$$A(y)2k|y|,$$
(4)
while for $`|y|y_0`$ the metric becomes flat:
$$A(y)\mathrm{constant}.$$
(5)
For instance, this is achieved in the GRS model by simply patching anti-de Sitter space to flat space at some point $`y=y_0`$:
$$A(y)=\{\begin{array}{cc}2k|y|\hfill & |y|y_0\hfill \\ 2ky_0\hfill & |y|y_0.\hfill \end{array}$$
(6)
However, more generally one can imagine backgrounds that smoothly interpolate between $`AdS_5`$ and flat space . Since the metric approaches the RS metric for $`|y|y_0`$, there is a brane located at $`y=0`$ on which the matter fields will live.
We now study gravitational fluctuations in the background (3). It is possible to choose a gauge for the gravitational fluctuations so that they have the form
$$ds^2=dy^2+e^{A(y)}(\eta _{\mu \nu }+h_{\mu \nu })dx^\mu dx^\nu ,$$
(7)
where, in addition, we impose the transverse-traceless conditions
$$^\nu h_{\mu \nu }=h_\mu ^\mu h=0.$$
(8)
In this gauge, the transverse-traceless fluctuations satisfy the simple equation
$$\left(e^A\mathrm{}^{(4)}+_y^22A^{}_y\right)h_{\mu \nu }=0,$$
(9)
which is nothing but the scalar wave equation for each of the components $`h_{\mu \nu }`$ in the background (3). At the brane, we have the usual matching condition:
$$_yh_{\mu \nu }|_{y=0}=0.$$
(10)
In the GRS model, there is an additional matching condition at the point $`y_0`$, where $`AdS_5`$ is patched onto flat space:
$$_yh_{\mu \nu }|_{y=y_0+}=_yh_{\mu \nu }|_{y=y_0}.$$
(11)
Until now, the only matter sources in the theory were those that were needed to produce the non-trivial background (3).<sup>3</sup><sup>3</sup>3The details of these background sources is not important for our present discussion. We only need assume that the sources vary linearly with a variation of the metric (see for a discussion of this). Now suppose we have additional sources corresponding to matter on the brane. The question is how this modifies the equations for the fluctuations. It turns out the answer, as described in the papers by Garriga and Tanaka , and by Giddings, Katz and Randall , is rather subtle but crucial for our analysis. We closely follow the former paper in what follows.
The point is that in the original coordinate system $`(x^\mu ,y)`$, the presence of the source on the brane actually bends the brane so that is no longer situated at $`y=0`$. In order to investigate this effect, it is useful to introduce a new coordinate system $`(\overline{x}^\mu ,\overline{y})`$ defined so that the brane is at $`\overline{y}=0`$. We require that the metric in the new coordinate system also has the form (7), and this means that the two coordinate systems must be related via
$$\overline{y}=y+\widehat{\xi }^y(x),\overline{x}^\mu =x^\mu +\widehat{\xi }^\mu (x)^\mu \widehat{\xi }^y(x)^y𝑑y^{}e^{A(y^{})}.$$
(12)
Notice that the functions $`\widehat{\xi }^y(x)`$ and $`\widehat{\xi }^\mu (x)`$ do not depend on $`y`$. The relation between the metric fluctuations in the two coordinates systems is then
$$h_{\mu \nu }=\overline{h}_{\mu \nu }+_\mu \widehat{\xi }_\nu +_\nu \widehat{\xi }_\mu 2_\mu _\nu \widehat{\xi }^y^y𝑑y^{}e^{A(y^{})}\eta _{\mu \nu }\widehat{\xi }^y_yA.$$
(13)
At the moment, we have not specified the function $`\widehat{\xi }^y(x)`$ and $`\widehat{\xi }^\mu (x)`$; however, they will be determined self-consistently, as we shall see below.
The additional sources will modify the matching condition at the brane (10). In the new coordinate system, where the brane is at $`\overline{y}=0`$, the Israel equations at the brane give,
$$_{\overline{y}}\overline{h}_{\mu \nu }|_{\overline{y}=0+}=\kappa ^2\left(S_{\mu \nu }\frac{1}{3}\eta _{\mu \nu }S\right),$$
(14)
where $`S_{\mu \nu }`$ is given in terms of the matter stress tensor via
$$T_{\mu \nu }^{\mathrm{brane}}=S_{\mu \nu }(x)\delta (\overline{y}),T_{\overline{y}\overline{y}}^{\mathrm{brane}}=T_{\mu \overline{y}}^{\mathrm{brane}}=0.$$
(15)
Using the relation between the metric fluctuations (13) we can now write down the boundary condition (14) in the original coordinates:<sup>4</sup><sup>4</sup>4Notice that to leading order in the source we can specify the following condition at $`y=0`$.
$$_yh_{\mu \nu }|_{y=0+}=\kappa ^2\left(S_{\mu \nu }\frac{1}{3}\eta _{\mu \nu }S\right)2_\mu _\nu \widehat{\xi }^y.$$
(16)
In the GRS model there is also the matching condition at $`y_0`$ (11). It is easy to show that this condition is simply
$$_{\overline{y}}\overline{h}_{\mu \nu }|_{\overline{y}=\overline{y}_0+}=_{\overline{y}}\overline{h}_{\mu \nu }|_{\overline{y}=\overline{y}_0},$$
(17)
and so in the original coordinates (11) remains unchanged.
The position of the brane in the original coordinate system is $`y=\widehat{\xi }^y(x)`$ and so the condition (16) includes, via $`\widehat{\xi }^y(x)`$, an effect from the bending of the brane. We now can now combine the equation of motion (9) and the boundary condition (16) to give
$$\left(e^A\mathrm{}^{(4)}+_y^22A^{}_y\right)h_{\mu \nu }=2\kappa ^2\mathrm{\Sigma }_{\mu \nu }\delta (y),$$
(18)
where we have defined the effective source
$$\mathrm{\Sigma }_{\mu \nu }=S_{\mu \nu }\frac{1}{3}\eta _{\mu \nu }S+\frac{2}{\kappa ^2}_\mu _\nu \widehat{\xi }^y.$$
(19)
In order to determine $`\widehat{\xi }^y`$, we now impose the fact that $`h_{\mu \nu }`$ is traceless: $`h=0`$. The implies that the right-hand side of (18) itself must be traceless and so
$$\mathrm{}^{(4)}\widehat{\xi }^y=\frac{\kappa ^2}{6}S,$$
(20)
with solution
$$\widehat{\xi }^y(x)=\frac{\kappa ^2}{6}d^4x^{}\mathrm{\Delta }_4(x,x^{})S(x^{}),$$
(21)
where $`\mathrm{\Delta }_4(x,x^{})`$ is the massless scalar Green’s function for 4D Minkowski space. If we now define the five-dimensional Green’s function
$$\left(e^A\mathrm{}^{(4)}+_y^22A^{}_y\right)\mathrm{\Delta }_5(x,y;x^{},y^{})=\delta ^{(4)}(xx^{})\delta (yy^{}),$$
(22)
then the fluctuation due to the source on the brane can be written
$$h_{\mu \nu }(x,y)=2\kappa ^2d^4x^{}\mathrm{\Delta }_5(x,y;x^{},0)\mathrm{\Sigma }_{\mu \nu }(x^{}).$$
(23)
The final thing that we need to do is to transform the fluctuation back into the coordinate system $`(\overline{x}^\mu ,\overline{y})`$ in which the brane is straight. In order to have a simple final expression for the fluctuation on the brane, we can use the additional freedom in the transformation (12), present in $`\widehat{\xi }^\mu (x)`$, to set,
$$\widehat{\xi }_\mu (x)=_\mu \left(\frac{1}{2k}\widehat{\xi }^y(x)2d^4x^{}\mathrm{\Delta }_5(x,0;x^{},0)\widehat{\xi }^y(x^{})\right).$$
(24)
Then from (13) and (23) we obtain our final result for the fluctuation evaluated on the brane:
$$\overline{h}_{\mu \nu }(x,0)=2\kappa ^2d^4x^{}\left\{\mathrm{\Delta }_5(x,0;x^{},0)\left(S_{\mu \nu }(x^{})\frac{1}{3}\eta _{\mu \nu }S(x^{})\right)\frac{k}{6}\mathrm{\Delta }_4(x,x^{})\eta _{\mu \nu }S(x^{})\right\},$$
(25)
where $`k=A^{}(0)/2`$. Note that in a flat background the last term in (25) would be absent, and the theory is simply that of five-dimensional gravity, as expected.<sup>5</sup><sup>5</sup>5We thank John Terning for raising this issue. The first term is exactly what one would have naïvely expected, while the second term arises from the effect of the bending of the brane .
Notice that the result (25) is written in terms of the scalar Green’s function $`\mathrm{\Delta }_5(x,y;x^{},y^{})`$ for the background (3) which allows us to make contact with the auxiliary quantum mechanical system of . First of all, we relate the coordinate $`z`$ to $`y`$ via
$$\frac{dz}{dy}=e^{A(y)/2}.$$
(26)
The wavefunctions $`\psi _m(z)`$, which satisfy the Schrödinger equation
$$\frac{d^2\psi _m}{dz^2}+\left(\frac{9}{16}A^2\frac{3}{4}A^{\prime \prime }\right)\psi _m=m^2\psi _m,$$
(27)
where $`{}_{}{}^{}d/dz`$, and then give the Green’s function that we need in (25). It is expressed as a sum over the bound-states and an integral over the continuum of KK modes:
$$\mathrm{\Delta }_5(x,0;x^{},0)=\frac{d^4p}{(2\pi )^4}e^{ip(xx^{})}\left\{\underset{m}{}\frac{\psi _m(0)^2}{p^2+m^2}+𝑑m\frac{\psi _m(0)^2}{p^2+m^2}\right\}.$$
(28)
The Schrödinger equation (27) always admits a zero energy state
$$\widehat{\psi }_0(z)=N_0e^{3A(z)/4},$$
(29)
which potentially describes the 4D graviton. In the RS model
$$\widehat{\psi }_0(z)=\frac{k^{1/2}}{(k|z|+1)^{3/2}},$$
(30)
is normalizable and so
$$\mathrm{\Delta }_5(x,0;x^{},0)=k\mathrm{\Delta }_4(x,x^{})+\mathrm{},$$
(31)
where the ellipsis represents the contribution from the KK modes. The contribution of the massless graviton fluctuation in (25) is then,
$$\overline{h}_{\mu \nu }(x,0)=2\kappa ^2kd^4x^{}\mathrm{\Delta }_4(x,x^{})\left(S_{\mu \nu }(x^{})\frac{1}{3}\eta _{\mu \nu }S(x^{})\right)+\mathrm{}.$$
(32)
However, this is not the usual propagator of a massless 4D graviton. Fortunately, (32) does not include the effect of the brane bending term in (25), which provides an additional contribution that precisely has the effect of changing $`\frac{1}{3}\frac{1}{2}`$ in (32) and so yields the usual massless 4D graviton propagator. So the brane bending effect is crucial even in the original RS model .
In a quasi-localized gravity model, the transverse space is asymptotically flat, i.e. $`A(y)`$ constant for $`|y|y_0`$. In this case, the state $`\widehat{\psi }_0(z)`$ is not normalizable and there are only contributions to the Green’s function from continuum modes. However, in we argued that if the scale $`y_0`$ is sufficiently large then $`\widehat{\psi }_0(z)`$ appears as a sharp resonance at the bottom of the continuum at $`m=0`$. In this case, we have approximately, for small $`m`$,
$$\psi _m(0)^2=\frac{𝒜}{m^2+\mathrm{\Delta }m^2}+\mathrm{}.$$
(33)
The height of the resonance is given by
$$\frac{𝒜}{\mathrm{\Delta }m^2}=\widehat{\psi }_0(0)^2=e^{3(A(0)A(\mathrm{}))/2},$$
(34)
while the width, to leading order, can be determined by the fact that as $`y_0\mathrm{}`$, the model should reduce to the RS model where $`\widehat{\psi }_0(z)`$ is normalizable. So as $`\mathrm{\Delta }m0`$, (33) should approximate $`k\delta (m)`$ giving
$$\frac{𝒜}{\mathrm{\Delta }m}=\frac{2k}{\pi }.$$
(35)
In the GRS model (34) and (35) give
$$\mathrm{\Delta }m=\frac{2k}{\pi }e^{3ky_0}.$$
(36)
For small, but non-zero $`\mathrm{\Delta }m`$, i.e. large $`y_0k^1`$, the resonance can mimic the effect of the bound-state in the RS model. For $`|xx^{}|\mathrm{\Delta }m^1`$, we can approximate the effect of the resonance by a delta function and hence naïvely we would expect a contribution to the graviton propagator as in (32). However, just as in the RS model itself we have to include the effect of the brane bending. This provides an additional contribution which precisely has the effect of changing $`\frac{1}{3}\frac{1}{2}`$ in (32) and yields the usual massless 4D graviton propagator . So the effect of brane bending in the case of quasi-localized gravity on the brane is exactly the same as for the RS model. We recover the normal 4D graviton propagator in the regime where we can ignore the contribution from the rest of the continuum, $`|xx^{}|k^1`$, and when we can approximate the resonance by a delta function, $`|xx^{}|\mathrm{\Delta }m^1`$. The effect of the finite width will give a calculable correction to the graviton propagator that can be determined by integrating over the shape of the resonance, giving
$$\begin{array}{cc}\hfill \overline{h}_{\mu \nu }(x,0)=& 2\kappa ^2kd^4x^{}\{\mathrm{\Delta }_4(x,x^{})(S_{\mu \nu }(x^{})\frac{1}{2}\eta _{\mu \nu }S(x^{}))\hfill \\ & +\mathrm{\Delta }m\stackrel{~}{\mathrm{\Delta }}_4(x,x^{})(S_{\mu \nu }(x^{})\frac{1}{3}\eta _{\mu \nu }S(x^{}))\},\hfill \end{array}$$
(37)
where
$$\stackrel{~}{\mathrm{\Delta }}_4(x,x^{})=\frac{d^4p}{(2\pi )^4}\frac{e^{ip(xx^{})}}{p^2(p+\mathrm{\Delta }m)}.$$
(38)
We have studied theories in which gravity is quasi-localized to a brane. In these theories there is a massless graviton resonance which eventually decays into the bulk, with the result of altering the very long distance behavior of gravity. It has been suggested that there are phenomenological difficulties with such models as a result of the non-decoupling of massive graviton polarizations in the massless limit. In this letter we have shown that in fact such difficulties are absent. The reason is that the bending of the brane exactly compensates for the effects of the extra polarization in the massive graviton propagator (just as it compensates for the effect of the massless scalar in the case of the RS theory). Thus the graviton propagator at intermediate distances will be equal to the massless propagator of the Einstein theory (up to corrections that can be made arbitrarily small by making the width of the resonance small). Therefore, all classic predictions of general relativity including the bending of light around the Sun and the precession rate of the orbit of Mercury are correctly reproduced in these theories.
Theories with quasi-localized gravity open exciting possibilities for phenomenology. Because gravity is modified at both very short and very long distances in such models, there are phenomenological consequences for both particle physics at high energies and cosmology at large distances. These consequences merit further investigation.
We would like to thank Tanmoy Bhattacharya and John Terning for several useful conversations, and Gia Dvali, Ruth Gregory and Valery Rubakov for comments on the manuscript and for sharing their views on these theories with us. C.C. is an Oppenheimer Fellow at the Los Alamos National Laboratory. C.C., J.E. and T.J.H. are supported by the US Department of Energy under contract W-7405-ENG-36.
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# BOTTOM QUARK PHYSICS PAST, PRESENT, FUTURE aafootnote aTalk given at “Symposium on Probing Luminous and Dark Matter, honoring Adrian Melissinos”, Rochester, October, 1999.
## 1 Introduction
Let me start by reminding you what’s going on at all the major High Energy Physics laboratories around the world. At CERN, the LEP program is winding down, and the LHC (Large Hadron Collider) is the Lab’s future. 7 TeV protons on 7 TeV protons, a center of mass energy of 14 TeV. Four large detectors are planned. Two, ATLAS and CMS, will study high $`p__T`$ physics, searching for Higgs, SUSY, etc. One, ALICE, will collide high $`Z`$ nuclei (when protons aren’t being collided), and study the quark-gluon plasma. And one, LHC-B, will study bottom quark physics. It is a sobering thought that a 14 TeV accelerator will be used to study a 5 GeV object, 3 orders of magnitude down the energy scale. (But one should not forget that the Tevatron is used to study kaon physics, again 3 orders of magnitude down the energy scale.)
At DESY, the main facility is HERA, an electron-proton collider, with 800 GeV protons on 30 GeV electrons. The two principal detectors, H1 and ZEUS, study these collisions, investigating deep inelastic scattering over a kinematic range far broader than heretofore. But the proton beam will also be used, on a fixed target (wires in the fringe of the beam) for bottom quark physics, in the HERA-B experiment.
At KEK, in Japan, TRISTAN, an $`e^+e^{}`$ collider operating at a center-of-mass energy of 60 GeV has been shut down, and replaced by an asymmetric $`e^+e^{}`$ collider, 8 GeV on 3.5 GeV, a center of mass energy of 10 GeV, to do bottom quark physics, with the Belle experiment.
At SLAC, the SLC (SLAC Linear Collider), $`e^+e^{}`$ collisions at center-of-mass energies around 90 GeV, has been shut down, and replaced by PEP- II, an asymmetric $`e^+e^{}`$ collider, 9 GeV on 3 GeV, a center of mass energy of 10 GeV, to do bottom quark physics with the BaBar experiment. Thus, the study of the $`Z^0`$, a 90 GeV object, is giving way to the study of the $`b`$ quark, a 5 GeV object.
At Fermilab, the main facility is the Tevatron, which collides 1 TeV protons against 1 TeV antiprotons, for a center-of-mass energy of 2 TeV. There are two general purpose detectors, operated by two large collaborations, CDF and D$`\mathrm{}`$. The primary goal of the running recently completed was the discovery of the top quark. Goals for the next running period (Run II) include precise measurements of top quark and $`W`$ boson masses, and searches for “new physics” – Higgs, SUSY, etc. But CDF has had an active program in bottom quark physics, and foresees an expanded program in Run II. A displaced vertex trigger is being implemented, in part to strengthen the $`b`$ physics program. D$`\mathrm{}`$ has done little $`b`$ physics so far, lacking a magnetic field in the central tracking volume. They are remedying this for Run II, and anticipate an active $`b`$ physics program. And serious consideration is being given to a third detector, B-TeV, which would be a dedicated bottom quark experiment.
Finally, Cornell’s Laboratory for Nuclear Studies, with a symmetric $`e^+e^{}`$ collider (CESR), has been doing bottom quark physics for two decades.
So, bottom quark physics must be interesting, because all the major labs have it as part of their program. Why is bottom quark physics so interesting? (The cynic might argue that the labs are into bottom quark physics because it’s affordable. There is perhaps some truth in this. But it doesn’t explain LHC-B. It doesn’t explain the interest in bottom quark physics within CDF, nor SLAC’s preference for studying a 5 GeV object over a 90 GeV object.) Why is bottom quark physics interesting? A primary goal of my talk will be to answer that question for you.
Bottom quark physics can be conveniently divided into three eras;
* The Early Days – 1977-88, further divided into Discovery – 1977-80, and Roughing out the Qualitative Features – 1980-88
* Beginnings of Precision Measurements and Rare Decay Studies – 1989-98
* The ‘Factory’ Era – 1999-??
In Section 2, I’ll discuss the early days.
In Section 3, I’ll point out a change in objective that took place around 1990, and give a brief review of the flavor sector of the Standard Model.
Then, in Sections 4, 5, and 6, I’ll discuss three of the “hot topics” in $`b`$ physics today: determination of $`|V_{ub}/V_{cb}|`$, rare hadronic $`B`$ decays, and the radiative penguin decay $`bs\gamma `$.
## 2 The Early Days
### 2.1 Discovery – 1977-80
The $`b`$ quark was discovered in its hidden form (“hidden beauty”, “covered bottom”) at Fermilab, in 1977, by Leon Lederman and collaborators. They measured the mass distribution of dimuon pairs from collisions of 400 GeV protons on a nuclear fixed target, and observed a structure consisting of two or more peaks in the 9.4-10.0 GeV region (see Fig. 1). The immediate (and correct) interpretation was a bound system of a quark-antiquark pair, charge $``$1/3 quarks. The bound system was named the Upsilon $`(\mathrm{{\rm Y}})`$.
The DORIS $`e^+e^{}`$ storage ring at DESY, at the time of the $`\mathrm{{\rm Y}}`$ discovery, had insufficient energy to produce $`\mathrm{{\rm Y}}`$’s. The machine energy was increased, and in 1978, straining their RF, physicists at DORIS observed two narrow resonances, $`\mathrm{{\rm Y}}`$(1S) and $`\mathrm{{\rm Y}}`$(2S). They could go no higher.
The CESR $`e^+e^{}`$ storage ring at LNS, Cornell, gave first luminosity to the CLEO and CUSB detectors in October, 1979. The $`\mathrm{{\rm Y}}`$(1S) and $`\mathrm{{\rm Y}}`$(2S) resonances were quickly located, and in December, in time to be “added in proof” to the Lab’s Christmas card, the $`\mathrm{{\rm Y}}`$(3S) was discovered (see Fig. 2).
The three resonances, $`\mathrm{{\rm Y}}`$(1S), $`\mathrm{{\rm Y}}`$(2S), $`\mathrm{{\rm Y}}`$(3S) were all narrow, with widths less than the instrumental resolution (beam energy spread). The production rates, leptonic decay branching fractions, level spacings, all matched very well with the bound $`b\overline{b}`$, charge $``$1/3 quark interpretation.
While there was no doubt, by then, about the existence of the bottom quark, the studies needed to determine further properties were of its weak decays. These could not be obtained from ‘hidden beauty,’ because a bound $`b\overline{b}`$ system decays via the strong interaction, with $`b`$ and $`\overline{b}`$ quarks annihilating each other, forming gluons or a virtual photon. For studies of the weak decay of the bottom quark, “bare bottom,” or “naked beauty” was needed.
Bare bottom was discovered at CESR by the spring of 1980. A scan, measuring cross section for production of hadronic events vs center-of-mass energy, above the $`\mathrm{{\rm Y}}`$(1S), $`\mathrm{{\rm Y}}`$(2S), $`\mathrm{{\rm Y}}`$(3S), revealed another resonance. This one was measurably broad (see Fig. 2), indicative of a rapid decay into $`b`$-flavored mesons, $`\mathrm{{\rm Y}}`$(4S) $`B\overline{B}`$. The compelling evidence for bare bottom came from the yield of muons and electrons, which also peaked at the $`\mathrm{{\rm Y}}`$(4S) resonance (see Fig. 3), indicating the decay sequence $`\mathrm{{\rm Y}}`$(4S) $`B\overline{B}`$ (via the strong interaction), followed by $`BX\mathrm{}V`$ (via the weak interaction). Leptons, a tell-tale signature of a weak decay, established bare bottom.
### 2.2 Roughing out the Qualitative Features – 1980-88
A series of measurements, from 1980 to 1988, determined the qualitative features of the $`b`$ quark.
2.2.1 Semileptonic Decay Branching Fraction
If the $`b`$ decays by a charged current interaction, $`bW_{\mathrm{virtual}}c`$ or
$`bW_{\mathrm{virtual}}u`$, then by simple counting of the $`W_{\mathrm{virtual}}`$ final states ($`\overline{u}d`$, $`\overline{c}s`$, $`\tau \nu `$, $`\mu \nu `$, $`e\nu `$), allowing for a factor of 3 for color for $`\overline{u}d`$ and $`\overline{c}s`$, one predicts a semileptonic decay branching fraction of 1/9. Phase space suppresses $`\overline{c}s`$ and $`\tau \nu `$, and hadronic final state interactions enhance $`\overline{u}d`$ and $`\overline{c}s`$, leading to a theoretical prediction for the semileptonic decay branching fraction of $``$ 12%. Early measurements were in qualitative agreement. (Aside – now, in the precision era, the measurements appear to be 1-2% below the theory, and that difference is not understood.)
2.2.2 Ruling out Topless Models
Giving that the bottom quark exists, is there a top quark? That was a very real question in the early 1980’s, because searches at PEP and PETRA had come up empty, and it was (then) hard to imagine that top was more than 2-3 times heavier than bottom. Producing “topless models” became an industry among theorists. Shooting them down became an experimental responsibility.
The simplest of the topless models had $`b`$ a weak isospin singlet, decaying by flavor-mixing with $`s`$ and $`d`$. In this case, the GIM mechanism would be inoperative, and there would be flavor-changing neutral decays of $`b`$, in particular $`bs\mathrm{}^+\mathrm{}^{}`$. Kane and Peskin derived a lower limit on the ratio $`\mathrm{\Gamma }(bs\mathrm{}^+\mathrm{}^{})/\mathrm{\Gamma }(bc\mathrm{}\nu )`$, for this topless model. CLEO (1984) and Mark J (1983) showed that the ratio was below the Kane-Peskin limit, ruling out that model.
A more complicated topless model had $`b`$ a weak isospin singlet, but decaying not by flavor mixing but by some new mechanism – exotic decays, which gave rise to enhanced yields of charged leptons and/or neutrinos and/or baryons. CLEO (1983) knocked that model off, by measuring yields of $`\mu `$, $`e`$, $`p`$, and missing energy.
The last stand of topless models was a particularly ugly one due to Henry Tye. It had $`b`$ in a right-handed doublet with $`c`$. Its decays mimicked $`bcW_{\mathrm{virtual}}^{}`$ reasonably well. However, its predicted production asymmetry, in $`e^+e^{}(\gamma \mathrm{or}Z^0)b\overline{b}`$, was very different, in the $`\gamma Z^0`$ interference region, from the predictions for $`b`$ a left-handed doublet with $`t`$. Experiments at PETRA (1985) established the left-handed doublet nature of $`b`$, killing the final topless model. Although it wouldn’t be discovered for another 10 years, by 1985 it was clear that top had to exist.
2.2.3 $`|V_{ub}/V_{cb}|`$
Does $`b`$ decay predominantly to $`c(bcW_{\mathrm{virtual}}^{})`$ or to $`u(buW_{\mathrm{virtual}}^{})`$? While there was a bias favoring $`bc`$, as of 1980 there was no strong theoretical argument favoring $`bc`$, nor any experimental evidence.
First evidence came from the kaon yield in $`B`$ decay (CLEO, CUSB, 1982), which was large, as would be expected for a $`bcs`$ sequence. The yield implied $`|V_{ub}/V_{cb}|^2<0.15`$.
Next evidence came from the lepton momentum spectrum. Since $`u`$ is lighter than $`c`$, $`bu\mathrm{}\nu `$ will have a stiffer lepton spectrum than $`bc\mathrm{}\nu `$ (see Fig. 4). By measuring the lepton spectrum and fitting to a mix of $`bu\mathrm{}\nu `$ and $`bc\mathrm{}\nu `$, CLEO (1984) established that $`|V_{ub}/V_{cb}|^2<0.04`$. By concentrating on the endpoint region of the spectrum, with more data, CLEO (1987) established that $`|V_{ub}/V_{cb}|^2<0.02`$. Finally, with still more data, CLEO (1990) saw leptons beyond the $`bc\mathrm{}\nu `$ endpoint, establishing that $`|V_{ub}/V_{cb}|^2>0`$.
2.2.4 B Reconstruction
Although there was no doubt about the existence of the $`b`$ quark in its bare form, and thus no doubt about the existence of $`b`$-flavored hadrons, it was nonetheless important to ‘reconstruct’ them, to assemble the decay products and show that they came from the decay, e.g., of a $`B`$ meson. Aside from the aesthetics of “it’s got to be there, so you must show that it is there”, $`B`$ reconstruction was needed to determine the $`B`$ meson mass. CLEO did this in 1983.
2.2.5 b Lifetime
The $`\mathrm{{\rm Y}}`$(4S) has a mass just slightly above $`B\overline{B}`$ threshold. As a result, the only decay of $`\mathrm{{\rm Y}}`$(4S) is $`\mathrm{{\rm Y}}`$(4S) $`B\overline{B}`$. There are no extra particles to confuse the situation. However (at a symmetric $`e^+e^{}`$ collider) the $`B`$ and $`\overline{B}`$ are moving very slowly, with momenta $``$300 MeV/c, $`\beta 0.06`$. This is not a suitable environment for determining $`b`$ lifetime.
For $`e^+e^{}`$ collisions at higher energies, the $`b`$-flavored hadrons will be moving faster, but the signal-to-noise will be less favorable (1 in 11, rather than 1 in 4), and the events will be more complicated (many particles in addition to the two $`b`$-flavored hadrons). For measuring the $`b`$ lifetime, the higher speed of the $`b`$-flavored hadrons overcomes the disadvantages just mentioned, and makes lifetime measurements possible. In 1983, MAC and Mark II, at PEP, made such measurements. They found the lifetime to be long, $`1`$ ps, implying $`|V_{cb}|0.05`$. Had $`|V_{cb}|`$ been more like $`|V_{us}|=0.22`$, the lifetime would have been a factor of 20 shorter, and probably would have then been too short to measure. This result, the long $`b`$ lifetime, was the first surprise in $`b`$ physics.
2.2.6 $`B^0\overline{B}^0`$ Mixing
It was recognized early on that the $`B^0`$ could, in principle, transform itself into a $`\overline{B}^0`$, just as the $`K^0`$ transforms into a $`\overline{K}^0`$. The diagram is a box diagram, with $`W^+,W^{}`$ on two parallel sides of the box, and $`u,c,t,\overline{u},\overline{c},\overline{t}`$ on the other two sides. In the limit of equal $`u,c,t`$ quark masses, the summed diagrams vanish, via the GIM mechanism. For “reasonable” values of the top quark mass – say 20 GeV, the rate for $`B^0\overline{B}^0`$ mixing would be immeasurably small.
In 1987, the UA1 experiment at CERN, in $`\overline{p}p`$ collisions, saw like-sign dilepton pairs, which they interpreted as $`B^0\overline{B}^0`$ mixing. As the extremely complicated environment of $`\overline{p}p`$ collisions made interpretation difficult, few people took this result seriously.
Later that same year, the ARGUS experiment at DESY, in the clean environment of $`e^+e^{}`$ collisions at the $`\mathrm{{\rm Y}}`$(4S), also saw like-sign lepton pairs. This could not be ignored. The most natural interpretation (and the correct one) was that top was a lot heavier than people had thought. This interpretation was slow in being accepted.
### 2.3 The Early Days – Summary
By 1990, there was a clear answer to “What are the basic features of the bottom quark?”.
* It is a member of a left-handed weak isospin doublet, with a (very heavy) top quark.
* It decays dominantly to the charm quark, via a charged current interaction $`bcW_{\mathrm{virtual}}^{}`$. $`|V_{cb}|0.04`$, so the coupling of the third generation to the second is smaller than the coupling of the second generation to the first $`(|V_{us}|=0.22)`$.
* Its decay to the up quark, $`buW_{\mathrm{virtual}}^{}`$, is small but not zero, so the coupling of third generation to first is the smallest of the three couplings.
* Because the top quark is so massive, the GIM mechanism breaks down for loop and box diagrams involving $`b`$. One consequence is the observed large rate for $`B^0\overline{B}^0`$ mixing. Another should be measurable rates for penguin decays. Lets look!
## 3 A Change in Objective
In “the early days,” the objective was to determine the basic features of the bottom quark. By 1990, that had been done, and the emphasis shifted.
Now, in recent times, and in the future, the objective is to use the bottom quark to probe the Standard Model, and search for physics ‘Beyond the Standard Model’.
There are two approaches to this probing and searching. These are:
* “Overdetermining the CKM Matrix”
* Measuring rates for Electroweak Penguins
I’ll examine each of these a bit later, but first a brief review of the flavor sector of the Standard Model.
### 3.1 The Flavor Sector of the Standard Model
Quarks come in left-handed weak isospin doublets, and decay via emission of (real or virtual) $`W^\pm `$ bosons.
$$\left(\genfrac{}{}{0pt}{}{t}{b}\right)tbW^+,\left(\genfrac{}{}{0pt}{}{c}{s}\right)csW^+$$
$`(3.1\text{-}1)`$
Thus $`t`$ decays to $`b`$ and a real $`W^+`$, while $`c`$ decays to $`s`$ and a virtual $`W^+`$.
Question: How do the lower, lighter members of the doublets decay? The $`b`$ quark can’t decay $`btW^{}`$; that violates energy conservation.
Answer: The mass eigenstates $`(b,s,d)`$ and the weak interaction eigenstates $`(b^{},s^{},d^{})`$ are slightly different. The $`b`$ quark “flavor mixes” with $`s`$ and $`d`$ quarks, according to the CKM matrix (Cabibbo, Kobayashi, Maskawa):
$$\left(\begin{array}{c}d^{}\\ s^{}\\ b^{}\end{array}\right)=\left(\begin{array}{ccc}V_{ud}& V_{us}& V_{ub}\\ V_{cd}& V_{cs}& V_{cb}\\ V_{td}& V_{ts}& V_{tb}\end{array}\right)\left(\begin{array}{c}d\\ s\\ b\end{array}\right)$$
(3.1-2) $``$ $``$
wk. int. mass
So, the $`b`$ quark (mass eigenstate) is a mixture of $`b^{}`$ (dominant component, amplitude $`V_{tb}`$), $`s^{}`$ (amplitude $`V_{cb}`$), and $`d^{}`$ (amplitude $`V_{ub}`$), and decays $`bcW_V^{}`$ with an amplitude proportional to $`V_{cb}`$, due to its $`s^{}`$ component, and decays $`buW_V^{}`$ with an amplitude proportional to $`V_{ub}`$, due to its $`d^{}`$ component.
The CKM matrix is a unitary matrix, which places $`n^2`$ constraints on an $`n\times n`$ matrix.
Because the phase of each quark state is arbitrary, $`2n1`$ phases in the CKM matrix can be transformed away.
Thus, if there were only two families of quarks, $`(u,d)`$, and $`(c,s)`$, the CKM matrix would be a $`2\times 2`$ matrix, with 4 complex elements, 8 parameters. The unitarity of the CKM matrix reduces 8 to 4, and the arbitrariness of the quark state phases reduces 4 to 1, a single parameter, the Cabibbo angle. The CKM matrix for two families is described by a single parameter, and can be made real.
But there are (at least) three families of quarks, $`(u,d),(c,s),\mathrm{and}(t,b)`$. The CKM matrix is a $`3\times 3`$ matrix, and the arithmetic goes 9 complex elements $``$ 18 parameters $``$9, for unitarity, $``$5, for arbitrary phases $`=`$ 4. The CKM matrix for three families is described by four parameters: 3 angles and 1 phase. This phase cannot be transformed away. If the phase is nonzero, weak decays will not be CP invariant.
Thus
$$3\mathrm{families}\mathrm{CP}\mathrm{violation}.$$
This was Kobayashi and Maskawa’s insight in 1973, before the charm quark had been discovered, let alone any members of the third family.
3.1.1 CP Violation
CP violation was observed in neutral kaon decay in 1964, by Christenson, Cronin, Fitch and Turlay. Given Kobayashi and Maskawa’s insight, that implies 3 families.
The $`b`$ quark, hence a third family, was observed in 1977, by Lederman and collaborators. That implies CP violation.
So, why is everyone making such a fuss about CP violation? It’s expected, observed, explained, isn’t it? There are two reasons why CP violation is now considered a “big deal”.
A) The CP violation given by the phase in the CKM matrix appears to be too small to account for the baryon-antibaryon asymmetry of the Universe at early times. Cosmology requires “New Physics,” and it must be CP-violating New Physics.
B) Measurements using CP-violating $`b`$ decays can help determine (overdetermine) the CKM matrix, hence probe the correctness of the Standard Model (or see New Physics).
Reason B) has CP violation in $`b`$ decay a useful tool for probing, searching, while reason A) has it as the primary object of study.
My own view is that too much attention is give to A) (hey, it’s got a lot of PR value), and not enough to B). But, in any case, whether you prefer to focus on A) or B), what studies you’ll perform will be much the same. CP violation in $`b`$ decay should be, and will be, studied.
3.1.2 Penguins, $`B^0\overline{B}^0`$ Mixing and the GIM Mechanism
The GIM mechanism causes flavor-changing neutral currents to vanish at tree level. It also suppresses FCNC beyond tree level, for loop and box diagrams. Let’s work this through for an important example, $`bs\gamma `$. (The same argument applies to gluonic penguins $`bsg`$.)
The Feynman diagrams for $`bs\gamma `$ are shown in Fig. 5 The overall amplitude is the sum of the three diagrams, with $`u`$, $`c`$, and $`t`$ quark inside the loop. The CKM factors are as shown on the figure, and thus the amplitude is
$$A=A\left(m_u^2\right)V_{ub}V_{us}^{}+A\left(m_c^2\right)V_{cb}V_{cs}^{}+A\left(m_t^2\right)V_{tb}V_{ts}^{}$$
$`(\mathrm{3.1.2}\text{-}1)`$
The amplitudes $`A(m_u^2),A(m_c^2),A(m_t^2)`$ depend only on masses, their flavor dependence having been removed by factoring out the CKM pieces. But, from unitarity
$$V_{ub}V_{us}^{}+V_{cb}V_{cs}^{}+V_{tb}V_{ts}^{}=0\left(\mathrm{Unitarity}\right)$$
$`(\mathrm{3.1.2}\text{-}2)`$
Thus, if $`m_u=m_c=m_t`$, then $`A=0`$. That’s GIM, the cancelation of the different terms in the sum. There is suppression, and the closer the three masses are to each other, the more the suppression.
But, $`m_tm_c`$, $`m_u`$. So, the cancelation is far from complete. The amplitude $`A`$ is proportional to $`m_t^2`$, and, since $`m_t`$ is a lot larger than originally expected, penguins in $`b`$ decay are also a lot larger than originally expected.
A similar argument applies to $`B^0\overline{B}^0`$ mixing. The Feynman diagrams are shown in Fig. 6. There is a double sum over $`u,c,t;\overline{u},\overline{c},\overline{t}`$. If $`m_u=m_c=m_t`$, the sum is zero, due to the unitarity of the CKM matrix. The heavy top badly breaks GIM, with an amplitude for mixing proportional to $`m_t^2`$.
### 3.2 Why is Bottom Quark Physics so Interesting?
We’re now in a position to answer the question “Why is bottom quark physics so interesting, such a good probe of New Physics?”
The answer is, “Because the TOP quark is so massive!”
The massive top quark gives rise to substantial $`\overline{B}^0\overline{B}^0`$ mixing, and substantial rates from loop diagrams (Penguins). Both of these are powerful tools for testing the Standard Model, for searching for New Physics.
3.2.1 Using $`B^0\overline{B}^0`$ Mixing to Learn Weak Phases
Consider a decay of a neutral $`B`$, with $`B^0`$ and $`\overline{B}^0`$ reaching the same final state, $`B^0X`$ and $`\overline{B}^0X`$. Examples are $`B^0\psi K^0\psi K_s^0`$ and $`\overline{B}^0\psi \overline{K}^0\psi K_s^0`$; and $`B^0\pi ^+\pi ^{}`$ and $`\overline{B}^0\pi ^+\pi ^{}`$. A particle born as a $`B^0`$ has two routes to this final state: i) The direct one $`B^0X`$, and ii) the indirect one, through $`B^0\overline{B}^0`$ mixing, $`B^0\overline{B}^0X`$. The amplitudes for these two routes will add coherently, and interfere. Similarly, the amplitudes for $`\overline{B}^0X`$ and $`\overline{B}^0B^0X`$ will add coherently and interfere. Immediately after birth, a particle born as a $`B^0`$ will be a $`B^0`$, but, over time, it will mix into $`\overline{B}^0`$, and so time development is the key. By tagging particle flavor at birth, comparing $`|A(B^0X)+A(B^0\overline{B}^0X)|^2`$ with $`|A(\overline{B}^0X)+A(\overline{B}^0B^0X)|^2`$, studying the time development of both, one can determine
$$\mathrm{sin}\left(\varphi _{\mathrm{Mixing}}+2\varphi _{B^0X}\right)$$
The expected value of $`\varphi _{B^0\psi K_s}`$ is zero, while the Standard Model value of $`\varphi _{\mathrm{Mixing}}`$ is $`2\beta `$, so, studying $`B\psi K_s^0`$ is the much talked about “measurement of $`\mathrm{sin}2\beta `$”.
Note that one must study time development. This class of measurements, time development of tagged $`B^0,\overline{B}^0`$ to a common final state, is the rationale behind asymmetric $`e^+e^{}`$ colliders at the $`\mathrm{{\rm Y}}`$(4S). The asymmetric initial state energies has the center-of-mass moving in the lab, so the decay time can be measured.
3.2.2 Electroweak Penguins as Probes of New Physics at High Mass Scales
The Standard Model diagram for $`bs\gamma `$ is shown in Fig. 7. The photon may be emitted from any of the charged lines. The top quark internal line is shown, because it is the excess of $`m_t^2`$ above $`m_c^2,m_u^2`$, that breaks the GIM mechanism. The mass scale of the diagram is set by the masses of the particles in the loop – $`m_t`$ and $`M_W,100`$ GeV.
Consequently, contributions from New Physics (e.g., charged Higgs, SUSY particles; see Fig. 7) will show up for New Physics masses in that same range. So, penguins probe for New Physics up to masses $``$500 GeV.
This argument, given for electroweak penguins, applies also to gluonic penguins. However, electroweak penguins, in particular $`bs\gamma `$, has the advantage that its rate can be calculated, within the Standard Model, and Beyond the Standard Model, to a precision of 10%.
3.2.3 Learning Weak Phases from Penguin-Tree Interference
Many rare $`B`$ decays involve both penguin and tree amplitudes, while some related decays are pure penguin, or pure tree. By studying relative rates and CP asymmetries, one can sort the phases out, and determine weak phases.
As an example, consider $`B^{}K^{}\pi ^0`$, $`B^+K^+\pi ^0`$, and $`B^\pm K^0\pi ^\pm `$. The first two involve both penguin and tree amplitudes, while the last is pure penguin. The amplitudes for the three processes are
$$A(B^{}K^{}\pi ^0)=A_P+A_Te^{i\varphi _W}e^{i\varphi _S}$$
$$A(B^+K^+\pi ^0)=A_P+A_Te^{i\varphi _W}e^{i\varphi _S}$$
$`(\mathrm{3.2.3}\text{-}1)`$
$$A(B^+K^0\pi ^+)=A_P$$
where $`\varphi _W`$ is the difference in weak phase between penguin and tree, and $`\varphi _S`$ is the difference in strong phase between penguin and tree. Squaring amplitudes, one sees, for the first two modes
$$|A|^2=A_P^2+2A_PA_T[\mathrm{cos}\varphi _W\mathrm{cos}\varphi _S\mathrm{sin}\varphi _W\mathrm{sin}\varphi _S]+A_T^2$$
$`(\mathrm{3.2.3}\text{-}2)`$
while the square for the third mode is just $`|A_P|^2`$.
Thus, the rate difference, i.e., CP asymmetry, gives $`\mathrm{sin}\varphi _W\mathrm{sin}\varphi _S`$, while the rate sum, compared with the third mode, gives $`\mathrm{cos}\varphi _W\mathrm{cos}\varphi _S`$. Of course, there are complications to the naive picture just presented, due to electroweak penguins, color-suppressed trees, and long distance rescatterings. But, the decays do depend on strong and weak phases, roughly as indicated, and by studying several rare decays one can learn weak phases.
### 3.3 How Does One Determine Elements of the CKM Matrix?
$$\left(\begin{array}{ccc}V_{ud}& V_{us}& V_{ub}\\ V_{cd}& V_{cs}& V_{cb}\\ V_{td}& V_{ts}& V_{tb}\end{array}\right)$$
$`(3.3\text{-}1)`$
Rates for nuclear beta decay, compared to the rate for muon decay, gives a very precise determination of the magnitude of $`V_{ud}`$. Kaon and hyperon decay rates give good determinations of the magnitude of $`V_{us}`$. Assorted studies of charm decay give rough measurements of $`|V_{cd}|`$ and $`|V_{cs}|`$. However, since the third generation is only weakly coupled to the first two, these studies determine only a single parameter, $`\lambda =\mathrm{sin}\theta _{\mathrm{Cabibbo}}`$.
Studies of $`b`$ decay determine two more parameters. In particular, the rate for $`bc\mathrm{}\nu `$ determines $`|V_{cb}|`$, and the rate for $`bu\mathrm{}\nu `$ determines $`|V_{ub}|`$.
Can $`|V_{ts}|`$, $`|V_{td}|`$ be determined from studies of top decay? Not soon! The rate for $`tW^+s`$ is proportional to $`|V_{ts}|^2`$, and the rate for $`tW^+d`$ is proportional to $`|V_{td}|^2`$. Measuring those rates would give $`|V_{ts}|`$ and $`|V_{td}|`$. But the expected value for $`|V_{ts}|^2`$ is $`2\times 10^3`$, while that for $`|V_{td}|^2`$ is $`10^4`$, while $`|V_{tb}|^21`$, giving a dominant decay $`tWb`$. It will be a while (quite a while) before top decay branching fractions at the $`10^310^4`$ level are measured.
So, for the foreseeable future, the situation is this. We can determine three magnitudes in the CKM matrix – $`\lambda `$, $`|V_{cb}|`$, $`|V_{ub}|`$ – from tree-level processes, theoretically secure, relatively free from possible New Physics contributions, reliably giving what they claim to determine. All else will come from loops, boxes, places where New Physics is likely to enter. Thus, if an “overdetermination of the CKM matrix” finds an inconsistency, that does not mean a problem with the CKM matrix, but rather that the relation to the CKM matrix of some measurable has been changed by New Physics. For example, if $`\mathrm{sin}2\beta `$ as determined from the time development of tagged $`B^0\psi K_s^0`$ disagrees with expectations, that would mean that the phase of the $`B^0\overline{B}^0`$ mixing amplitude is not $`\mathrm{sin}2\beta `$, but has been altered by New Physics contributions to mixing.
Let’s rewrite the CKM matrix in a $`b`$-centric fashion. Taking $`|V_{cb}|=0(\lambda ^2)`$, $`|V_{ub}|=0(\lambda ^3)`$, and enforcing unitarity, we have, correct to $`0(\lambda ^3)`$
$$\left(\begin{array}{ccc}1\lambda ^2/2& \lambda & |V_{ub}|e^{i\gamma }\\ \lambda & 1\lambda ^2/2& |V_{cb}|\\ \lambda |V_{cb}||V_{ub}|e^{i\gamma }& |V_{cb}|& 1\end{array}\right)$$
$`(3.3\text{-}2)`$
Since $`\lambda `$ is already determined with high precision, this form makes apparent the urgency of good determinations of $`|V_{cb}|`$ and $`|V_{ub}|`$.
$`|V_{cb}|`$ is obtained from measurements of the $`B`$ meson lifetime, and either the rate for $`BD^{}\mathrm{}\nu `$ extrapolated to the point where $`D^{}`$ is at rest, or the rate for $`BX\mathrm{}\nu `$ inclusive, plus information on the $`b`$ quark mass and HQET Operator Product Expansion parameter $`\lambda _1`$. The $`b`$ lifetime and semileptonic decay branching fraction are well measured. CLEO has in hand data on $`BD^{}\mathrm{}\nu `$, and on moments of hadronic mass and lepton energy in $`BX\mathrm{}\nu `$ sufficient for $`\pm 4`$% determinations of $`|V_{cb}|`$, separately by each method. For now, $`|V_{cb}|=(39.5\pm 3.6)\times 10^3`$.
$`|V_{ub}|`$ is less well determined, and $`\gamma `$ even less well determined.
## 4 What is $`\mathbf{|}𝑽_{𝒖𝒃}\mathbf{/}𝑽_{𝒄𝒃}\mathbf{|}`$?
### 4.1 Limitations of Previous Approaches
In Section 2.2.3, I described progress during the early days in placing upper-limits on, and finally establishing a nonzero value for, $`|V_{ub}/V_{cb}|`$. All the approaches tried then had serious limitations. The kaon yield approach was really a measurement of $`|V_{cb}|`$, limiting $`|V_{ub}|`$ by $`|V_{cb}|`$’s deviation from 1.0. Since the total number of kaons produced per $`bcW_V`$ decay is uncertain at greater than the ten percent level, this approach was quickly discarded, once it was realized that $`|V_{ub}/V_{cb}|`$ was in the sub-ten-percent range.
Fitting the measured lepton spectrum in $`B`$ semileptonic decay to the predicted spectra for $`bc\mathrm{}\nu `$ and $`bu\mathrm{}\nu `$ hits its limit because, with the $`bu\mathrm{}\nu `$ rate less than 5% of the $`bc\mathrm{}\nu `$ rate, minor errors in modeling of the $`bc\mathrm{}\nu `$ spectrum cause major errors in the $`bu\mathrm{}\nu `$ yield. This approach has also been discarded.
The endpoint approach avoids sensitivity to the $`bc\mathrm{}\nu `$ modeling because it limits the focus to the lepton momentum range where $`bc\mathrm{}\nu `$ is small or zero. But here there is sensitivity to the modeling of $`bu\mathrm{}\nu `$. The fraction of the $`bu\mathrm{}\nu `$ spectrum in the endpoint windows cannot be reliably calculated, and its uncertainty limits accuracy of $`|V_{ub}|`$ by this method to $``$20%. While results from this approach are currently one of the two ways now giving useful results, future improvements to the 10% range and below seem unlikely. (Note added in proof: Leibovich, Low, and Rothstein, hep-ph/9909404 v2, show how to determine the fraction, using a measurement of the photon spectrum from $`bs\gamma `$.)
### 4.2 Neutrino Detection
The difficulty in studying $`bu\mathrm{}\nu `$ is the neutrino. If that particle were ‘detected’, its momentum measured, then the decay would cause no problems. Consequently, several of us in CLEO are attempting a new approach, ‘detecting’ the neutrino in a semileptonic decay via the missing 4-momentum in the event. Given a ‘detected’ neutrino, one can then carry out full reconstruction of exclusive semileptonic decays, or look at the mass distribution in inclusive semileptonic decays.
4.2.1 Exclusive Decays $`B\pi \mathrm{}\nu `$, $`B\rho \mathrm{}\nu `$
With the neutrino ‘detected’, the decays $`B\pi \mathrm{}\nu `$, $`B\rho \mathrm{}\nu `$, $`B\omega \mathrm{}\nu `$ have no undetected particles, and so the standard $`B`$ reconstruction technique is applicable. The summed energy of the decay products of the candidate $`B`$ are compared with the beam energy, giving a difference $`\mathrm{\Delta }E`$ which should peak at zero. The summed vector momenta of the decay products of the candidate $`B`$, $`𝑷_{\mathrm{𝐜𝐚𝐧𝐝}}`$, are used to calculate the “beam constrained mass” $`M_{m\mathrm{}\nu }=\sqrt{E_{\mathrm{beam}}^2P_{\mathrm{cand}}^2}`$, which should peak at the $`B`$ mass.
We completed and published an analysis for $`B\pi \mathrm{}\nu `$, and $`B\rho /\omega \mathrm{}\nu `$ some time ago (PRL 77, 5000 (16 Dec. 1996)), based on a 4$`fb^1`$ data sample. The plots of mass and energy difference are shown in Fig. 8. The branching fraction accuracy (statistical plus systematic) gave a 12% uncertainty in $`V_{ub}`$, and that uncertainty should fall as $`1/\sqrt{}`$.
The big issue is model dependence – how the branching fraction for the exclusive modes are related to $`V_{ub}`$. Of course, they are proportional to $`|V_{ub}|^2`$, with the constants of proportionality related to form factors. It is through the uncertainty in the form factors that model dependence enters. For the 1996 analysis, we estimated this at $`\pm `$20% in $`|V_{ub}|`$. This will improve with more data, which will allow measurement of the $`q^2`$ dependence of the decays, providing constraints on models for form factors. It will also improve with better form factor calculations, from lattice gauge QCD and other techniques. Finally, studies of the decays $`D\pi \mathrm{}\nu `$, $`D\rho \mathrm{}\nu `$, where $`|V_{cd}|`$ is quite well known, can also help. One can expect an accuracy in $`|V_{ub}|`$ from CLEO’s existing, 14$`fb^1`$ data sample, in the $`\pm `$15% range, or better, depending on how much progress can be made on the model dependence.
4.2.2 Inclusive Decays, $`BX_u\mathrm{}\nu `$
Given a ‘detected’ neutrino, and a (really) detected charged lepton, one can calculate the mass of the hadronic system $`X`$ in the decay $`BX\mathrm{}\nu `$:
$$M_X^2=M_B^2+M_\mathrm{}\nu ^22E_BE_\mathrm{}\nu +2P_BP_\mathrm{}\nu \mathrm{cos}\theta _{\mathrm{}\nu ,B}$$
All quantities in this equation are known except $`\theta _{\mathrm{}\nu ,B}`$, the angle between the $`B`$ meson and the $`\mathrm{}\nu `$ system (everything evaluated in the lab frame). The total lack of knowledge of $`\theta _{\mathrm{}\nu ,B}`$ results in a smearing in the determination of $`M_X^2`$, which is reasonably small since $`P_B`$ is small ($``$300 MeV/c).
The game plan, then, is to measure the $`M_X^2`$ distribution, given neutrino and charged lepton, and then fit that distribution with a sum of $`bu\mathrm{}\nu `$ and $`bc\mathrm{}\nu `$. The contributions from $`bc\mathrm{}\nu `$ will include $`D,D^{},`$ and heavier stuff. The contributions from $`bu\mathrm{}\nu `$ will dominantly be below the $`D`$ meson mass, consisting of $`X_u`$ objects like $`\pi ,\rho ,A_1,A_2,`$ etc. A calculation of the expected $`X_u`$ mass distribution is possible, for example from a naive spectator model, or more properly from HQET and OPE. If one could measure $`M_X`$ to high precision, separating $`bu\mathrm{}\nu `$ from $`bc\mathrm{}\nu `$ would be easy, and an inclusive measurement of $`|V_{ub}|`$, with relatively little model dependence, would be possible.
Unfortunately, the measurement of $`M_X^2`$ so far achieved has rather poor resolution, due to the inaccuracy in determining the neutrino vector momentum. This inaccuracy is not so much from the inaccuracy in measurement of individual particles, but rather from ‘messups’ (inefficiencies in detecting charged particles and photons, false tracks and photons), and also from undetectable particles ($`K`$-long, neutrons, second neutrinos in the event). The consequence of the poor resolution in $`M_X`$ is that there is a low-mass-side tail to $`bc\mathrm{}\nu `$, which swamps the $`bu\mathrm{}\nu `$ contribution.
An analysis has been completed (Scott Roberts’ Ph.D. thesis, University of Rochester, 1997), but not submitted for journal publication. To suppress the $`bc\mathrm{}\nu `$ component, we required $`P_{\mathrm{}}>2.0`$ GeV/c, a momentum bite a factor of 2 bigger than the $`P_{\mathrm{}}>2.3`$ GeV/c typical of an endpoint analysis. The choice of 2.0 GeV/c was a compromise between reducing model dependence (wanting a lower momentum cut) and suppressing $`bc\mathrm{}\nu `$ (wanting a higher momentum cut).
The measured $`M_X^2`$ distribution is shown in Fig. 9. The fitted components from $`bc\mathrm{}\nu `$ and $`bu\mathrm{}\nu `$ are also shown. In the region where $`bu\mathrm{}\nu `$ is substantial, the $`bc\mathrm{}\nu `$ background is about twice the $`bu\mathrm{}\nu `$ signal. Taking faith in our modeling of the $`bc\mathrm{}\nu `$ background (though allowing a systematic error for its uncertainty), we obtained a fit, from a 5$`fb^1`$ data sample, which gives $`|V_{ub}|`$ to $`\pm `$16% – statistical plus systematic. We did not carry out a careful study of model dependence, but since the lepton momentum bite is twice as large as that for the endpoint analysis, one would expect a model dependence that is twice as small – $`\pm `$10% instead of $`\pm `$20%.
The value of $`|V_{ub}|`$ obtained from the fit is quite reasonable, and the combined nominal error, ($`\pm `$16% with $`\pm `$10%) are competitive. But the plot is certainly not very convincing. The $`bc\mathrm{}\nu `$ component is just too large in the $`bu\mathrm{}\nu `$ signal region. And we would like to push the lepton momentum cut down, say to 1.8 or 1.6 GeV/c, which would make the $`bc\mathrm{}\nu `$ background several times larger. So, our plan is not to publish this analysis, but to work on it some more – a lot more.
* We will use the full CLEO II data sample of 14$`fb^1`$, a factor of 3 increase from that in Scott Roberts’ analysis. (This is the easy one.)
* We will improve the accuracy with which neutrinos are ‘detected’ and their momenta determined, by upgrading our algorithm for distinguishing between showers in the electromagnetic calorimeter caused by photons and by hadrons; by improving various aspects of charged particle tracking; and by pushing to lower momentum our electron identification capabilities (we veto events with more than one charged lepton, hence more than one neutrino).
* Finally, we will study the correctness of our simulation of the $`bc\mathrm{}\nu `$ component, to be sure we are correctly modeling the low-mass tail. (For example, we will fake $`K`$-long events by finding events with $`K`$-shorts, then pitching the $`K`$-short, and see if the $`M_X^2`$ spectra so obtained for data and Monte Carlo agree.)
The original motivation for neutrino ‘detection’ was for studying inclusive decays, with its use for exclusive decays an afterthought. We still view the inclusive approach as the best hope for a measurement of $`|V_{ub}|`$ to $`\pm `$10% accuracy.
## 5 Rare Hadronic $`𝑩`$ Decays
### 5.1 Introduction
I should start this section by saying what I mean, and indeed what is typically meant, by “rare”, as it refers to $`B`$ decays. A “rare” $`B`$ decay is one which involves penguin or box diagrams. With this definition, it is easy to see why the field of rare $`B`$ decays is ahead of the field of rare kaon decays, why $`bs`$ processes have been studied, while $`sd`$ processes much less so. The CKM factor for $`bs`$ penguins is $`V_{tb}V_{ts}^{}`$, while that for the dominant, $`bc`$ is $`V_{cb}`$. $`|V_{tb}V_{ts}^{}/V_{cb}|^21`$. For kaon decays, the penguin with top quark in the loop has a CKM factor $`V_{ts}V_{td}^{}`$, while that for the dominant, $`su`$ tree is $`V_{us}`$. $`|V_{ts}V_{td}^{}/V_{us}|^23\times 10^6`$, so the branching fractions for rare $`B`$ decays are typically 5-6 orders of magnitude larger than those for rare kaon decays – $`10^510^6`$ vs $`10^{11}10^{12}`$.
As we saw in Section 3.2.3, rare decays involving penguins often also involve $`bu`$ trees (see Fig. 10). The example given there was $`BK\pi `$, a “Cabibbo-allowed penguin”, i.e., a $`bs`$ penguin. The tree diagram there is the “Cabibbo-suppressed $`bu`$ tree”, $`buW_V^{}`$, $`W_V^{}\overline{u}s`$. The Cabibbo-suppression is in the decay of the virtual $`W`$, $`W^{}\overline{u}s`$, rather than the Cabibbo-allowed decay $`W^{}\overline{u}d`$. This same mix of Cabibbo-allowed penguin plus (sometimes) Cabibbo-suppressed $`bu`$ tree occurs for $`BK^{}\pi `$, $`BK\rho `$, $`BK\omega `$, $`BK^{}\rho `$. The Cabibbo-allowed penguin diagram contributes to all of these, while the Cabibbo-suppressed $`bu`$ tree is absent for the charge modes involving neutral $`K`$ or $`K^{}`$.
There is also a class of decays involving Cabibbo-allowed $`bu`$ trees, i.e., $`buW_V^{},W_V^{}\overline{u}d`$, and Cabibbo-suppressed penguins, i.e. $`bd`$ penguins (Fig. 10c,d). Examples include $`B\pi \pi `$, $`B\pi \rho `$, $`B\pi \omega `$, $`B\rho \rho `$. In fashion similar to the ‘allowed penguin, suppressed tree’ class, there are particular modes for which the Cabibbo-suppressed penguin is absent, e.g. $`B^\pm \pi ^\pm \pi ^0`$.
So, Penguin-Tree interference is the rule rather than the exception in rare hadronic $`B`$ decays. And the exceptions, modes which are pure allowed penguins, or pure allowed $`bu`$ trees, help to sort out the interference.
As mentioned in Section 3.2.3, the simple picture is complicated by electroweak penguins $`bsZ^0`$ (we’ve been talking about gluonic penguins $`bsg`$), color-suppressed trees, long distance rescattering. It will require careful study of many rare decays before a precise value of the CKM phase $`\gamma `$ can be obtained. But as we will see, some qualitative information can already be obtained.
### 5.2 The Data Sample
CLEO has 10 million events of the form $`e^+e^{}\mathrm{{\rm Y}}\left(4\mathrm{S}\right)B+\overline{B}`$, and has recently completed analysis of several of the rare decay modes. The reconstruction is conventional, with the summed energy of the decay products of the candidate $`B`$ compared with the beam energy, and the summed vector momenta of the decay products of the candidate $`B`$ used to calculate the beam-constrained mass.
There are substantial backgrounds to rare $`B`$ decays, not from the dominant $`bc`$ tree decays, but from the continuum background process $`e^+e^{}q\overline{q}`$, $`q=u,d,c,s`$. These backgrounds are 2-jet-like, and are suppressed by a maximum likelihood fit, inputting many ‘shape variables’.
Examples of $`B`$ mass plots and $`\mathrm{\Delta }E`$ plots are shown in Fig. 11.
### 5.3 Results
5.3.1 $`BK\pi ,\pi \pi ,KK`$
Results for these modes are given in Table 1 All four $`BK\pi `$ modes have been convincingly seen. Only one of the three $`B\pi \pi `$ modes has been convincingly seen, though the evidence for $`B\pi ^\pm \pi ^0`$ is fairly good. No $`BKK`$ mode has been seen, nor were they expected to be.
5.3.2 $`B`$ Decays Involving $`\eta `$ or $`\eta ^{}`$
Results for the decay $`B(\eta \mathrm{or}\eta ^{})`$ ($`K`$ or $`K^{}`$) are shown in Table 2. One sees that $`B\eta ^{}K`$ is big, much larger than all the others. $`B\eta K^{}`$ is seen, and is larger than $`B\eta K`$.
The interpretation of these results is far from clear.
* The $`\eta ^{}`$ could perfectly well contain a $`c\overline{c}`$ component, and if it did, a Cabibbo-allowed $`bc`$ tree could contribute (as it does for $`B\psi K)`$. This situation would lead to enhanced branching fractions for both $`B\eta ^{}K`$ and $`B\eta ^{}K^{}`$.
* As pointed out by Lipkin, there will be contributions from the gluonic penguins $`bsg,gs\overline{s}`$ and $`bsg,gu\overline{u}/d\overline{d}`$, and these diagrams will interfere. Lipkin argues that this will enhance $`B\eta ^{}K`$ and $`B\eta K^{}`$ relative to $`B\eta ^{}K^{}`$ and $`B\eta K`$.
The data show some features of both suggestions, but at present there is no quantitative understanding.
5.3.3 Decays $`B`$ Pseudoscalar Vector
Only a smattering of these have been seen so far, e.g. $`B^\pm \omega \pi ^\pm `$, $`B\rho \pi `$, $`BK^{}\pi `$. CLEO’s analyses of the Pseudoscalar-Vector modes are finished for the full data sample of 10 million $`B\overline{B}`$ events for about half of the decay modes. Results are given in Table 3.
### 5.4 Search for Direct CP Violation in $`B`$ Decay
If some $`B`$ decay has contributions from two (or more) amplitudes $`A_1`$, $`A_2`$, with relative weak phase $`\varphi _W`$, and relative strong phase $`\varphi _S`$, i.e., a total decay amplitude $`A=|A_1|+|A_2|e^{i\varphi _S}e^{i\varphi _W}`$, then there will be direct CP violation in the decay, which will show up as a charge asymmetry
$$𝒜=\frac{\mathrm{\Gamma }(\overline{B}X)\mathrm{\Gamma }(B\overline{X})}{\mathrm{\Gamma }(\overline{B}X)+\mathrm{\Gamma }(B\overline{X})}$$
$`(5.4\text{-}1)`$
If $`|A_2||A_1|`$, then
$$𝒜2\frac{|A_2|}{|A_1|}\mathrm{sin}\varphi _S\mathrm{sin}\varphi _W$$
$`(5.4\text{-}2)`$
For penguin-tree interference, one expects $`|A_2/A_1|0.2`$. It’s less clear what to expect for $`\mathrm{sin}\varphi _S`$, but in the absence of some enhancement due to long range rescattering it will be small, probably less than 0.25. So we expect $`𝒜`$ less than 0.1. CLEO results, for five decay modes that have been convincingly seen, and are self-tagging, are given in Table 4. There are no surprises. All asymmetries are consistent with zero. There is not yet sufficient sensitivity to see CP violations at the level expected. Since the errors are dominantly statistical, and are based on 10 million $`B\overline{B}`$ pairs, it will likely be a while before nonzero asymmetries are established.
### 5.5 Interpretation of Rare Hadronic $`B`$ Decays
Now that CLEO has roughed out the rare $`B`$ decay terrain, what does it all mean? Recall that in Section 3.2.3, the motivation for studying rare $`B`$ decays was given as using them to determine weak phases. What can the existing rare $`B`$ decay data tell us about weak phases, in particular, about $`\gamma `$, the phase of $`V_{ub}^{}`$?
CLEO’s visiting theorist George Hou and CLEO members Jim Smith and Frank Würthwein have addressed that question. They assume naive factorization, use effective-theory matrix elements, and ignore annihilation type diagrams. With these assumptions, (and some more, mentioned below) they are able to express the amplitudes for all two-body rare $`B`$ decays in terms of a relatively small number of parameters.
The quark-level process $`bq_1\overline{q}_2q_3`$ is described, in effective theory, by ten parameters $`a_1\mathrm{}a_{10}`$. These are calculable within a QCD framework, and Hou et al. take two sets of values from the literature.
The binding of $`q_1\overline{q}_2`$ into mesons is described by decay constants $`f`$, $`(f_\pi ,f_K`$, $`f_K^{}`$, $`f_\rho ,f_\omega ,f_\varphi )`$ which are known.
The binding of $`q_3`$ and the spectator antiquark into a meson is described by form factors. For $`BPP`$, there is a single form factor $`F_0^{BP}`$ (but $`P=\pi ,K`$), while for $`BPV`$ there are two more, $`F_1^{BP},P=\pi ,K`$, and $`A_0^{BV},V=\rho ,\omega `$. Hou et al. lean on SU(3), with breaking, to relate $`F_0^{BK}`$ to $`F_0^{B\pi }`$, and $`A_0^{B\omega }`$ to $`A_0^{B\rho }`$. They also use the relation $`F_0^{BP}=F_1^{BP}`$, valid at $`q^2=0`$. They thus describe the decays of interest with just two form factors $`F_0^{B\pi }`$ and $`A_0^{B\rho }`$, rather than six.
Two of the penguin terms $`(a_6,a_8)`$ depend on the quark mass $`(m_s\mathrm{or}m_d)`$ and Hou et al. allow a free parameter $`R_{su}`$ to describe this dependence.
Hou et al. thus use five free parameters: $`F_0^{B\pi },A_0^{B\rho },R_{su},|V_{ub}/V_{cb}|,`$ and $`\gamma `$. They constrain $`|V_{ub}/V_{cb}|`$ by including the difference from its measured values, $`0.08\pm 0.02`$, as a term in $`\chi ^2`$.
They fit 14 branching fractions: $`K^{}\pi ^+,K^{}\pi ^0,\overline{K}^0\pi ^{},`$ $`\overline{K}^0\pi ^0,`$ $`\pi ^+\pi ^{},`$ $`\pi ^{}\pi ^0,`$ $`\rho ^0\pi ^{},`$ $`\omega \pi ^{},`$ $`\rho ^{}\pi ^\pm ,K^{}\pi ^+,\omega K^{},\omega \overline{K}^0,\varphi K^{},`$ and $`\varphi \overline{K}^0`$. They leave the $`\eta ^{}`$ and $`\eta `$ decay modes out of the fit, as something strange is happening with these modes. The $`\chi ^2`$ of the fit, as a function of $`\gamma `$, is shown in Fig. 12. The fit gives $`\gamma =114\pm 23`$ degrees.
The error just quoted is that from the branching fraction errors only, and does not include anything for theoretical uncertainty. Those must be estimated and included before a serious number for $`\gamma `$ can be quoted. However, from this exercise, so far, we can see that the data contain information sufficient for a precise determination of $`\gamma `$, given adequate theoretical understanding. Further, they argue for a large value of $`\gamma `$. I’ll take the liberty of assuming that the theoretical error won’t be more than $`\pm 50^{}`$, and interpret the rare $`B`$ results as saying $`\gamma >60^{}`$.
## 6 The Radiative Penguin Decay $`𝒃\mathbf{}𝒔𝜸`$
In Section 3.2.2, I argued that electroweak penguin processes, in particular $`bs\gamma `$, probe for New Physics up to masses $``$500 GeV. What’s been learned so far?
### 6.1 The Exclusive Decay $`BK^{}(890)\gamma `$
The observation of $`BK^{}(890)\gamma `$, in 1993, was the first clear observation of a penguin process. That analysis combined conventional $`B`$ reconstruction techniques with continuum suppression techniques, and used a likelihood ratio approach for further evidence. While the existence of the radiative penguin process $`bs\gamma `$ was clearly established by this analysis, it did not provide a good measurement of the inclusive rate (the theoretically interesting quantity), since the theoretical estimates of $`\mathrm{\Gamma }(BK^{}\gamma )/\mathrm{\Gamma }(bs\gamma )`$ ranged from 5% to 90%. A direct measurement of $`bs\gamma `$ was called for.
### 6.2 Branching Fraction for $`bs\gamma `$
The inclusive decay $`bs\gamma `$ gives a monoenergetic photon in the $`b`$ quark rest frame. That monoenergetic line is Doppler broadened by the motion of the $`b`$ quark in the $`B`$ meson frame, and the motion of the $`B`$ meson in the lab frame. But it remains a relatively narrow distribution. In Fig. 13, I show the photon energy distribution expected from $`bs\gamma `$, along with that expected from other $`B`$ decay processes. The $`bs\gamma `$ decays extend beyond those from other $`B`$ decay processes and a study of the photon spectrum above 2.0 GeV should cleanly give $`bs\gamma `$.
But wait. There are other curves shown on Fig. 13. One is the photon spectrum from initial state radiation in continuum production, $`e^+e^{}q\overline{q}\gamma `$. The other, the spectrum of $`\gamma `$’s from $`\pi ^0`$ decay in continuum production, $`e^+e^{}q\overline{q}\pi ^0X\gamma \gamma X`$. The sum of these two processes is more than two orders of magnitude larger than $`bs\gamma `$, at the $`bs\gamma `$ peak. Continuum suppression is absolutely essential.
In our 1995 measurement of the rate for $`bs\gamma `$, we used two different methods for continuum suppression. The first used eight carefully chosen event-shape variables. While no individual variable has strong discriminating power, each possesses some. We combined the eight variables into a single variable $`r`$, which tends toward $`+1`$ for $`bs\gamma `$ and tends towards $`1`$ for ISR and $`q\overline{q}`$. We used a neural network for the task of combining the eight variables into a single variable. This was CLEO’s first use of a neural network, and was single handedly pushed through the collaboration by Jesse Ernst, against strong opposition, much of it from his thesis advisor (me). That neural networks are now used extensively, and intelligently, within CLEO can be attributed to Jesse’s good understanding of the strengths and limitations of the technique.
The second method for continuum suppression has been dubbed “pseudoreconstruction”. In it, a high energy photon is combined with a kaon ($`K^\pm `$ or $`K_s^0`$) and 1-4 pions (of which one may be a $`\pi ^0`$), and tested for consistency with being a reconstructed $`B`$. (A $`\chi ^2`$ composed of $`B`$ mass and $`B`$ energy, $`\chi _B^2`$, is used for this test.) For those events with a pseudoreconstructed $`B`$, $`\theta _{tt}`$, the angle between the thrust axis of the candidate $`B`$ and the thrust axis of the rest of the event, gives additional discrimination against continuum background. In pseudoreconstruction, often one does not have the totally correct combination of particle (hence the “pseudo”), but this is not important (here), because the method is used only to suppress background, and not for a mode-by-mode $`B`$ reconstruction analysis.
In our 1995 result, we performed two separate analyses, the event-shape analysis and the pseudoreconstruction analysis, and averaged the branching fractions obtained from each (allowing for a small amount of event overlap). That result, $`(bs\gamma )=(2.32\pm 0.57\pm 0.35)\times 10^4`$, was based on a data sample of 3.0$`fb^1`$.
More recently, we’ve combined the two continuum suppression techniques into a single, unified analysis. For all events containing a high energy photon, we compute the neural net variable $`r`$. For the subset of events that pseudoreconstruct, with very loose requirements, we also calculate $`\chi _B^2`$ and $`\mathrm{cos}\theta _{tt}`$. For these events, we feed $`\chi _B^2`$, $`\mathrm{cos}\theta _{tt}`$, and $`r`$, into another neural network, obtaining a new net variable $`r_{\mathrm{comb}}`$. We assign a weight to each event, based on $`r_{\mathrm{comb}}`$ for pseudoreconstructed events, and on $`r`$ for those events which fail to reconstruct. In this way we’ve analyzed a 4.7$`fb^1`$ data sample. The photon energy spectrum obtained is shown in Fig. 14. The branching fraction obtained is
$$(bs\gamma )=(3.15\pm 0.35\pm 0.32\pm 0.26)\times 10^4$$
where the errors, in order, are statistical, systematic, and model dependent. This number is in excellent agreement with the NLO prediction of $`(3.28\pm 0.33)\times 10^4`$ of Chetyrkin, Misiak and Münz.
The comparison of experimental result with Standard Model prediction can be (has been) used to place restrictions on New Physics. For example, our conservative upper limit on the branching fraction, $`4.5\times 10^4`$, rules out a charged Higgs with Model II coupling for Higgs masses less than 200 GeV. (In SUSY, there would be additional particles, which could contribute with opposite sign, so the limitation is more complicated. However, a hunk of SUSY parameter space is ruled out.)
CLEO now has 14$`fb^1`$, 3 times the integrated luminosity used in the analysis just described. What’s holding us back? Well, look at the three errors on the branching fraction. Reducing the statistical error by $`1/\sqrt{3}`$ will do little good unless systematics and model dependence can be beaten down. That takes more time.
### 6.3 CP Asymmetry in $`bs\gamma `$
The CP asymmetry in $`bs\gamma `$, $`𝒜`$, defined by
$$𝒜\frac{|A(bs\gamma )|^2|A(\overline{b}\overline{s}\gamma )|^2}{|A(bs\gamma )|^2+|A(\overline{b}\overline{s}\gamma )|^2},$$
is very small, less than 1%, in the Standard Model. So, observing a nonzero value would be clear evidence for New Physics.
Suppose, in addition to the Standard Model decay amplitude for $`bs\gamma `$, $`A_{SM}`$, there is a New Physics amplitude, which differs in weak phase from $`A_{SM}`$ by $`\theta _W`$, and in strong phase by $`\theta _S`$. Then
$$A(bs\gamma )=A_{SM}+A_{\mathrm{New}}e^{i\theta _S}e^{i\theta _W};$$
$$A(\overline{b}\overline{s}\gamma )=A_{SM}+A_{\mathrm{New}}e^{i\theta _S}e^{i\theta _W}.$$
The $`(b/\overline{b})`$ averaged branching fraction, $``$, is
$$=\frac{1}{2}[|A(bs\gamma )|^2+|A(\overline{b}\overline{s}\gamma )|^2]A_{SM}^2(1+2\rho \mathrm{cos}\theta _S\mathrm{cos}\theta _W+\rho ^2),$$
where $`\rho =A_{\mathrm{New}}/A_{SM}`$. The CP asymmetry $`𝒜2\rho \mathrm{sin}\theta _S\mathrm{sin}\theta _W`$.
If one is sensitive to branching fraction differences of 20%, then one can detect New Physics amplitudes that are 10% of the Standard Model amplitude, if $`\theta _W`$ is near zero or 180 degrees, but cannot detect New Physics amplitudes smaller than 45% of the SM amplitude, if $`\theta _W`$ is near $`90^{}`$. For $`\theta _W`$ near $`90^{}`$, $`𝒜2\rho \mathrm{sin}\theta _S`$. So, if one were sensitive to CP asymmetries of 0.10, then one would have sensitivity to this New Physics for $`\rho \mathrm{sin}\theta _S>0.05`$.
So, there is a portion of New Physics parameter space, albeit small, where New Physics will show up as a CP asymmetry, but not as a branching fraction difference. This is discussed in general by A. Kagan and M. Neubert (hep-ph/9803368), and as applied to SUSY by Aoki, Cho, and Oshimo (hep-ph/9811251). Asymmetries in the 0.05-0.20 range are mentioned.
How might CLEO measure CP asymmetries in $`bs\gamma `$? By pseudoreconstruction! But wait a minute, didn’t I just say, in Section 6.2, that “In pseudoreconstruction, often one does not have the totally correct combination of particles (hence the ‘pseudo’), but this is not important, because the method is used only to suppress background $`\mathrm{}`$”? Well, yes. It still isn’t necessary to get the totally correct combination of particles, but it is necessary to get the flavor – $`b`$ or $`\overline{b}`$ – right. It turns out we get the flavor right about 92% of the time. It is straightforward to correct for the 8% mistake rate, a 19% scaling up of the measured asymmetry. With the 4.7$`fb^1`$ data sample used for the most recent branching fraction analysis, we obtain a corrected asymmetry of
$$𝒜=0.16\pm 0.14\pm 0.05$$
So, no evidence for CP violation, but errors that are uncomfortably large. The errors shown are statistical and systematic, in that order. With relatively little work, the systematic error can be reduced substantially, so even with 3 times the luminosity (our in-hand 14$`fb^1`$), the measurement will be statistics limited. An error of $`\pm 0.08`$ should be straightforward to achieve. We’re looking for ways to push that down, giving consideration to lepton tagging as a possibility.
## 7 Summary
In “the Early Days”, the basic features of the bottom quark were established:
* A left-handed doublet with a very heavy top.
* Decaying dominantly to charm, $`bcW_V^{}`$. Coupling to the second generation, $`|V_{cb}|0.04`$, smaller than the coupling between second generation and first, $`|V_{us}|=0.22`$.
* Decay to up, $`buW_V^{}`$, suppressed relative to decay to charm, but not zero.
In “Recent Times”, the emphasis is on testing the Standard Model, searching for New Physics. There are two approaches: measuring rates for electroweak penguins, and “overdetermining the CKM matrix”. Lets see where we now stand on each, and where we are going.
On electroweak penguins, the branching fraction for $`bs\gamma `$ has been measured to $`\pm `$17%, and is in good agreement with the Standard Model. In the near future, with data already in hand, the accuracy should be improved, to $`\pm `$10%. At that point, the error on the measurement will be about equal to the error on the theoretical prediction, and further progress will be slower in coming.
The CP asymmetry has been measured to an accuracy of $`\pm `$0.14, and should soon improve to $`\pm `$0.08. Further improvements are straightforward, as the error is purely statistical, and the large data samples to be accumulated by BaBar, Belle, and CLEO, in the next 3-4 years, should give a sensitivity to asymmetries in the 0.05 range.
The electroweak process $`bs\mathrm{}^+\mathrm{}^{}`$ has not yet been seen, though seeing it may not be far off. Very large data samples will be required to study the various distributions that this 3-body final state makes available. Possibly this is the role for hadron colliders.
Concerning “overdetermining the CKM matrix”, I would argue that the CKM matrix has now been “determined”. Figure 15 shows the famous unitarity triangle, the one from the unitarity condition obtained by multiplying the first column of the CKM matrix by the complex conjugate of the third column:
$$V_{ud}V_{ub}^{}+V_{cd}V_{cb}^{}+V_{td}V_{tb}^{}V_{ub}^{}\lambda |V_{cb}|+V_{td}=0$$
$`(7\text{-}1)`$
The base of the triangle, $`\lambda |V_{cb}|`$, is known to $`\pm `$10% (soon to be $`\pm `$4%). The left-hand leg, $`|V_{ub}|`$, is known to $`\pm `$25%. The one sigma error band for $`|V_{ub}|`$ is shown. The angle $`\gamma `$ probably lies between $`60^{}`$ and $`90^{}`$, the lower limit coming from the Hou, Smith, Würthwein analysis of CLEO data, the upper limit from $`B_s\overline{B}_s`$ mixing. I’ve shown more conservative limits on Fig. 15, $`45^{}`$ and $`110^{}`$. That, I claim, “determines” the CKM matrix. With that, one can predict the angle $`\beta `$ to be $`20^{}\pm 5^{}\pm 2^{}`$, where the first error comes from the uncertainty in $`|V_{ub}|`$, and the second from the uncertainty in $`\gamma `$. (Note that the uncertainty in the prediction of $`\beta `$ comes dominantly from $`|V_{ub}|`$, not $`\gamma `$.)
The first “overdetermination” of the CKM matrix comes from the CP violating parameter $`ϵ`$ in neutral kaon decay. The band it defines nicely intersects the allowed region.
The next “overdetermination” will be BaBar and Belle’s measurements of $`\beta `$. With 30$`fb^1`$ data samples, they expect to measure $`\beta `$ to $`\pm 5^{}`$. It will be interesting to see how their results compare with CLEO’s prediction of ($`20\pm 5)^{}`$, and also how their error on measured $`\beta `$ will compare with CLEO’s predictions based on improved measurement of $`|V_{ub}|`$. Interesting times ahead!
## 8 Acknowledgements
I have benefitted immeasurably from countless interactions and discussions with my collaborators in CLEO over the past two decades. I wish to thank them for this, and absolve them of any blame for my rash statements in this paper.
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# Measurement of Inclusive 𝐷_𝑠^± Photoproduction at HERA
## 1 Introduction
Inclusive $`D_s^\pm `$ photoproduction cross sections at HERA are presented for photon-proton centre-of-mass energies in the range $`130<W<280\text{GeV}`$. The $`D_s^\pm `$ mesons were reconstructed through the decay chain $`D_s^\pm \varphi \pi ^\pm (K^+K^{})\pi ^\pm `$. This analysis supplements recent measurements of inclusive photoproduction of $`D^\pm `$ mesons at HERA . The high-statistics measurement by the ZEUS collaboration was performed in the same $`W`$ range as given above. The measured cross sections were compared to next-to-leading order (NLO) calculations with fragmentation parameters extracted from $`D^\pm `$ production in $`e^+e^{}`$ annihilation. The experimental results were found generally to lie above the NLO expectations, in particular in the forward (proton) direction.
The study of $`D_s^\pm `$ photoproduction provides another test of perturbative QCD (pQCD) calculations of charm production which is experimentally independent of the $`D^\pm `$ measurement. Furthermore, from a ratio of the $`D_s^\pm `$ and $`D^\pm `$ cross sections the strangeness-suppression factor, $`\gamma _s`$, in charm fragmentation can be determined. A comparison of the cross-section ratio and $`\gamma _s`$ with those obtained in charm production in $`e^+e^{}`$ annihilation tests the universality of charm fragmentation.
## 2 Experimental Conditions
The measurements were performed at the HERA $`ep`$ collider in the ZEUS detector during 1996/1997. In this period HERA collided positrons with energy $`E_e=27.5\text{GeV}`$ and protons with energy $`E_p=820\text{GeV}`$. The integrated luminosity used in this analysis is $`38\text{pb}^1`$. A detailed description of the detector can be found elsewhere .
Charged particles were measured in the central tracking detector, CTD , which is a drift chamber consisting of $`72`$ concentric sense-wire layers covering the polar angle<sup>1</sup><sup>1</sup>1The ZEUS coordinate system is right-handed and has the nominal interaction point at $`X=Y=Z=0`$, with the $`Z`$-axis pointing in the proton beam direction and the $`X`$-axis horizontal, pointing towards the center of HERA. region $`15^{}<\theta <164^{}`$. The CTD operates in a magnetic field of $`1.43\text{T}`$ provided by a thin superconducting solenoid. The transverse momentum resolution for full length tracks is $`\sigma _p_{}/p_{}=0.0058p_{}0.00650.0014/p_{}`$ ($`p_{}`$ in GeV). To estimate the energy loss, $`dE/dx`$, of tracks, the truncated mean of the sense-wire pulse-heights was recorded for each track, discarding the $`10\%`$ lowest and $`30\%`$ highest pulses.
The solenoid is surrounded by the uranium–scintillator sampling calorimeter (CAL) , which is almost hermetic and consists of 5918 cells, each read out by two photomultipliers. Under test beam conditions, the CAL has a relative energy resolution of $`0.18/\sqrt{E}`$ ($`E`$ in GeV) for electrons and $`0.35/\sqrt{E}`$ for hadrons.
The luminosity was determined from the rate of the bremsstrahlung process $`e^+pe^+\gamma p`$, where the photon was measured by a lead–scintillator calorimeter located at $`Z=107\text{m}`$.
## 3 Event Selection and $`𝑫_𝒔^\mathbf{\pm }`$ Reconstruction
The ZEUS detector uses a three-level trigger system . At the first level (FLT), the calorimeter cells were combined to define regional and global sums that were required to exceed various CAL energy thresholds. At the second level (SLT), beam-gas events were rejected by cutting on the quantity $`\mathrm{\Sigma }_i(Ep_Z)_i>8\text{GeV}`$, where the sum runs over all calorimeter cells and $`p_Z`$ is the $`Z`$ component of the momentum vector assigned to each cell of energy $`E`$. At the third level (TLT), at least one combination of CTD tracks was required to be within wide mass windows around the nominal $`\varphi `$ and $`D_s`$ meson<sup>2</sup><sup>2</sup>2Here $`D^{}`$ and $`D_s`$ refer to $`D^\pm `$ and $`D_s^\pm `$, respectively. mass values, assuming $`\pi `$ and $`K`$ masses as appropriate. The $`D_s`$ transverse momentum, $`p_{}^{D_s}`$, was required to be greater than $`2.8\text{GeV}`$.
Photoproduction events were selected by requiring that no scattered positron was identified in the CAL . The Jacquet–Blondel estimator of $`W`$, $`W_{\text{JB}}=\sqrt{2E_p\mathrm{\Sigma }_i(Ep_Z)_i}`$, was required to be between 115 and 250 GeV. The lower limit was due to the SLT requirement, while the upper one suppressed remaining DIS events with an unidentified scattered positron in the CAL . After correcting for detector effects, the most important of which were energy losses in inactive material in front of the CAL and particle losses in the beam pipe , this $`W_{\text{JB}}`$ range corresponds to an interval of true $`W`$ of $`130<W<280\text{GeV}`$. Under these conditions, the photon virtuality, $`Q^2`$, is limited to values less than $`1\text{GeV}^2`$. The corresponding median $`Q^2`$ was estimated from a Monte Carlo (MC) simulation to be about $`3\times 10^4\text{GeV}^2`$.
The MC sample used for this analysis was prepared with the PYTHIA 6.1 generator. The proportions of direct- and resolved-photon events corresponded to the PYTHIA cross sections ($`50\%`$ each), where charm excitation processes were included in the resolved component . The MRSG and GRV-G HO parametrisations were used for the proton and photon structure functions, respectively. The MC events were processed through the standard ZEUS detector- and trigger-simulation programs and through the event reconstruction package used for offline data processing. The data and MC distributions were found to be in good agreement.
The $`D_s`$ mesons were reconstructed through the decay mode $`D_s^\pm \varphi \pi ^\pm `$, which has a branching ratio $`B_{D_s\varphi \pi }=0.036\pm 0.009`$ . The $`\varphi `$ was identified via its decay mode $`\varphi K^+K^{}`$, with $`B_{\varphi K^+K^{}}=0.491\pm 0.008`$ . The analysis was restricted to the pseudorapidity range $`1.5<\eta ^{D_s}<1.5`$, for which the CTD acceptance is high. Here $`\eta ^{D_s}\mathrm{ln}(\mathrm{tan}\frac{\theta }{2})`$, where $`\theta `$ is the polar angle with respect to the proton beam direction. The kinematic region in $`p_{}^{D_s}`$ was limited to $`3<p_{}^{D_s}<12\text{GeV}`$. The lower cut was required to comply with the $`p_{}^{D_s}`$ cut applied at the TLT. The upper cut was due to the limited statistics.
The three $`D_s`$ decay tracks were required to originate from the event vertex, which was measured with a resolution of 0.4 cm in the $`Z`$ direction and 0.1 cm in the $`XY`$ plane. Only tracks with polar angles $`20^{}<\theta <160^{}`$ and transverse momenta $`p_{}>0.75\text{GeV}`$ were considered in the analysis.
The decay of the pseudoscalar $`D_s`$ meson to the $`\varphi `$ (vector) plus $`\pi `$ (pseudoscalar) final state results in an alignment of the spin of the $`\varphi `$ meson with respect to the direction of motion of the $`\varphi `$ relative to the $`D_s`$. Consequently, the distribution of $`\mathrm{cos}\theta _K^{}`$, where $`\theta _K^{}`$ is the angle between one of the kaons and the pion in the $`\varphi `$ rest frame, followed a $`\mathrm{cos}^2\theta _K^{}`$ shape, implying a flat distribution for $`\mathrm{cos}^3\theta _K^{}`$. In contrast, the $`\mathrm{cos}\theta _K^{}`$ distribution of the combinatorial background was flat and its $`\mathrm{cos}^3\theta _K^{}`$ distribution peaked at zero. A cut of $`|\mathrm{cos}^3\theta _K^{}|>0.15`$ suppressed the background by a factor of approximately two while reducing the signal by 15%.
The $`dE/dx`$ information was used to allow partial $`K`$ and $`\pi `$ separation. For each kaon (pion) candidate a likelihood function $`\mathrm{}_{K(\pi )}\mathrm{exp}\{\frac{1}{2}[(dE/dx)_{\text{meas}}(dE/dx)_{K(\pi )}]^2/\sigma ^2\}`$ was determined, where $`(dE/dx)_{K(\pi )}`$ is the expected value for a kaon (pion), and $`\sigma `$ is the $`dE/dx`$ resolution, which is inversely proportional to $`\sqrt{n}`$, where $`n`$ is the number of hits entering the truncated mean. The parametrisation for the $`dE/dx`$ expectation was obtained from a fit to an independent inclusive track sample (Fig. 1a). A normalised likelihood function $`L_i\mathrm{}_i/_j\mathrm{}_j`$ was defined, where the sum extends over the considered particle hypotheses $`\pi `$, $`K`$ and $`p`$. Provided that the number of hits was sufficiently high ($`n>7`$), low-likelihood hypotheses were rejected if $`\mathrm{}_i<0.05`$, unless $`L_i>0.12`$. The $`dE/dx`$ cuts reduced the combinatorial background by approximately 20%. The signal loss was determined by means of a Monte Carlo simulation, using $`dE/dx`$ parameters obtained from the data. The overall loss was $`2.1\%`$.
For the $`D_s`$ candidates, $`p_{}^{D_s}/E_{}^{\theta >10^{}}>0.18`$ was required, where $`E_{}^{\theta >10^{}}`$ is the transverse energy outside a cone of $`\theta =10^{}`$ defined with respect to the proton direction. This cut removed more than $`20\%`$ of the background while preserving about $`95\%`$ of the $`D_s`$ signal, as verified by MC studies.
The $`K^+K^{}\pi ^\pm `$ mass distribution with the above cuts is shown in Fig. 1b for events in the $`\varphi `$ mass range, $`1.0115<M(K^+K^{})<1.0275\text{GeV}`$. The fraction of events with more than one entry in the $`D_s`$ mass region was less than 1%. The mass distribution was fitted to a sum of a Gaussian with the $`D_s`$ mass and width as free parameters, and an exponential function describing the non-resonant background. In order to avoid a possible contribution from $`D^\pm \varphi \pi ^\pm `$, the fit was not extended below $`1.895\text{GeV}`$. The fit yielded $`339\pm 48`$ $`D_s`$ mesons. The mass value obtained was $`M_{D_s}=1967\pm 2\text{MeV}`$, in agreement with the PDG value . The width of the signal was $`\sigma _{D_s}=12.5\pm 1.9\text{MeV}`$, in agreement with the MC estimation.
A clear $`\varphi `$ signal is seen in the $`M(K^+K^{})`$ distribution (Fig. 1c) for the $`D_s`$ region, $`1.94<M(K^+K^{}\pi ^\pm )<2.00\text{GeV}`$. The fit function for the $`\varphi `$ was a relativistic P-wave Breit-Wigner with variable mass and a fixed full-width of $`4.43\text{MeV}`$ , convoluted with a Gaussian function whose width, $`\sigma _\varphi `$, was a free parameter of the fit. The background was parametrised with the functional form $`a[M(K^+K^{})2m_K]^b`$, where $`m_K`$ is the $`K^\pm `$ mass. The fit yielded $`M(\varphi )=1019.5\pm 0.3\text{MeV}`$, in agreement with the PDG value , and $`\sigma _\varphi =1.7\pm 0.4\text{MeV}`$, in agreement with MC estimation. The number of $`\varphi `$ mesons originating from $`D_s`$ decays was estimated by a side-band subtraction, and was found to be in good agreement with the number of $`D_s`$ obtained from the above $`M(K^+K^{}\pi ^\pm )`$ fit.
## 4 Measurement of Inclusive $`𝑫_𝒔`$ Cross Sections
The inclusive $`D_s`$ cross section is given by:
$$\sigma _{epD_sX}=\frac{N_{D_s}}{AB_{D_s(\varphi K^+K^{})\pi }},$$
where $`N_{D_s}`$ is the fitted number of $`D_s`$ mesons, $``$ is the integrated luminosity, $`B_{D_s\left(\varphi K^+K^{}\right)\pi }=0.0177\pm 0.0044`$ is the combined $`D_s(\varphi K^+K^{})\pi `$ decay branching ratio and $`A`$ is the acceptance calculated with the MC sample (Section 3). The MC sample contains all events with the $`D_s(\varphi K^+K^{})\pi `$ decay channel as well as small admixtures from other $`D_s`$ decay modes and from other charm particle decays. Thus these contributions were taken into account in the acceptance correction procedure.
The total $`D_s`$ cross section in the kinematic region $`Q^2`$$`<1\text{GeV}^2`$, $`130<W<280\text{GeV}`$, $`3<p_{}^{D_s}<12`$ GeV and $`1.5<\eta ^{D_s}<1.5`$ was measured to be $`\sigma _{epD_sX}=3.79\pm 0.59\text{(stat.)}{}_{\mathrm{\hspace{0.17em}0.46}}{}^{+\mathrm{\hspace{0.17em}0.26}}\text{(syst.)}\pm 0.94\text{(br.)}\text{nb}`$, where the last error is due to the 25% uncertainty in the branching ratio $`B_{D_s\varphi \pi }`$. The differential cross sections $`d\sigma /dp_{}^{D_s}`$ and $`d\sigma /d\eta ^{D_s}`$ are given in Table 1.
In Fig. 2 they are compared with the distributions for the $`D^{}`$ in the same kinematic region. The overall normalisation uncertainty due to the luminosity measurement (1.7%) is not included in the cross section errors.
### 4.1 Systematic Uncertainties
A detailed study of possible sources of systematic uncertainties was carried out for all measured cross sections by shifting the nominal analysis parameters as described below, taking into account resolution effects. For each variation, except for the first one, which is due to the fit systematics, the $`D_s`$ mass and width were fixed to the values found for the nominal cut values. The following sources of systematic errors were considered:
* the uncertainties in the determination of the number of $`D_s`$ mesons were estimated by using a quadratic polynomial function for the background parametrisation instead of an exponential one and by varying the range of the $`K^+K^{}\pi ^\pm `$ mass distribution in the fit procedure;
* to estimate the uncertainties in the tracking procedure, the track selection cuts, including $`M(K^+K^{})`$ and $`\mathrm{cos}\theta _K^{}`$, were shifted by at least the expected resolutions from the nominal values (Section 3);
* the cut on $`p_{}^{D_s}/E_{}^{\theta >10^{}}`$ was changed by $`\pm 10\%`$;
* the MC simulation was found to reproduce the absolute energy scale of the CAL to within $`\pm 3\%`$ . The corresponding uncertainty in the cross section was calculated, including a shift of 3% in the FLT CAL energy thresholds (Section 3);
* the fraction of resolved photon processes in the PYTHIA MC sample was varied between 40% and 60%;
* the $`dE/dx`$ likelihood cuts were changed in the range $`0.02<l_i<0.10`$ and $`0.08<L_i<0.15`$;
* the uncertainty associated with the correction to the true $`W`$ range was determined by moving the $`W_{\text{JB}}`$ boundary values by the estimated resolution of $`\pm 7\%`$.
None of the above was dominant in the total systematic uncertainty to the inclusive $`D_s`$ cross section. All systematic errors were added in quadrature, yielding a total uncertainty of $`{}_{12}{}^{+7}\%`$, compared to a statistical error of 16%.
## 5 Comparison with QCD Calculations
The $`D_s`$ cross sections were compared to two types of pQCD calculations . The fractions of $`c`$ quarks hadronising as $`D^{}`$ and $`D_s`$ mesons were used as input to each calculation. The values $`f(cD^+)=0.235\pm 0.007\mathit{}0.007)`$ and $`f(cD_s^+)=0.101\pm 0.009\mathit{}0.025)`$ were extracted by a least-squares procedure from all relevant existing $`e^+e^{}`$ experimental data . The errors in brackets are due to uncertainties in the charm hadron decay branching ratios. They affect experimental and theoretical cross-section calculations in the same way and can be ignored in the comparison.
The NLO calculation of charm photoproduction in the fixed-order approach of Frixione et al. assumes that gluons and light quarks $`(u,d,s)`$ are the only active partons in the structure functions of the proton and the photon. In this approach there is no explicit charm excitation component, which can be important in charm photoproduction at HERA , and charm is only produced dynamically in hard pQCD processes. This calculation is expected to be valid when the $`c`$-quark transverse momentum, $`p_{}`$, is not much larger than the $`c`$-quark mass, $`m_c`$.
The structure function parametrisations used in the NLO calculation were MRSG for the proton and GRV-G HO for the photon. The renormalisation scale used was $`\mu _R=m_{}\sqrt{m_c^2+p_{}^2}`$, and the factorisation scales of the photon and proton structure functions were set to $`\mu _F=2m_{}`$. The pole mass definition is used in this calculation for the $`c`$-quark mass with a nominal value $`m_c=1.5\text{GeV}`$. The Peterson fragmentation function was used for charm fragmentation in this calculation. The Peterson parameter $`ϵ=0.035`$ was obtained for $`D^{}`$ in a NLO fit to ARGUS data . A recent NLO fit to ARGUS data yields an $`ϵ(D_s)`$ value equal to $`ϵ(D^{})`$ within the fit uncertainties. Using the same value for both channels leads to the same predictions, except for the difference in $`f(cD^+)`$ and $`f(cD_s^+)`$, which enter the calculation as scale factors. The NLO prediction for the total inclusive $`D_s`$ cross section of $`2.18\text{nb}`$ is smaller by $`2`$ standard deviations compared to the measured cross section. Scaling the Peterson parameter $`ϵ`$ with the squared ratio of the constituent quark masses , $`m_s=0.5\text{GeV}`$ and $`m_{u,d}=0.32\text{GeV}`$ , leads to $`ϵ(D_s)=0.085`$, which yields a NLO prediction 22% lower than that with $`ϵ(D_s)=0.035`$.
In Fig. 2, the NLO calculations are compared to the differential cross sections. The thick curves correspond to the nominal values of $`\mu _R`$ and $`m_c`$, as defined above. For the thin curves, a rather extreme value for the $`c`$-quark mass, $`m_c=1.2\text{GeV}`$, and a $`\mu _R`$ value of $`0.5m_{}`$ have been used. The $`D_s`$ cross section decreases steeply with rising $`p_{}`$, as in the $`D^{}`$ case. The NLO calculation reproduces within errors the shape of the $`p_{}^{D_s}`$ distribution but underestimates the data for the nominal parameter set. For the $`\eta ^{D_s}`$ distribution the NLO predictions are below the data in the central and forward regions. A similar effect was observed when $`d\sigma /d\eta ^D^{}`$ distributions were compared with various NLO predictions over a wide range of $`W`$ values and photon virtualities .
Recently Berezhnoy, Kiselev and Likhoded (BKL) have suggested another model which does not employ any specific fragmentation function. In this tree-level pQCD $`O(\alpha \alpha _s^3)`$ calculation, the $`(c,\overline{q})`$ state produced in pQCD is hadronised, taking into account both colour-singlet and colour-octet contributions. Using the experimental value for $`f(cD^+)`$, a $`c`$-quark mass $`m_c=1.5\text{GeV}`$, a light constituent quark mass $`m_q=0.3\text{GeV}`$ and the CTEQ4M proton structure function parametrisation , the ratio of the colour-octet and colour-singlet components was tuned to describe the ZEUS $`D^{}`$ photoproduction cross sections .<sup>3</sup><sup>3</sup>3A comparison of the $`D^{}`$ data with the BKL calculations can be found in the BKL paper . This ratio and the experimental value for $`f(cD_s^+)`$ were used to obtain predictions for $`D_s`$ photoproduction. In this case, a strange quark mass, $`m_s=0.5\text{GeV}`$, was used instead of $`m_q`$. The calculated cross sections for $`D_s`$ and $`D_{s}^{}{}_{}{}^{\pm }`$ were combined, since the inclusive $`D_s`$ channel includes fully the prompt $`D_{s}^{}{}_{}{}^{\pm }`$ meson production.
In Fig. 3,
the BKL calculations are compared to the $`D_s`$ differential cross-section measurements. The agreement with the data is better than that of the NLO calculation with the nominal parameters, but the shape of the $`\eta ^{D_s}`$ distribution is not matched by the BKL prediction. The total predicted $`D_s`$ cross section in the kinematic range of the measurement ($`3.37\text{nb}`$) is compatible with the measured inclusive cross section.
## 6 $`𝑫_𝒔`$ to $`𝑫^{\mathbf{}}`$ cross-section ratio and $`𝜸_𝒔`$
In the $`D_s`$ kinematic region, as defined in Section 4, the $`D^{}`$ cross section was measured to be $`\sigma _{epD^{}X}=9.17\pm 0.35\text{(stat.)}{}_{\mathrm{\hspace{0.17em}0.39}}{}^{+\mathrm{\hspace{0.17em}0.40}}\text{(syst.)}\text{nb}`$ . This yields a ratio $`\sigma _{epD_sX}/\sigma _{epD^{}X}=0.41\pm 0.07\text{(stat.)}_{0.05}^{+0.03}\text{(syst.)}\pm 0.10\text{(br.)}`$, where common systematic errors ($`W_{JB}`$ and CAL energy scale) have been removed and the last error is the uncertainty in $`B_{D_s\varphi \pi }`$. This ratio is in good agreement with the ratio $`f(cD_s^+)/f(cD^+)=0.43\pm 0.04\pm 0.11\text{(br.)}`$ obtained from results of $`e^+e^{}`$ experiments (see Section 5). The last error originates from the uncertainty in $`B_{D_s\varphi \pi }`$ and can be ignored in the comparison.
The strangeness-suppression factor, $`\gamma _s`$, is the ratio of probabilities to create $`s`$ and $`u,d`$ quarks during the fragmentation process. In simulation programs based on the Lund string fragmentation scheme , $`\gamma _s`$ is a free parameter with a default value of $`0.3`$. By varying the value of $`\gamma _s`$ in the PYTHIA simulation , a direct relation can be obtained between $`\gamma _s`$ and the $`D_s`$ to $`D^{}`$ cross-section ratio. Fixing the value of $`f(cD^+)`$ in PYTHIA to 0.235 yields
$$\gamma _s=0.27\pm 0.04\text{(stat.)}{}_{\mathrm{\hspace{0.17em}0.03}}{}^{+\mathrm{\hspace{0.17em}0.02}}\text{(syst.)}\pm 0.01\text{(frac.)}\pm 0.07\text{(br.)}.$$
The third error is due to the uncertainty in $`f(cD^+)`$, while the forth one results from the uncertainty in $`B_{D_s\varphi \pi }`$. Adding all errors in quadrature, except the last one, gives $`\gamma _s=0.27\pm 0.05\pm 0.07\text{(br.)}`$.
Previously, $`\gamma _s`$ was measured mainly from the ratio of $`K`$ to $`\pi `$ production and from the momentum spectrum of $`K`$ mesons in hadron-hadron and $`e^+e^{}`$ collisions, as well as in deep inelastic scattering (DIS) experiments . The most accurate measurement, obtained in $`p\overline{p}`$ collisions , is $`\gamma _s=0.29\pm 0.02\text{(stat.)}\pm 0.01\text{(syst.)}`$. The DIS results require a lower value, $`\gamma _s0.2`$. Recent results from $`e^+e^{}`$ collisions are in some disagreement with each other. The SLD preliminary result is $`\gamma _s=0.26\pm 0.06`$, while OPAL finds $`\gamma _s=0.422\pm 0.049\pm 0.059`$. Previous $`\gamma _s`$ measurements are therefore in good agreement with that of this analysis, except the latest OPAL value, which is about 2 standard deviations higher.
Existing $`\gamma _s`$ values obtained from heavy-meson production (charm and beauty) in $`e^+e^{}`$ collisions are centred around 0.3. For charm production, the ratio $`2f(cD_s^+)/[f(cD^+)+f(cD^0)]`$ was used as a measure of $`\gamma _s`$. Using for the above fractions the more recent values quoted in leads to $`\gamma _s=0.26\pm 0.03\pm 0.07\text{(br.)}`$, where the latter uncertainty originates from $`B_{D_s\varphi \pi }`$ and can be ignored in the comparison with the ZEUS result.
The results presented here on the $`D_s`$ to $`D^{}`$ cross-section ratio and on $`\gamma _s`$, taken together with charm production data in $`e^+e^{}`$ annihilation, tend to support the universality of charm fragmentation.
## 7 Summary and Conclusions
The first measurement at HERA of inclusive $`D_{s}^{}{}_{}{}^{\pm }`$ photoproduction has been performed with the ZEUS detector. The cross section for $`Q^2<1\text{GeV}^2`$, $`130<W<280\text{GeV}`$, $`3<p_{}^{D_s}<12\text{GeV}`$ and $`1.5<\eta ^{D_s}<1.5`$ is $`\sigma _{epD_sX}=3.79\pm 0.59\text{(stat.)}{}_{\mathrm{\hspace{0.17em}0.46}}{}^{+\mathrm{\hspace{0.17em}0.26}}\text{(syst.)}\pm 0.94\text{(br.)}\text{nb}`$. The differential cross sections $`d\sigma /dp_{}^{D_s}`$ and $`d\sigma /d\eta ^{D_s}`$ are generally above fixed-order NLO calculations, as was the case with the results previously obtained for $`D^{}`$ photoproduction in the same kinematic region. The BKL calculation, using the octet-to-singlet ratio tuned to fit the ZEUS $`D^{}`$ data, describes the $`D_s`$ cross sections reasonably well, but the shape of the $`\eta ^{D_s}`$ distribution is not matched by the BKL prediction. The cross-section ratio, $`\sigma _{epD_sX}/\sigma _{epD^{}X}`$, in the kinematic region as defined above is $`0.41\pm 0.07\text{(stat.)}_{0.05}^{+0.03}\text{(syst.)}\pm 0.10\text{(br.)}`$, in good agreement with the ratio of $`c`$ quarks hadronising into $`D_s`$ and $`D^{}`$ mesons, extracted from $`e^+e^{}`$ experiments. From this ratio, the strangeness-suppression factor in charm photoproduction, within the LUND string fragmentation model, has been calculated to be $`\gamma _s=0.27\pm 0.04\text{(stat.)}{}_{\mathrm{\hspace{0.17em}0.03}}{}^{+\mathrm{\hspace{0.17em}0.02}}\text{(syst.)}\pm 0.01\text{(frac.)}\pm 0.07\text{(br.)}`$, in good agreement with the $`\gamma _s`$ value extracted from charm production in $`e^+e^{}`$ annihilation.
## 8 Acknowledgements
We would like to thank the DESY Directorate for their strong support and encouragement. The remarkable achievements of the HERA machine group were essential for the successful completion of this work and are greatly appreciated. We would like to thank C. Oleari for discussions and for providing us with his latest results for $`ϵ(D_s)`$ and S. Frixione and A. Berezhnoy for providing us with their QCD calculations.
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# Quantum spinor field in the FRW universe with a constant electromagnetic background
## I Introduction
The article is a natural continuation of our paper . Here we generalize the latter consideration to the case of a massive spinor field which is placed in Friedmann-Robertson-Walker (FRW) Universe of special type with constant electromagnetic field. First, we ought to repeat briefly a motivation for such an activity.
It is quite well known fact that quantum field theory in an external background is, generally speaking, theory with unstable vacuum. The vacuum instability leads to many interesting features, among which particle creation from the vacuum is one of the most beautiful non-perturbative phenomenon. One has to treat it exactly in regard to the external field. The latter has been realized long ago by Schwinger . The particle creation effects together with back reaction issue are important in black hole physics and in dynamics of the early Universe (EU) (see, for example, and references therein).
In quantum field theory with unstable vacuum it is necessary to construct different kinds of Green functions (GF), e.g. besides the causal GF (out-in GF) one has to use so called in-in GF, out-out GF, and so on (for a review and technical details see ). General methods of such GF construction in electromagnetic background have been developed in . A possible generalization of the formalism to an external gravitational background has been given in Ref. . Since Ref. it has been known that causal (out-in) GF may be presented as a proper-time integral over a real infinite contour. At the same time, in the unstable vacuum case the in-in and out-out GF differ from the causal one. It was shown there are examples of external fileds (electromagnetic with constant uniform invariants) when these functions may be presented by the same proper-time integrals (with the same integrand) but over another contours in the complex proper-time plane. Then, it is not difficult to compare contributions from the in-in GF and from the causal one. The complete set of GF mentioned is necessary for the construction of the Furry picture in interacting theories, and even in noninteracting cases one has to use them to define, for example, the back reaction of particles created and to construct different kinds of effective actions (EA) (for a general introduction to EA in background field method, see , and for review of modern generalizations, see ). The such proper-time representation of GF may be the necessary step in the study of chiral symmetry breaking in QED and the four-fermion models under the action of gravitational and electromagnetic fields (see and references therein). ¿From another point, in-in GF which gives the origin to in-in EA maybe used in more realistic theories, like GUT theories with scalars, spinors and vectors in order to analyse the properties of above electro-gravitational background in the EU. For example, one of extremely interesting questions there is: can we realise the asymptotic conformal invariance phenomenon (which means that theory becomes approximately conformally invariant at large curvature) even for in-in EA, or in other words for mean values in EU. Taking into account our recent study of in-in GF structure for scalars it looks quite interesting next application of above calculation.
It may be likely that EU is filled with some type of electromagnetic fields. For example, recently (see and references therein) the possibility of existence and role of primordial magnetic fields in EU have been discussed. From another point the possibility of existence of electromagnetic field in the EU has been discussed long ago in . It has been shown there that the presence of the electrical field in the EU increases significantly the gravitational particle creation from the vacuum. In principle, this process may be considered as a source for the dominant part of the Universe mass.
Bearing in mind the above cosmological motivations it is becoming interesting to study the quantum field theory in curved background with electromagnetic field (of a special form in order to solve the problem analytically). In the present paper we are going to consider a massive spinor field placed in the expanding FRW Universe with the scale factor $`\mathrm{\Omega }(\eta )`$ (in terms of the conformal time) $`\mathrm{\Omega }^2(\eta )=b^2\eta ^2+a^2`$. Such a scale factor corresponds to the expanding radiation-dominated FRW Universe. In terms of physical time $`t`$ the corresponding metric may be written as follows:
$$ds^2=dt^2\mathrm{\Omega }^2(t)(dx^2+dy^2+dz^2),$$
(1)
where for small times $`|t|a^2/b,\mathrm{\Omega }^2(t)a^2[1+(bt/a^2)^2]`$, and for large times $`|t|a^2/b,\mathrm{\Omega }^2(t)2b|t|`$ (see ). Moreover, such FRW Universe will be filled by the constant electromagnetic field.
Thus, we start from the theory of massive spinor in above background. Making a conformal transformation we remain with QED in flat background but with time-dependent mass (QED-$`\mathrm{\Omega }`$ theory). In the Sect.II special sets of exact solutions of Dirac equation in QED-$`\mathrm{\Omega }`$ theory are constructed and classified as corresponding to particles and antiparticles at $`t\pm \mathrm{}`$. In the Sect.III, using these solutions, representations for out-in, in-in, and out-out spinor GF are explicitly constructed as proper-time integrals over the corresponding contours in complex proper-time plane. As far as we know, it is a first explicit example for the proper-time representations for complete set of spinor GF in gravitational-electromagnetic background. In the Sect.IV we are interested in to reveal global features of the theory. The vacuum-to-vacuum transition amplitudes and number of created particles are found and vacuum instability is discussed. It is seen the creation process is a coherent effect of both fields. The all mean values of the current and energy-momentum tensor are presented in the same manner as the proper-time integrals, and evaluated. The different approximations for them are investigated. The back reaction produced by both of particles created from a vacuum and polarization of an unstable vacuum estimated in different regimes. It is shown a behaviour of such components in time are quite different.
## II Classified sets of exact solutions
In this Section we study exact solutions of the Dirac equation in an external constant uniform electromagnetic background and in a time-dependent mass-like potential, which effectively reproduces effects of a gravitational background (solutions of the Dirac equation of QED-$`\mathrm{\Omega }`$ theory),
$`(𝒫_\mu \gamma ^\mu M\mathrm{\Omega })\psi \left(x\right)=0,\mathrm{\Omega }=\mathrm{\Omega }(x_0)=\sqrt{a^2+b^2x_0^2},`$ (2)
$`𝒫_\mu =i_\mu qA_\mu \left(x\right),[\gamma ^\mu ,\gamma ^\nu ]_+=2\eta ^{\mu \nu },\eta ^{\mu \nu }=\mathrm{diag}(1,1,1,1),`$ (3)
where $`x^0=\eta `$ is conformal time, $`q`$ is charge of a particle, for example,$`q=|e|`$ for electron. The time-independent inner product of the solutions of the equation (2) may be chosen as
$$(\psi ,\psi ^{})=\overline{\psi }(x)\gamma ^0\psi ^{}(x)𝑑𝐱.$$
(4)
As usual, it is convenient to present $`\psi (x)`$ in the following form
$$\psi (x)=\left(𝒫_\mu \gamma ^\mu +M\mathrm{\Omega }\right)\varphi (x).$$
(5)
Then the functions $`\varphi `$ have to obey the squared Dirac equation,
$`\left(𝒫^2\left(M\mathrm{\Omega }\right)^2{\displaystyle \frac{q}{2}}\sigma ^{\mu \nu }F_{\mu \nu }+iM_0\mathrm{\Omega }\gamma ^0\right)\varphi (x)=0,`$ (6)
$`F_{\mu \nu }=_\mu A_\nu (x)_\nu A_\mu (x),\sigma ^{\mu \nu }={\displaystyle \frac{i}{2}}[\gamma ^\mu ,\gamma ^\nu ].`$ (7)
The external electromagnetic field in our case consists of a constant uniform electric ($`E`$) and parallel to it magnetic ($`H`$) fields,
$$F_{03}=E,F_{\mu \nu }^{}=H\left(\delta _\mu ^2\delta _\nu ^1\delta _\nu ^2\delta _\mu ^1\right).$$
(8)
For such a field we select the following potentials:
$$A_0=0,A_3=Ex^0,A_i=A_i^{}=Hx_2\delta _i^1,i=1,2.$$
(9)
For $`b=0`$ we have a usual flat-space case with the mass $`m=aM`$. In this case particle-antiparticle classified solutions of the equation and all the GF were found in . The case $`b0`$ is of special interest for us. In the case one can consider spinor field in the conformally-flat Universe (with scale factor $`\mathrm{\Omega }`$) filled by a constant uniform electromagnetic field. Making a standard conformal transformation of the gravitational metric and spinor field, we arrive to a theory in flat space-time with a time-dependent mass (QED-$`\mathrm{\Omega }`$). The corresponding field equation is given by (2) (an electromagnetic field should not be transformed under the conformal transformation). Note that the such a conformal transformation may be used also for interacting theories . Thus, Eq. (2) is actually relevant to the quantum spinor field in the expanding FRW Universe with the external constant electromagnetic field.
Remember that the expansion law with $`a0`$ is necessarily connected with a nonzero energy density of the cosmological substratum in the early, radiation-dominated phase of the Universe. The minimum value of $`\mathrm{\Omega }`$ ($`\mathrm{\Omega }_{min}=a`$) is caused by the strong interaction. It is difficult to find solutions of the equation (6) with the last spin term. Nevertheless, in many interesting cases with strong background part of this term (for region of small $`(x_0)^2(a/b)^2`$) may be treated perturbatively starting from the reduced squared Dirac equation:
$$\left[𝒫^2\left(M\mathrm{\Omega }\right)^2\frac{q}{2}\sigma ^{\mu \nu }F_{\mu \nu }+ibM\gamma ^0\right]\varphi (x)=0,$$
(10)
where the asymptotic form of the last spin term is used ($`_0\mathrm{\Omega }b`$ for large enough $`x_0`$). If the intensity of an electric field is more than the characterizing parameter of a mass-like potential, $`(qE)^2(bM)^2`$, or its potential is not intense, $`a^2M/b1`$, then both the last spin term of equation (6) and of the equation (10) can be disregarded. If $`(qE)^2(bM)^2,`$ one can not neglect the contribution from the term $`iM_0\mathrm{\Omega }\gamma ^0`$. However, our main interest is to study an intense gravitational background when $`a^2M/b1`$. Then, as it may be seen from explicit solutions of the equation (10), one can calculate any spin contributions to matrix elements of operators using solutions of equation (10) with a relative accuracy of the order of $`a^2M/b.`$ And one may use these sets as a basis to construct solutions of the equation (6) perturbatively, considering the term $`iM(_0\mathrm{\Omega }b)\gamma ^0`$ as a perturbation. Moreover, the large time ($`(x_0)^2(a/b)^2`$ ) asymptotic solutions of both equations (6) and (10) are the same.
To construct the above mentioned generalized Furry picture for both QED and QED-$`\mathrm{\Omega }`$ one has to find special sets of classified solutions of the equation (2), namely, two complete and orthonormal sets of solution: $`\left\{{}_{\pm }{}^{}\psi _{\left\{n\right\}}^{}(x)\right\}`$, which describes particles (+) and antiparticles ($``$) in the initial time instant ($`x^0\mathrm{})`$, and $`\left\{{}_{}{}^{\pm }\psi _{\left\{n\right\}}^{}(x)\right\}`$, which describes particles (+) and antiparticles ($``$) in the final time instant ($`x^0+\mathrm{}).`$ According to the general approach , which can be easily adapted to QED-$`\mathrm{\Omega }`$, such solutions obey the following asymptotic conditions
$`H_{o.p.}(x^0){}_{\zeta }{}^{}\psi _{\left\{n\right\}}^{}(x)={}_{\zeta }{}^{}\epsilon {}_{\zeta }{}^{}\psi _{\left\{n\right\}}^{}(x),,\mathrm{sgn}{}_{\zeta }{}^{}\epsilon =\zeta ,x^0\mathrm{},`$ (11)
$`H_{o.p.}(x^0){}_{}{}^{\zeta }\psi _{\left\{l\right\}}^{}(x)={}_{}{}^{\zeta }\epsilon {}_{}{}^{\zeta }\psi _{\left\{l\right\}}^{}(x),\mathrm{sgn}{}_{}{}^{\zeta }\epsilon =\zeta ,x^0+\mathrm{},\zeta =\pm ,`$ (12)
where $`\zeta ,\left\{n\right\}`$ and $`\zeta ,\left\{l\right\}`$ are complete sets of quantum numbers which characterize solutions $`{}_{\zeta }{}^{}\psi _{\left\{n\right\}}^{}(x)`$ and $`{}_{}{}^{\zeta }\psi _{\left\{l\right\}}^{}(x)`$ respectively, $`H_{o.p.}\left(x^0\right)=\gamma ^0(M\mathrm{\Omega }\gamma ^i𝒫_i)`$ is one-particle Dirac Hamiltonian; $`{}_{}{}^{+}\epsilon `$, $`{}_{+}{}^{}\epsilon `$ are particle quasi-energies and $`|^{}\epsilon |`$ and $`|_{}\epsilon |`$ are antiparticles quasi-energies. All the information about the processes of particles scattering and creation by an external field (in zeroth order with respect to the radiative corrections) can be extracted from the decomposition coefficients, which form the matrices $`G({}_{\zeta }{}^{}|{}_{}{}^{\zeta ^{}})`$,
$${}_{}{}^{\zeta }\psi (x)={}_{+}{}^{}\psi (x)G({}_{+}{}^{}|{}_{}{}^{\zeta })+{}_{}{}^{}\psi (x)G({}_{}{}^{}|{}_{}{}^{\zeta }).$$
(13)
The matrices $`G({}_{\zeta }{}^{}|{}_{}{}^{\zeta ^{}})`$ obey the following relations,
$`G({}_{\zeta }{}^{}|{}_{}{}^{+})G({}_{\zeta }{}^{}|{}_{}{}^{+})^{}+G({}_{\zeta }{}^{}|{}_{}{}^{})G({}_{\zeta }{}^{}|{}_{}{}^{})^{}=𝐈,`$ (14)
$`G({}_{+}{}^{}|{}_{}{}^{+})G({}_{}{}^{}|{}_{}{}^{+})^{}+G({}_{+}{}^{}|{}_{}{}^{})G({}_{}{}^{}|{}_{}{}^{})^{}=0,`$ (15)
where $`𝐈`$ is the identity matrix. All GF in the formalism may be also constructed with the sets $`\left\{{}_{\pm }{}^{}\psi _{\left\{n\right\}}^{}(x)\right\}`$ and $`\left\{{}_{}{}^{\pm }\psi _{\left\{n\right\}}^{}(x)\right\}`$. Below we are going to present such solutions.
The functions $`\varphi (x)`$ can be written in the following form:
$$\varphi _{p_3p_1n\xi r}(x)=\varphi _{p_3n\xi r}(x_{})\varphi _{p_1nr}(x_{})v_{\xi r},$$
(16)
where $`x_{}^\mu =(0,x^1,x^2,0),x_{}^\mu =(x^0,0,0,x^3)`$; $`\{p_3,p_1,n,\xi ,r\}`$ is a complete set of quantum numbers. Among them $`p_3`$ and $`p_1`$ are momenta of the continuous spectrum, $`n`$ is an integer quantum number, $`\xi =\pm 1`$ and $`r=\pm 1`$ are spin quantum numbers; $`v_{\xi r}`$ are some constant orthonormal spinors, $`v_{\xi r}^{}v_{\xi r^{}}=\delta _{rr^{}}.`$ The eq.(10) allows one to subject these spinors to some supplementary conditions,
$`\mathrm{\Xi }v_{\xi r}=\xi v_{\xi r},\mathrm{\Xi }=\gamma ^0(qE\gamma ^3bM)/\rho ,\rho =\sqrt{(qE)^2+(bM)^2},`$ (17)
$`Rv_{\xi r}=rv_{\xi r},R=\text{sgn}\left(qH\right)i\gamma ^1\gamma ^2.`$ (18)
If $`H0,`$ the function $`\varphi _{p_1nr}(x_{})`$ has the form
$`\varphi _{p_1nr}(x_{})=\left({\displaystyle \frac{\sqrt{|qH|}}{2^{n+1}\pi ^{\frac{3}{2}}n!}}\right)^{1/2}\mathrm{exp}\left\{ip_1x^1{\displaystyle \frac{X^2}{2}}\right\}_n\left(X\right),X=\sqrt{|qH|}\left(x^2+{\displaystyle \frac{p^1}{qH}}\right),`$
where $`_n(x)`$ are Hermite polynomials with integer $`n=0,1,\mathrm{}`$. If $`H=0,`$ the discrete quantum number $`n`$ has to be replaced by the momentum $`p_2,`$ and the corresponding function has the form $`\varphi _{p_1nr}(x_{})=(2\pi )^1\mathrm{exp}\left\{i\left(p_1x^1+p_2x^2\right)\right\}.`$ Let us present the function $`\varphi (x_{})`$ as follows
$$\varphi _{p_3n\xi r}(x_{})=(2\pi )^{1/2}e^{ip_3x^3}\varphi _{p_3n\xi r}(x^0),\varphi _{p_3n\xi r}(x^0)=\varphi _{p_3n\xi r}(x^0,p_z)|_{p_z=0}$$
(19)
where $`\varphi _{p_3n\xi r}(x^0,p_z)`$ is a solution of equation
$$\left[\left(i\frac{}{\stackrel{~}{\eta }}\right)^2\left(p_z\rho \stackrel{~}{\eta }\right)^2\rho \lambda i\rho \xi \right]\varphi _{p_3n\xi r}(x^0,p_z)=0,$$
(20)
with $`\stackrel{~}{\eta }=x^0\rho ^2qEp_3,\rho \lambda =p_3^2(bM/\rho )^2+\omega +a^2M^2,`$
$`\omega =\{\begin{array}{cc}|qH|(2n+1r),n=0,1,\mathrm{},\hfill & H0\hfill \\ p_1^2+p_2^2,\hfill & H=0\hfill \end{array}.`$
One can form two complete sets $`\left\{{}_{\pm }{}^{}\varphi _{p_3n\xi r}^{}(x^0,p_z)\right\}`$ and $`\left\{{}_{}{}^{\pm }\varphi _{p_3n\xi r}^{}(x^0,p_z)\right\}`$ of the solutions of equation (20) using the functions
$`{}_{+}{}^{}\varphi _{p_3n\xi r}^{}(x^0,p_z)=C_\xi D_{\nu \xi /2}[\pm (1i)\tau ],\tau ={\displaystyle \frac{1}{\sqrt{\rho }}}\left(\rho \stackrel{~}{\eta }p_z\right),`$ (21)
$`{}_{}{}^{+}\varphi _{p_3n\xi r}^{}(x^0,p_z)=C_\xi ^{}D_{\nu 1+\xi /2}\left[\pm (1+i)\tau \right],\nu ={\displaystyle \frac{i\lambda }{2}}{\displaystyle \frac{1}{2}}.`$ (22)
Similar solutions were first presented in . Then, solutions of equation (10) $`\varphi (x)`$ can be constructed as follows,
$`{}_{\pm }{}^{}\varphi _{p_3p_1n\xi r}^{}(x)`$ $`=`$ $`{}_{\pm }{}^{}\varphi _{p_3p_1n\xi r}^{}(x,p_z)|_{p_z=0},`$ (23)
$`{}_{\pm }{}^{}\varphi _{p_3p_1n\xi r}^{}(x,p_z)`$ $`=`$ $`(2\pi )^{1/2}e_\pm ^{ip_3x^3}\varphi _{p_3n\xi r}(x^0,p_z)\varphi _{p_1nr}(x_{})v_{\xi r},`$ (24)
and in the same form with $`(\pm )`$ indices above.
One can verify that the solutions of the Dirac equation with different $`\xi `$, namely, $`\left(𝒫_\mu \gamma ^\mu +M\mathrm{\Omega }\right){}_{\pm }{}^{}\varphi _{p_3,p_1,n,+1,r}^{}(x)`$ and $`\left(𝒫_\mu \gamma ^\mu +M\mathrm{\Omega }\right){}_{\pm }{}^{}\varphi _{p_3,p_1,n,1,r}^{}(x),`$ or $`\left(𝒫_\mu \gamma ^\mu +M\mathrm{\Omega }\right){}_{}{}^{\pm }\varphi _{p_3,p_1,n,+1,r}^{}(x)`$ and $`\left(𝒫_\mu \gamma ^\mu +M\mathrm{\Omega }\right){}_{}{}^{\pm }\varphi _{p_3,p_1,n,1,r}^{}(x)`$ are linearly dependent for each sign ”$`+`$” or ”$``$”. Thus, to construct the complete sets we may use only the following sets of solutions:
$`{}_{\pm }{}^{}\psi _{p_3p_1nr}^{}(x)`$ $`=`$ $`\left(𝒫_\mu \gamma ^\mu +M\mathrm{\Omega }\right)_\pm \varphi _{p_3,p_1,n,+1,r}(x),`$ (25)
$`{}_{}{}^{\pm }\psi _{p_3p_1nr}^{}(x)`$ $`=`$ $`\left(𝒫_\mu \gamma ^\mu +M\mathrm{\Omega }\right)^\pm \varphi _{p_3,p_1,n,+1,r}(x).`$ (26)
Choosing the coefficients $`C`$ and $`C^{}`$ in (22) as follows: $`C_{+1}=(2\rho )^{1/2}\mathrm{exp}\left(\pi \lambda /8\right)`$ and $`C_{+1}^{}=(\rho \lambda )^{1/2}\mathrm{exp}\left(\pi \lambda /8\right),`$ one gets two complete sets $`\{{}_{\pm }{}^{}\psi _{p_3p_1nr}^{}(x\}`$ and $`\left\{{}_{}{}^{\pm }\psi _{p_3p_1nr}^{}(x)\right\}`$ of orthonormalized solutions of the equation (2). These solutions are classified as particles (+) and antiparticles ($``$) at $`x^0\pm \mathrm{}`$ according to the asymptotic forms of the corresponding quasienergies, $`{}_{\zeta }{}^{}\epsilon =\zeta \rho |x^0|`$ and $`{}_{}{}^{\zeta }\epsilon =\zeta \rho |x^0|`$ (see for additional arguments advocating such a classification). It matches with classification of similar solutions in QED.
According to the above discussion the solutions (25) and (26) of the Dirac equation may serve in an intense gravitational background, $`a^2M/b1`$, and in the other cases mentioned after Eq. (10). Satisfying the Cauchy conditions one can see the solutions (25) are valid with $`x^0<0`$, $`|x^0|a/b`$ and the solutions (26) are valid with $`x^0>0`$, $`|x^0|a/b`$ for any intensity of the background.
To find the matrices $`G({}_{\zeta }{}^{}|{}_{}{}^{\zeta ^{}})`$ defined by (13), it is convenient to use an asymptotic form of the solutions. Using (10) and (18) we get
$$G({}_{\zeta }{}^{}|{}_{}{}^{\zeta ^{}})_{ll^{}}=\delta _{l,l^{}}g({}_{\zeta }{}^{}|{}_{}{}^{\zeta ^{}}),l=(p_3,p_1,n,r),l^{}=(p_3{}_{}{}^{},p_1{}_{}{}^{},n^{},r^{}),$$
(27)
where
$$g(_\zeta |^\zeta ^{})={}_{\zeta }{}^{}\varphi _{p_3,n,+1,r}^{}(x^0,p_z)i\underset{0}{\overset{}{}}(i_0\rho \stackrel{~}{\eta }){}_{}{}^{\zeta ^{}}\varphi _{p_3,n,+1,r}^{}(x^0,p_z).$$
(28)
## III Green functions
Let us start with out-in GF which is the causal propagator
$$S^c(x,x^{})=c_v^1i<0,out|T\psi (x)\overline{\psi }(x^{})|0,in>,c_v=<0,out|0,in>.$$
(29)
Here $`\psi (x)`$ is quantum spinor field satisfying the Dirac equation (2), $`|0,in>`$ and $`|0,out>`$ are the initial and the final vacuum, and $`c_v`$ is the vacuum-to-vacuum transition amplitude. The propagator $`S^c(x,x^{})`$ obeys the equation
$$\left(𝒫_\mu \gamma ^\mu M\mathrm{\Omega }\right)S^c(x,x^{})=\delta ^{(4)}(xx^{}).$$
(30)
Another important singular function is the commutation function $`S(x,x^{})=i[\psi (x),\overline{\psi }(x^{})]_+.`$ It obeys the homogeneous Dirac equation (2) and the initial condition $`S(x,x^{})|_{x_0=x_0^{}}=i\gamma ^0\delta (𝐱𝐱^{}).`$ Besides, the following GF are studied :
$`S_{in}^c(x,x^{})=i<0,in|T\psi (x)\overline{\psi }(x^{})|0,in>,S_{in}^{\overline{c}}(x,x^{})=i<0,in|\psi (x)\overline{\psi }(x^{})T|0,in>,`$ (31)
$`S_{in}^{}(x,x^{})=i<0,in|\psi (x)\overline{\psi }(x^{})|0,in>,S_{in}^+(x,x^{})=i<0,in|\overline{\psi }(x^{})\psi (x)|0,in>,`$ (32)
$`S_{out}^c(x,x^{})=i<0,out|T\psi (x)\overline{\psi }(x^{})|0,out>.`$ (33)
Here the $`T`$-product acts on both sides: it orders the field operators to the right of its and antiorders them to the left. The functions $`S_{in}^c`$ and $`S_{out}^c`$ obey the equation (30), $`S^{}`$ satisfy the equation (2), and $`S_{in}^{\overline{c}}`$ obeys the equation $`\left(𝒫_\mu \gamma ^\mu M\mathrm{\Omega }\right)S_{in}^{\overline{c}}(x,x^{})=\delta ^{(4)}(xx^{}).`$
One can express the GF via the solutions (25) and (26) :
$`S^c(x,x^{})`$ $`=`$ $`\theta \left(x_0x_0^{}\right)S^{}(x,x^{})\theta \left(x_0^{}x_0\right)S^+(x,x^{}),`$ (34)
$`S(x,x^{})`$ $`=`$ $`S^{}(x,x^{})+S^+(x,x^{}),`$ (35)
$`S_{in}^c(x,x^{})`$ $`=`$ $`\theta \left(x_0x_0^{}\right)S_{in}^{}(x,x^{})\theta \left(x_0^{}x_0\right)S_{in}^+(x,x^{}),`$ (36)
$`S_{in}^{\overline{c}}(x,x^{})`$ $`=`$ $`\theta \left(x_0^{}x_0\right)S_{in}^{}(x,x^{})\theta \left(x_0x_0^{}\right)S_{in}^+(x,x^{}),`$ (37)
$`S_{out}^c(x,x^{})`$ $`=`$ $`\theta \left(x_0x_0^{}\right)S_{out}^{}(x,x^{})\theta \left(x_0^{}x_0\right)S_{out}^+(x,x^{}),`$ (38)
where
$`S^{}(x,x^{})`$ $`=`$ $`i{\displaystyle _{\mathrm{}}^+\mathrm{}}dp_3dp_1{\displaystyle \underset{nr}{}}{}_{}{}^{+}\psi _{p_3p_1nr}^{}(x)g\left({}_{+}{}^{}|_{}^{+}\right)_+^1\overline{\psi }_{p_3p_1nr}\left(x^{}\right),`$ (39)
$`S^+(x,x^{})`$ $`=`$ $`i{\displaystyle _{\mathrm{}}^+\mathrm{}}dp_3dp_1{\displaystyle \underset{nr}{}}{}_{}{}^{}\psi _{p_3p_1nr}^{}\left(x\right)\left[g\left({}_{}{}^{}|_{}^{\mathrm{\_}}\right)^1\right]^{}{}_{}{}^{}\overline{\psi }_{p_3p_1nr}^{}\left(x^{}\right),`$ (40)
$`S_{in}^{}(x,x^{})`$ $`=`$ $`i{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑p_3𝑑p_1{\displaystyle \underset{nr}{}}{}_{\pm }{}^{}\psi _{p_3p_1nr}^{}\left(x\right)_\pm \overline{\psi }_{p_3p_1nr}\left(x^{}\right),`$ (41)
$`S_{out}^{}(x,x^{})`$ $`=`$ $`i{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑p_3𝑑p_1{\displaystyle \underset{nr}{\overset{\pm }{}}}\psi _{p_3p_1nr}\left(x\right)^\pm \overline{\psi }_{p_3p_1nr}\left(x^{}\right).`$ (42)
$`_{nr}`$ means the summation over all discrete quantum numbers $`n,r`$ (and the integration over the continuous $`p_2`$ if $`H=0`$). Using the relations between GF and between the matrices $`G({}_{\zeta }{}^{}|{}_{}{}^{\zeta ^{}})`$ one can present the functions $`S^{},`$ $`S_{in}^{}`$ and $`S_{out}^{}`$ as follows
$`\pm S^{}(x,x^{})=S^c(x,x^{})\pm \theta ((x_0x_0^{}))S(x,x^{}),`$ (43)
$`\pm S_{in}^{}(x,x^{})=S_{in}^c(x,x^{})\pm \theta ((x_0x_0^{}))S(x,x^{}),`$ (44)
$`\pm S_{out}^{}(x,x^{})=S_{out}^c(x,x^{})\pm \theta ((x_0x_0^{}))S(x,x^{}),`$ (45)
$`S_{in}^c(x,x^{})=S^c(x,x^{})S^a(x,x^{}),S_{out}^c(x,x^{})=S^c(x,x^{})S^p(x,x^{}),`$ (46)
$`S^a(x,x^{})=i{\displaystyle _{\mathrm{}}^+\mathrm{}}dp_3dp_1{\displaystyle \underset{nr}{}}{}_{}{}^{}\psi _{p_3p_1nr}^{}(x)\left[g(_+|^{})g(_{}|^{})^1\right]^{}{}_{+}{}^{}\overline{\psi }_{p_3p_1nr}^{}(x^{}),`$ (47)
$`S^p(x,x^{})=i{\displaystyle _{\mathrm{}}^+\mathrm{}}dp_3dp_1{\displaystyle \underset{nr}{\overset{+}{}}}\psi _{p_3p_1nr}(x)\left[g(_+|^+)^1g(_+|^{})\right]{}_{}{}^{}\overline{\psi }_{p_3p_1nr}^{}(x^{}).`$ (48)
Let us consider the functions $`S^\pm `$ and $`S^{a,p}`$. The coefficients (28) do not depend on $`p_z`$, thus, one can present the functions $`S^{}`$ and $`S^{a,p}`$ in the following convenient form
$$S^{,a,p}(x,x^{})=_{\mathrm{}}^+\mathrm{}𝑑z_{\mathrm{}}^+\mathrm{}\frac{dp_3}{2\pi }e^{ip_3y^3}S_Q^{,a,p},y_\mu =x_\mu x_\mu ^{},$$
(49)
where
$`S_Q^{}`$ $`=`$ $`i{\displaystyle _{\mathrm{}}^+\mathrm{}}dp_zdp_1{\displaystyle \underset{nr}{}}{}_{}{}^{+}\psi _{p_3p_1nr}^{}(\stackrel{~}{\eta },x_{},z,p_z)g\left({}_{+}{}^{}|_{}^{+}\right)_+^1\overline{\psi }_{p_3p_1nr}(\stackrel{~}{\eta ^{}},x_{}^{},z^{},p_z),`$ (50)
$`S_Q^+`$ $`=`$ $`i{\displaystyle _{\mathrm{}}^+\mathrm{}}dp_zdp_1{\displaystyle \underset{nr}{}}{}_{}{}^{}\psi _{p_3p_1nr}^{}(\stackrel{~}{\eta },x_{},z,p_z)\left[g\left({}_{}{}^{}|_{}^{\mathrm{\_}}\right)^1\right]^{}{}_{}{}^{}\overline{\psi }_{p_3p_1nr}^{}(\stackrel{~}{\eta ^{}},x_{}^{},z^{},p_z),`$ (51)
$`S_Q^a`$ $`=`$ $`i{\displaystyle _{\mathrm{}}^+\mathrm{}}dp_zdp_1{\displaystyle \underset{nr}{}}{}_{}{}^{}\psi _{p_3p_1nr}^{}(\stackrel{~}{\eta },x_{},z,p_z)\left[g(_+|^{})g(_{}|^{})^1\right]^{}{}_{+}{}^{}\overline{\psi }_{p_3p_1nr}^{}(\stackrel{~}{\eta ^{}},x_{}^{},z^{},p_z),`$ (52)
$`S_Q^p`$ $`=`$ $`i{\displaystyle _{\mathrm{}}^+\mathrm{}}dp_zdp_1{\displaystyle \underset{nr}{\overset{+}{}}}\psi _{p_3p_1nr}(\stackrel{~}{\eta },x_{},z,p_z)\left[g(_+|^+)^1g(_+|^{})\right]{}_{}{}^{}\overline{\psi }_{p_3p_1nr}^{}(\stackrel{~}{\eta ^{}},x_{}^{},z^{},p_z),`$ (53)
and
$`{}_{\pm }{}^{}\psi _{p_3p_1nr}^{}(\stackrel{~}{\eta },x_{},z,p_z)=\left(\gamma ^0i_0+\gamma ^3\left(p_3qEx^o\right)+\gamma _{}\left(iqA\right)+M\mathrm{\Omega }\right)_\pm \varphi _{p_3p_1nr},`$ (54)
$`{}_{\pm }{}^{}\varphi _{p_3p_1nr}^{}=_\pm \varphi _{p_3nr}(\stackrel{~}{\eta },z,p_z)\varphi _{p_1nr}(x_{})v_{+1,r},_\pm \varphi _{p_3nr}(\stackrel{~}{\eta },z,p_z)={\displaystyle \frac{1}{\sqrt{2\pi }}}e_\pm ^{ip_zz}\varphi _{p_3,n,+1,r}(x^0,p_z),`$ (55)
and in the same form with $`(\pm )`$ indices above. Within the same approximation one can rewrite the functions $`S_Q^{,a,p}`$ as follows:
$$S_Q^{,a,p}=\left(\gamma ^0i_0+\gamma ^3\left(p_3qEx^o\right)+\gamma _{}\left(iqA\right)+M\mathrm{\Omega }\right)_Q^{,a,p},$$
(56)
where the functions $`\mathrm{\Delta }_Q^{,a,p}`$ obey the equation
$`\left[\left(i{\displaystyle \frac{}{\stackrel{~}{\eta }}}\right)^2\left(i{\displaystyle \frac{}{z}}\rho \stackrel{~}{\eta }\right)^2+𝒫_{}^2{\displaystyle \frac{p_3^2\left(bM^2\right)}{\rho ^2}}i\rho \mathrm{\Xi }{\displaystyle \frac{q}{2}}F_{\mu \nu }^{}\sigma ^{\mu \nu }\right]\mathrm{\Delta }_Q^{,a,p}=0.`$
They may be easily expressed via the solutions (55) according to the Eqs. (53). The functions $`\mathrm{\Delta }_Q^{,a,p}`$ are just the GF of the squared Dirac equation in electromagnetic background , where $`\stackrel{~}{\eta }`$ is the time, $`z`$ is the coordinate along the electric field, the mass $`m_Q^2=p_3^2(bM)^2/\rho ^2`$, the potential of the electromagnetic field is $`A_z=\stackrel{~}{\eta }\rho /q`$, and the spin term $`qE\gamma ^0\gamma ^3`$ is changed to $`\rho \mathrm{\Xi }`$. In what follows we are going to use the following representations :
$`\pm \mathrm{\Delta }_Q^{}=\mathrm{\Delta }_Q^c\pm \theta (y_0)\mathrm{\Delta }_Q,`$ (57)
$`\mathrm{\Delta }_Q^c={\displaystyle _{\mathrm{\Gamma }_c}}f_Q𝑑s,\mathrm{\Delta }_Q=\text{sgn}(y_0){\displaystyle _{\mathrm{\Gamma }_c\mathrm{\Gamma }_2\mathrm{\Gamma }_1}}f_Q𝑑s,`$ (58)
$`\mathrm{\Delta }_Q^a={\displaystyle _{\mathrm{\Gamma }_a}}f_Q𝑑s+\theta (z^{}z){\displaystyle _{\mathrm{\Gamma }_3+\mathrm{\Gamma }_2\mathrm{\Gamma }_a}}f_Q𝑑s,`$ (59)
$`\mathrm{\Delta }_Q^p={\displaystyle _{\mathrm{\Gamma }_a}}f_Q𝑑s+\theta (zz^{}){\displaystyle _{\mathrm{\Gamma }_3+\mathrm{\Gamma }_2\mathrm{\Gamma }_a}}f_Q𝑑s,`$ (60)
where $`\theta (0)=1/2`$, the contours of the integration are indicated on the Fig.1, and
$`f_Q=e^{s\rho \mathrm{\Xi }}\mathrm{exp}\left({\displaystyle \frac{i}{2}}qF_{\mu \nu }^{}\sigma ^{\mu \nu }s\right)f_Q^{(0)},`$ (61)
$`f_Q^{(0)}=\mathrm{exp}\left\{iq{\displaystyle _x^{}^x}A_\mu ^{}𝑑x^\mu \right\}f_Q^{}\left(zz^{}\right)f_{},`$ (62)
$`f_{}=(4\pi )^2{\displaystyle \frac{qH}{\mathrm{sin}(qHs)}}\mathrm{exp}\left\{{\displaystyle \frac{i}{4}}y_{}qF\mathrm{coth}(qFs)y_{}\right\},`$ (63)
$`f_Q^{}\left(z\right)={\displaystyle \frac{\rho }{\mathrm{sinh}(\rho s)}}\mathrm{exp}\left\{i{\displaystyle \frac{\rho }{2}}(\stackrel{~}{\eta }+\stackrel{~}{\eta ^{}})zim_Q^2s+i{\displaystyle \frac{\rho }{4}}\left[z^2(\stackrel{~}{\eta }\stackrel{~}{\eta ^{}})^2\right]\mathrm{coth}(\rho s)\right\}.`$ (64)
One can calculate in (49) all Gaussian integrals over $`p_3`$ and $`z`$. Thus, one gets
$$S^{(\mathrm{})}(x,x^{})=\left(\gamma ^\mu 𝒫_\mu +M\mathrm{\Omega }\right)^{(\mathrm{})}(x,x^{}),$$
(65)
$`\mathrm{\Delta }^c(x,x^{})={\displaystyle _{\mathrm{\Gamma }_c}}f(x,x^{},s)𝑑s,\mathrm{\Delta }(x,x^{})=\text{sgn}(y_0){\displaystyle _\mathrm{\Gamma }}f(x,x^{},s)𝑑s,`$ (66)
$`\mathrm{\Delta }^a(x,x^{})=\mathrm{\Delta }^{(1)}(x,x^{})\mathrm{\Delta }^{(2)}(x,x^{}),`$ (67)
$`\mathrm{\Delta }^p(x,x^{})=\mathrm{\Delta }^{(1)}(x,x^{})+\mathrm{\Delta }^{(2)}(x,x^{}),`$ (68)
$`\mathrm{\Delta }^{(1)}(x,x^{})={\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Gamma }_3+\mathrm{\Gamma }_2+\mathrm{\Gamma }_a}}f(x,x^{},s)𝑑s,`$ (69)
$`\mathrm{\Delta }^{(2)}(x,x^{})={\displaystyle _{\mathrm{\Gamma }_3+\mathrm{\Gamma }_2\mathrm{\Gamma }_a}}f_r(x,x^{},s)𝑑s,`$ (70)
where $`f_r(x,x^{},s)=\pi ^{1/2}/2\gamma (1/2,\alpha )f(x,x^{},s)`$, $`\gamma (1/2,\alpha )`$ is the incomplete gamma-function, and $`\alpha =e^{i\pi /2}\left(4s(bM)^2\omega \right)^1\left[\left(x_0+x_{}^{}{}_{0}{}^{}\right)s(bM)^2+qEy^3\right]^2`$. Here
$`f(x,x^{},s)=\mathrm{exp}\left(\rho \mathrm{\Xi }si{\displaystyle \frac{q}{2}}\sigma ^{\mu \nu }F_{\mu \nu }^{}s\right)f^{(0)}(x,x^{},s),`$ (71)
$`f^{(0)}(x,x^{},s)=\mathrm{exp}\left\{iq{\displaystyle _x^{}^x}A_\mu 𝑑x^\mu \right\}f_{}f_{},`$ (72)
$`f_{}={\displaystyle \frac{\rho }{\mathrm{sinh}(\rho s)\omega ^{1/2}}}\mathrm{exp}\{i{\displaystyle \frac{qE}{2}}(x_0+x_{}^{}{}_{0}{}^{})y^3i{\displaystyle \frac{\rho }{4}}(x_0x_{}^{}{}_{0}{}^{})^2\mathrm{coth}(\rho s)`$ (73)
$`i(aM)^2s+i{\displaystyle \frac{\rho }{4\omega }}y_3^2\mathrm{coth}(\rho s){\displaystyle \frac{i}{4\omega }}[(bM)^2s(x_0+x_{}^{}{}_{0}{}^{})^2+2qEy^3(x_0+x_{}^{}{}_{0}{}^{})]\},`$ (74)
where $`\omega =s\mathrm{coth}(\rho s)(bM)^2/\rho +(qE)^2/\rho ^2`$. One can see that
$`i{\displaystyle \frac{d}{ds}}f(x,x^{},s)=\left(M^2\mathrm{\Omega }^2𝒫^2+{\displaystyle \frac{q}{2}}\sigma ^{\mu \nu }F_{\mu \nu }ibM\gamma ^0\right)f(x,x^{},s),`$ (75)
$`\underset{s+0}{lim}f(x,x^{},s)=i\delta ^{(4)}(xx^{}).`$ (76)
Thus, $`f(x,x^{},s)`$ is the Fock-Schwinger function of the QED-$`\mathrm{\Omega }`$ theory, and $`s`$ is the Fock-Schwinger proper-time in this case. The contour $`\mathrm{\Gamma }_c\mathrm{\Gamma }_2\mathrm{\Gamma }_1`$ in representation (70) for $`\mathrm{\Delta }`$ transformed into $`\mathrm{\Gamma }`$ (see Fig.2) after the integration over $`p_D`$ and $`z`$. The result is consistent with the general expression for the commutation function obtained in . The function $`f^{(0)}(x,x^{},s)`$ which has appeared in (71) coincides with the Fock-Schwinger function of the scalar case. Due to that we can widely use results presented in .
If $`b0`$, then the function $`f(x,x^{},s)`$ has three singular points on the complex region between contours $`\mathrm{\Gamma }_c\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_a\mathrm{\Gamma }_3`$, which are distributed at the imaginary axis: $`\rho s_0=0,`$ $`\rho s_1=i\pi `$ and $`\rho s_2=ic_2.`$ The latter point is connected with zero value of the function $`\omega `$. We get an equation for $`c_2`$ from the condition $`\omega =0`$, $`c_2\mathrm{tan}(c_2\pi /2)\left(qE/(bM)\right)^2=0,`$ where $`\pi /2<c_2<\pi `$. The position of this point depends on the ratio $`qE/(bM)`$, e.g. $`c_2\pi `$ as $`bM/(qE)0`$ and $`c_2\pi /2`$ as $`qE/(bM)0`$. Notice, in the case $`E=0`$ it is convenient to put $`c_2=\pi /2+0`$ since the contour $`\mathrm{\Gamma }_2`$ is also passed above the singular point $`s_2`$. If $`b=0`$, then $`\omega =1`$, and the function $`f(x,x^{},s)`$ has only two singular points $`s_0`$ and $`s_1`$ on the above mentioned complex region. In this degenerate case we have known representation .
Similarly to Ref. one may demonstrate that the function $`\mathrm{\Delta }^{(2)}`$ from (70) can be presented via proper-time integrals with the Fock-Schwinger kernel $`f(x,x^{},s)`$. Then, using (75), one can verify that the representations (70) for $`\mathrm{\Delta }^{(1)}`$ and $`\mathrm{\Delta }^{(2)}`$ obey the equation (10). Thus, all the $`\mathrm{\Delta }`$-functions considered here, excluding those marked by the index “c”, are solutions of the equation (10). The important difference between the functions $`\mathrm{\Delta }^c`$, $`\mathrm{\Delta }^{(1)}`$ and $`\mathrm{\Delta }`$, $`\mathrm{\Delta }^{(2)}`$ is that the first ones are symmetric under simultaneous change of sign in $`x_0,x_0^{},x_3,x_3^{}`$ and the second ones change sign in this case.
Using the kernel (71) one can express out-in effective action in the QED-$`\mathrm{\Omega }`$ theory,
$$\mathrm{\Gamma }_{outin}=\frac{1}{2}𝑑x_0^{\mathrm{}}s^1f(x,x,s)𝑑s.$$
(77)
## IV Vacuum instability, mean values of current, and energy momentum tensor. Discussion
Using the exact solutions and the GF constructed above, one may calculate different physical effects having both local and global character. The proper-time representation of the GF gives us a possibility to study all mean values of current and energy momentum tensor in the same manner.
All the information about the processes of particles creation, annihilation, and scattering in an external field (without radiative corrections) can be extracted from the matrices $`G({}_{\zeta }{}^{}|{}_{}{}^{\zeta ^{}})`$ (13). These matrices define a canonical transformation between in and out creation and annihilation operators in the generalized Furry representation ,
$$a^{}(out)=a^{}(in)G({}_{+}{}^{}|{}_{}{}^{+})+b(in)G({}_{}{}^{}|{}_{}{}^{+}),b(out)=a^{}(in)G({}_{+}{}^{}|{}_{}{}^{})+b(in)G({}_{}{}^{}|{}_{}{}^{}).$$
(78)
Here $`a_l^{}(in)`$, $`b_l^{}(in)`$, $`a_l(in)`$, $`b_l(in)`$ are creation and annihilation operators of in-particles and antiparticles respectively and $`a_l^{}(out)`$,$`b_l^{}(out)`$, $`a_l(out),b_l(out)`$ are ones of out-particles and antiparticles, $`l`$ are possible quantum numbers (in our case it $`l=p_1,p_3,n,r`$) . For example, the mean differential number of particles created (which are also equal to the number of pairs created) by the external field from the in-vacuum $`|0,in>`$ is
$$N_l=<0,in|a_l^{}(out)a_l(out)|0,in>=|g({}_{}{}^{}|{}_{}{}^{+})|^2$$
(79)
(for a review of gravitational particles creation, see ). The standard space coordinate volume regularization was used to get the latter formula, so that $`\delta (p_jp_j^{})\delta _{p_j,p_j^{}}`$. The probability for a vacuum to remain a vacuum is
$$P_v=|c_v|^2=\mathrm{exp}\left\{\underset{l}{}\mathrm{ln}\left(1N_l\right)\right\}.$$
(80)
Similar to the electric field case we get
$$N_l==e^{\pi \lambda },\text{if}\sqrt{\rho }T1,\text{and}\sqrt{\rho }T\lambda ,\text{and}\rho ^2T|qEp_3|,$$
(81)
where $`\lambda `$ is defined in (20). The latter conditions take place with large enough time for action of the electric-like field, $`T=x_{out}^0x_{in}^0`$. The result (81) coincides with one obtained in , and for $`b=0`$ it coincides with one obtained in . Evidently the creation process is a coherent effect of both the electromagnetic and gravitational fields. If the condition $`p_3^2(bM)^2/\rho ^3<<1`$ takes place (the gravitational field is in a sense weaker than the electric one), the $`p_3`$ dependence on $`N_l`$ is similar to the case $`b=0`$. Thus , one can estimate that $`𝑑p_3=(qE)^1\rho ^2T`$ . Then the particle creation per unit of time may be calculated similar to . In strong enough gravitational background the time dependence of the effect is nonlinear and needs to be studied specially.
To get the total number $`N`$ of particles created one has to sum over the quantum numbers $`l`$. The sum over the momenta can be easily transformed into an integral. Thus, if $`b=0`$ one gets result presented in . If $`b0,`$ the total number of pairs created per space coordinate volume has the form
$$\stackrel{~}{n}^{cr}=\frac{_lN_l}{𝑑𝐱}=\frac{\beta (1)}{4\pi ^2}\frac{\rho ^{3/2}}{bM}\mathrm{exp}\left\{\pi \frac{(aM)^2}{\rho }\right\},$$
(82)
where $`\beta (n)=qH\mathrm{coth}(n\pi qH/\rho ).`$ The observable number density of the created pairs in the asymptotic region $`x_0=x_0^{out}\mathrm{}`$ is given by the expression $`n^{cr}=\stackrel{~}{n}^{cr}/\mathrm{\Omega }^3(x_0).`$ If the electromagnetic field is absent and $`a=0`$ , these results coincide with ones in . In case $`b0`$ the expression (82) is growing infinitely. In this case the particles are created in main by the electric field, whereas the parameter $`b`$ plays a role of “cut-off” factor, which eliminates creation of particles with extremely high momenta along the electric field. It is seen from the expression (81). Thus, the limit $`b0`$ corresponds to the case of the electric field which acts for infinite time. Then the number of particles created is proportional to the time of the field action. As was already remarked above, in this case $`𝑑p_3=(qE)^1\rho ^2T.`$ Then, the parameter $`b`$ may be understood as $`(\sqrt{\rho })(TM)^1`$. The vacuum-to-vacuum transition probability can be calculated, using formula (80). Thus, we get an analog of the well-known Schwinger formula in the case under consideration,
$$P_v=\mathrm{exp}\left\{\mu \stackrel{~}{n}^{cr}𝑑𝐱\right\},\mu =\underset{n=0}{\overset{\mathrm{}}{}}\frac{\beta (n+1)}{(n+1)^{3/2}\beta (1)}\mathrm{exp}\left\{n\pi \frac{(aM)^2}{\rho }\right\}.$$
(83)
Now we are going to discuss vacuum matrix elements of current and of metric energy-momentum tensor (EMT) of spinor field. Making the conformal transformation the current in FRW universe may be presented in the following form
$$J_\mu =\mathrm{\Omega }^2(x^0)j_\mu ,j_\mu =\frac{q}{2}[\overline{\psi }(x),\gamma _\mu \psi (x)],$$
(84)
where $`j_\mu `$ is QED-$`\mathrm{\Omega }`$ current of the spinor field $`\psi (x)`$. The latter obeys the Dirac equation (2). The EMT of spinor field in FRW universe may be presented as
$`\tau _{\mu \nu }=\mathrm{\Omega }^2(x^0)T_{\mu \nu },T_{\mu \nu }={\displaystyle \frac{1}{2}}\left(T_{\mu \nu }^{can}+T_{\nu \mu }^{can}\right),,`$ (85)
$`T_{\mu \nu }^{can}={\displaystyle \frac{1}{4}}\left\{[\overline{\psi }(x),\gamma _\mu 𝒫_\nu \psi (x)]+[𝒫_\nu ^{}\overline{\psi }(x),\gamma _\mu \psi (x)]\right\},`$ (86)
where $`T_{\mu \nu }`$ is EMT in QED-$`\mathrm{\Omega }.`$ In the case of unstable vacuum one needs to calculate three types of matrix elements , depending on the problem in question:
$`<`$ $`j_\mu >^c=<0,out|j_\mu |0,in>c_v^1,<T_{\mu \nu }>^c=<0,out|T_{\mu \nu }|0,in>c_v^1,`$ (87)
$`<`$ $`j_\mu >^{in}=<0,in|j_\mu |0,in>,<T_{\mu \nu }>^{in}=<0,in|T_{\mu \nu }|0,in>,`$ (88)
$`<`$ $`j_\mu >^{out}=<0,out|j_\mu |0,out>,<T_{\mu \nu }>^{out}=<0,out|T_{\mu \nu }|0,out>,`$ (89)
Using GF, which were found before, one can present these matrix elements as follows:
$`<`$ $`j_\mu >^c=iq\text{tr}\left\{\gamma _\mu S^c(x,x)\right\}=iq\text{tr}\left\{\gamma _\mu (\gamma ^\kappa 𝒫_\kappa +M\mathrm{\Omega })\mathrm{\Delta }^c(x,x^{})\right\}|_{x=x^{}},`$ (90)
$`<`$ $`T_{\mu \nu }>^c=i/4\text{tr}\left\{(\gamma _\mu (𝒫_\nu +𝒫_{}^{}{}_{\nu }{}^{})+\gamma _\nu (𝒫_\mu +𝒫_{}^{}{}_{\mu }{}^{}))S^c(x,x^{})\right\}|_{x=x^{}}`$ (91)
$`=`$ $`i\text{tr}\left\{B_{\mu \nu }\mathrm{\Delta }^c(x,x^{})\right\}|_{x=x^{}},`$ (92)
$`<`$ $`j_\mu >^{in}=<j_\mu >^c+<j_\mu >^{(1)}+<j_\mu >^{(2)},`$ (93)
$`<`$ $`j_\mu >^{out}=<j_\mu >^c+<j_\mu >^{(1)}<j_\mu >^{(2)},`$ (94)
$`<`$ $`T_{\mu \nu }>^{in}=<T_{\mu \nu }>^c+<T_{\mu \nu }>^{(1)}+<T_{\mu \nu }>^{(2)},`$ (95)
$`<`$ $`T_{\mu \nu }>^{out}=<T_{\mu \nu }>^c+<T_{\mu \nu }>^{(1)}<T_{\mu \nu }>^{(2)},`$ (96)
$`<`$ $`j_\mu >^{(1,2)}=iq\text{tr}\left\{\gamma _\mu (\gamma ^\kappa 𝒫_\kappa +M\mathrm{\Omega })\mathrm{\Delta }^{(1,2)}(x,x^{})\right\}|_{x=x^{}},`$ (97)
$`<`$ $`T_{\mu \nu }>^{(1,2)}=i\text{tr}\left\{B_{\mu \nu }\mathrm{\Delta }^{(1,2)}(x,x^{})\right\}|_{x=x^{}},`$ (98)
$`B_{\mu \nu }=1/4\left\{\gamma _\mu \left(𝒫_\nu +𝒫_{}^{}{}_{\nu }{}^{}\right)+\gamma _\nu \left(𝒫_\mu +𝒫_{}^{}{}_{\mu }{}^{}\right)\right\}\left(\gamma ^\kappa 𝒫_\kappa +M\mathrm{\Omega }\right),`$ (99)
where $`𝒫_{}^{}{}_{\mu }{}^{}=i\frac{}{x^{}^\mu }qA_\mu (x^{}),`$ the GF are given by Eqs. ((70), and the relation $`\mathrm{\Delta }^c(x,x)=(1/2)\left[\mathrm{\Delta }^{}(x,x)\mathrm{\Delta }^+(x,x)\right]`$ is used.
To get convenient forms of $`<j_\mu >^{in}`$ and $`<T_{\mu \nu }>^{in}`$ we may rewrite $`\mathrm{\Delta }^{(1)}`$ as follows:
$`\mathrm{\Delta }^{(1)}(x,x^{})={\displaystyle \frac{1}{2}}\mathrm{\Delta }^{\mathrm{\Gamma }_2}(x,x^{})+\mathrm{\Delta }^{(3)}(x,x^{}),\mathrm{\Delta }^{\mathrm{\Gamma }_2}(x,x^{})={\displaystyle _{\mathrm{\Gamma }_2}}f(x,x^{},s)𝑑s,`$ (100)
$`\mathrm{\Delta }^{(3)}(x,x^{})={\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Gamma }_3+\mathrm{\Gamma }_a}}f(x,x^{},s)𝑑s.`$ (101)
The contour $`\mathrm{\Gamma }_2`$ does not pass through the singular points of the function $`f(x,x^{},s)`$, thus the contributions from $`\mathrm{\Delta }^{\mathrm{\Gamma }_2}`$ to EMT and to the current are finite. Displacing the contour $`\mathrm{\Gamma }_2`$ to the real axis we may present the difference $`\mathrm{\Delta }^c\mathrm{\Delta }^{\mathrm{\Gamma }_2}`$ as
$`\mathrm{\Delta }^c(x,x^{})\mathrm{\Delta }^{\mathrm{\Gamma }_2}(x,x^{})=\mathrm{\Delta }^{\overline{c}}(x,x^{})+\mathrm{\Delta }^\mathrm{\Gamma }(x,x^{}),`$ (102)
$`\mathrm{\Delta }^{\overline{c}}(x,x^{})={\displaystyle _0^{\mathrm{}}}f(x,x^{},s)𝑑s,\mathrm{\Delta }^\mathrm{\Gamma }(x,x^{})={\displaystyle _\mathrm{\Gamma }}f(x,x^{},s)𝑑s,`$ (103)
where $`\mathrm{\Delta }^{\overline{c}}`$ is related to the anticausal GF, $`S^{\overline{c}}(x,x^{})=(\gamma 𝒫+M\mathrm{\Omega })\mathrm{\Delta }^{\overline{c}}(x,x^{})`$, and the integral $`\mathrm{\Delta }^\mathrm{\Gamma }`$ can be expressed by means of the $`\mathrm{\Delta }`$-function related to the commutation function, $`S(x,x^{})=(\gamma 𝒫+M\mathrm{\Omega })\mathrm{\Delta }(x,x^{})`$ : $`\mathrm{\Delta }^\mathrm{\Gamma }(x,x^{})=\text{sgn}(x_0x_0^{})\mathrm{\Delta }(x,x^{}).`$ Taking into account the above definition the space-like part of the limit $`xx^{}0`$ in the EMT expressions have to be treated either as the limit $`x_0x_0^{}=+0`$ or as $`x_0x_0^{}=0`$, respectively. Then, in according to the initial condition for the commutation function we get $`\mathrm{\Delta }(x,x^{})|_{x_0=x_0^{}}=0,_0\mathrm{\Delta }(x,x^{})|_{x_0=x_0^{}}=\delta (𝐱𝐱^{}).`$ Thus, the contributions from $`\mathrm{\Delta }`$ into $`<T_{\mu \nu }>^{in}`$ and $`<j_\mu >^{in}`$ are zero. (All contributions from $`\mathrm{\Delta }`$, which might appear as a result of a change of the limit definition, are background independent. They may be eliminated by a renormalization.) It follows from (71) that
$$f(x,x^{},s^{})=\gamma ^0f(x^{},x,s)^{}\gamma ^0.$$
(104)
Changing $`ss`$ one can represent integral $`\mathrm{\Delta }^{\overline{c}}`$ from (103) in the form
$$\mathrm{\Delta }^{\overline{c}}(x,x^{})=_0^{\mathrm{}}\gamma ^0f(x^{},x,s)^{}\gamma ^0𝑑s,$$
(105)
to get
$`<`$ $`j_\mu >^{in}=Re<j_\mu >^c+<j_\mu >^{(2)}+<j_\mu >^{(3)},`$ (106)
$`<`$ $`T_{\mu \nu }>^{in}=Re<T_{\mu \nu }>^c+<T_{\mu \nu }>^{(2)}+<T_{\mu \nu }>^{(3)},`$ (107)
where the terms $`<j_\mu >^{(3)}`$ and $`<T_{\mu \nu }>^{(3)}`$ represent the contributions from $`\mathrm{\Delta }^{(3)}`$:
$`<`$ $`j_\mu >^{(3)}=iq\text{tr}\left\{\gamma _\mu (\gamma ^\kappa 𝒫_\kappa +M\mathrm{\Omega })\mathrm{\Delta }^{(3)}(x,x^{})\right\}|_{x=x^{}},`$ (108)
$`<`$ $`T_{\mu \nu }>^{(3)}=i\text{tr}\left\{B_{\mu \nu }\mathrm{\Delta }^{(3)}(x,x^{})\right\}|_{x=x^{}}.`$ (109)
Using the Eq. (104) one can verify that each of the terms $`<j_\mu >^{(2)},`$ $`<j_\mu >^{(3)},`$ $`<T_{\mu \nu }>^{(2)}`$ and $`<T_{\mu \nu }>^{(3)}`$ is real, as it follows from their initial definitions.
One can get similar expressions for the out-out-matrix elements of the current and EMT:
$`<`$ $`j_\mu >^{out}=Re<j_\mu >^c<j_\mu >^{(2)}+<j_\mu >^{(3)},`$ (110)
$`<`$ $`T_{\mu \nu }>^{out}=Re<T_{\mu \nu }>^c<T_{\mu \nu }>^{(2)}+<T_{\mu \nu }>^{(3)}.`$ (111)
All nondiagonal matrix elements of EMT are zero:
$$<T_{\mu \nu }>^c=<T_{\mu \nu }>^{(2)}=<T_{\mu \nu }>^{(3)}=0,\mu \nu .$$
(112)
Using the formulas
$`\mathrm{exp}\left(\rho \mathrm{\Xi }s\right)`$ $`=`$ $`\mathrm{cosh}\left(\rho s\right)+\mathrm{\Xi }\mathrm{sinh}\left(\rho s\right),`$
$`\mathrm{exp}\left(i{\displaystyle \frac{q}{2}}\sigma ^{\mu \nu }F_{\mu \nu }^{}s\right)`$ $`=`$ $`\mathrm{cos}\left(qHs\right)+\gamma ^2\gamma ^1\mathrm{sin}\left(qHs\right),`$
one may calculate the traces for the current and EMT:
$`\tau (s)=\text{tr}Y=4\mathrm{cosh}\left(\rho s\right)\mathrm{cos}\left(qHs\right),Y=\left\{\mathrm{exp}\left(\rho \mathrm{\Xi }si{\displaystyle \frac{q}{2}}\sigma ^{\mu \nu }F_{\mu \nu }^{}s\right)\right\},`$
$`\text{tr}\left\{\mathrm{\Xi }Y\right\}=\mathrm{tanh}(\rho s)\tau (s),\text{tr}\left\{\gamma ^0\gamma ^3Y\right\}={\displaystyle \frac{qE}{\rho }}\mathrm{tanh}(\rho s)\tau (s),`$
$`\text{tr}\left\{\gamma ^0Y\right\}={\displaystyle \frac{bM}{\rho }}\mathrm{tanh}(\rho s)\tau (s),\text{tr}\left\{\gamma ^2\gamma ^1Y\right\}=\mathrm{tan}(qHs)\tau (s),`$
and other traces are zero. Then, the current matrix elements can be found in the forms
$`<`$ $`j_\mu >^c={\displaystyle _{\mathrm{\Gamma }_c}}\alpha _\mu (s)\tau (s)f^{(0)}(x,x^{},s)ds|_{x=x^{}},`$ (113)
$`<`$ $`j_\mu >^{(3)}={\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Gamma }_3+\mathrm{\Gamma }_a}}\alpha _\mu (s)\tau (s)f^{(0)}(x,x^{},s)ds|_{x=x^{}},`$ (114)
$`<`$ $`j_\mu >^{(2)}={\displaystyle _{\mathrm{\Gamma }_3+\mathrm{\Gamma }_2\mathrm{\Gamma }_a}}\alpha _\mu (s)\tau (s)f_r^{(0)}(x,x^{},s)ds|_{x=x^{}},`$ (116)
$`\alpha _\mu (s)=iq\delta _\mu ^3\left[𝒫_3+𝒫_0{\displaystyle \frac{qE}{\rho }}\mathrm{tanh}(\rho s)\right],`$
where the only $`x^3`$ components (along the electric field) of the currents differ from zero and vanish in the absence of the electric field, as it presented to be. The nonzero components of $`<T_{\mu \nu }>^c`$ are
$`<`$ $`T_{\mu \nu }>^c={\displaystyle _{\mathrm{\Gamma }_c}}t_\mu (s)\tau (s)f^{(0)}(x,x,s)ds,\text{if}\mu =\nu ,`$ (117)
$`t_0(s)={\displaystyle \frac{\rho }{\mathrm{sinh}(2\rho s)}},t_2(s)=t_1(s)={\displaystyle \frac{qH}{\mathrm{sin}(2qHs)}},`$ (119)
$`t_3(s)={\displaystyle \frac{\rho }{2\omega }}\mathrm{coth}(\rho s){\displaystyle \frac{(qE)^2}{2\rho }}\mathrm{tanh}(\rho s)+i\left({\displaystyle \frac{qE}{2}}(x_0+x_0^{})(1\omega ^1)\right)^2,`$
and nonzero components of $`<T_{\mu \nu }>^{(2)}`$ and $`<T_{\mu \nu }>^{(3)}`$ can be written as follows:
$`<`$ $`T_{\mu \nu }>^{(3)}={\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Gamma }_3+\mathrm{\Gamma }_a}}l_\mu (s)\tau (s)f^{(0)}(x,x,s)ds,\text{if}\mu =\nu ,`$ (120)
$`<`$ $`T_{\mu \nu }>^{(2)}={\displaystyle _{\mathrm{\Gamma }_3+\mathrm{\Gamma }_2\mathrm{\Gamma }_a}}l_\mu (s)\tau (s)f_r^{(0)}(x,x^{},s)ds|_{x=x^{}},\text{if}\mu =\nu ,`$ (123)
$`l_1(s)=t_1(s),l_2(s)=t_2(s),l_3(s)=i\left(𝒫_3\right)^2{\displaystyle \frac{(qE)^2}{2\rho }}\mathrm{tanh}(\rho s),`$
$`l_0(s)=l_1(s)+l_2(s)+l_3(s)+i\left[2\left(\omega ^1b^2M^2sx_0\right)^2i\omega ^1b^2M^2s+M^2\mathrm{\Omega }^2\right]{\displaystyle \frac{\rho }{2}}\mathrm{tanh}(\rho s).`$
Here, $`l_0(s)`$ is obtained from $`t_0(s)`$, using Eq. (75) and taking into account that the $`\mathrm{\Delta }^{(2)}`$ function obeys the equation (10).
The components $`<j_\mu >^c`$ and $`<T_{\mu \nu }>^c`$ are expressed by the causal GF in the well-known Schwinger’s form. They are sources of information about local features of the theory. For the first time it is obtained exactly with respect to the given external background. The components $`<j_\mu >^{(2,3)}`$ and $`<T_{\mu \nu }>^{(2,3)}`$ are only related to global features of the theory and characterize the vacuum instability. Such expressions cannot be calculated in the frame of the perturbation theory with respect to the external background or in the frame of the WKB method. Then, such terms has never been studied. If the parameter $`b0`$, the only expression (117) for $`<T_{\mu \nu }>^c`$ has to be regularized and then renormalized. The expression (113) for $`<j_\mu >^c`$ is finite after the regularization lifting. The terms $`<j_\mu >^{(2,3)}`$ and $`<T_{\mu \nu }>^{(2,3)}`$ are also finite. Thus, as one can see from Eqs. (107) and (111), the divergence in $`<T_{\mu \nu }>^{in}`$ and $`<T_{\mu \nu }>^{out}`$ is the same as in $`<T_{\mu \nu }>^c.`$ That is consistent with the fact that the ultraviolet divergences have a local nature and result (as in the theory without an external field) from the leading local terms at $`s+0`$. Nevertheless, as $`b0`$ ( flat space limit) the terms $`<j_\mu >^{(2,3)}`$ and $`<T_{\mu \nu }>^{(2,3)}`$ have electric field dependent divergences defined by the dimensionless parameter $`\beta =bM/qE.`$ The nature of that is similar to the divergence of $`\stackrel{~}{n}^{cr}`$ for $`b0`$. In the case $`b=0`$ the density of excited vacuum states is proportional to the time $`T`$. After a special regularization with respect to time $`T`$ all of the terms become finite.
Now we are going to discuss current, energy density and pressure of the created particles. Let us introduce normalized values of current and EMT, which may be easily connected to observable values in some appropriate asymptotic region $`x_0=x_0^{as}`$,
$$j_\mu ^{cr}=\stackrel{~}{j}_\mu ^{cr}/\mathrm{\Omega }^3(x_0),T_{\mu \nu }^{cr}=\stackrel{~}{T}_{\mu \nu }^{cr}/\mathrm{\Omega }^3(x_0),$$
(124)
where according to the definitions (88) and (89) the corresponding densities of particles created per space-coordinates volume are
$`\stackrel{~}{j}_\mu ^{cr}={\displaystyle \frac{𝑑𝐱\left(<j_\mu >^{in}<j_\mu >^{out}\right)}{𝑑𝐱}},x_0=x_0^{as},`$ (125)
$`\stackrel{~}{T}_{\mu \nu }^{cr}={\displaystyle \frac{𝑑𝐱\left(<T_{\mu \nu }>^{in}<T_{\mu \nu }>^{out}\right)}{𝑑𝐱}},x_0=x_0^{as}.`$ (126)
Then, using representations (106), (107) and (110), (111) one gets from (125) and (126),
$$\stackrel{~}{j}_\mu ^{cr}=2<j_\mu >^{(2)},\stackrel{~}{T}_{\mu \nu }^{cr}=2<T_{\mu \nu }>^{(2)},x_0=x_0^{as}.$$
(127)
The component $`T_0^{cr0}=\mathrm{\Omega }^4\stackrel{~}{T}_{00}^{cr}`$ is the energy density and $`T_\mu ^{cr\nu }=\mathrm{\Omega }^4\stackrel{~}{T}_{\mu \nu }^{cr}`$ for $`\mu =\nu =1,2,3`$ are components of pressure which are measured by the cosmic observer relative to the measured volume. In contrast to that, measured energy and pressure taken per coordinate volume are given by $`\stackrel{ˇ}{T}_0^{cr0}=\mathrm{\Omega }^1\stackrel{~}{T}_{00}^{cr},\stackrel{ˇ}{T}_\mu ^{cr\nu }=\mathrm{\Omega }^1\stackrel{~}{T}_{\mu \nu }^{cr},\mu =\nu =1,2,3.`$
We are going to analyze contributions to the quantities (127) at $`x_0=x_0^{as}`$. The leading asymptotics in $`<j_\mu >^{(2)}`$ and $`<T_{\mu \nu }>^{(2)}`$ at $`x_0>>\sqrt{\rho }/(bM)`$ are determinated by the expressions
$$j_\mu >^{(2)}=\alpha _\mu (s_1)\text{tr}\left\{\mathrm{\Delta }^{(2)}(x,x^{})\right\}|_{x=x^{}},T_{\mu \mu }>^{(2)}=l_\mu (s_1)\text{tr}\left\{\mathrm{\Delta }^{(2)}(x,x^{})\right\}|_{x=x^{}}.$$
(128)
An asymptotic expression for $`\text{tr}\mathrm{\Delta }^{(2)}`$ with $`xx^{}`$ can be found using the method described in the App. A of Ref. for the scalar case. Then, one gets
$`\text{tr}\mathrm{\Delta }^{(2)}={\displaystyle \frac{i\stackrel{~}{n}^{cr}}{\rho (x_0+x_0^{})}}\mathrm{exp}\left\{iq\mathrm{\Lambda }+{\displaystyle \frac{iqE}{2}}(x_0+x_0^{})y^3{\displaystyle \frac{\rho ^3(y_3)^2}{4\pi (bM)^2}}+y_{}{\displaystyle \frac{qF}{4}}\mathrm{cot}\left(\pi qF/\rho \right)y_{}\right\},`$
where $`\stackrel{~}{n}^{cr}`$ is defined in (82). This expression is also valid when $`x_0<0,`$ since the function $`\mathrm{\Delta }^{(2)}`$ is odd in $`x_0`$ as $`x=x^{}`$ (see Sec.III ). Thus, one can see when $`x_0>>\sqrt{\rho }/(bM)`$ the leading terms in $`<j_\mu >^{(2)}`$ and $`<T_{\mu \nu }>^{(2)}`$ are
$`<`$ $`j_\mu >^{(2)}=q^2E\rho ^1\stackrel{~}{n}^{cr}\delta _\mu ^3;`$ (129)
$`<`$ $`T_{00}>^{(2)}=[\rho ^2x_0^2+a^2M^2+2qH\mathrm{sinh}^1(2\pi qH/\rho )]{\displaystyle \frac{\stackrel{~}{n}^{cr}}{\rho x_0}},`$ (130)
$`<`$ $`T_{11}>^{(2)}=<T_{22}>^{(2)}=qH\mathrm{sinh}^1(2\pi qH/\rho ){\displaystyle \frac{\stackrel{~}{n}^{cr}}{\rho x_0}},`$ (131)
$`<`$ $`T_{33}>^{(2)}=[(qEx_0)^2+{\displaystyle \frac{\rho ^3}{2\pi (bM)^2}}]{\displaystyle \frac{\stackrel{~}{n}^{cr}}{\rho x_0}}.`$ (132)
Comparing our results with the scalar theory case , when $`x_0\mathrm{}`$ we may see the difference only in the factor $`\stackrel{~}{n}^{cr}`$ which enters in the leading terms of $`<j_\mu >^{(2)},`$ $`<T_{00}>^{(2)}`$ and $`<T_{33}>^{(2)}.`$
Doubling the expressions (132) according to the Eqs. (127), one gets the mean densities for current and EMT of particles created. It turns out that these quantities are proportional to the density of total number of particles and antiparticles created ($`2\stackrel{~}{n}^{cr}`$) for the infinite time and do not change their structure with increasing of $`x_0`$. The latter means that one can consider all the expressions obtained at any fixed $`x_0`$ as asymptotic forms if $`x_0\sqrt{\rho }/bM`$. In a strong background $`a^2M^2/\rho 1`$ and, therefore, such a time has to obey the condition $`x_0>>a/b`$. Thus, in the strong background our asymptotic conformal time $`x_0`$, which is large enough in quantum sense explained, corresponds to the large cosmological time $`t`$.
Note, that one can neglect the second term in the brackets of the expression (132) for $`<T_{33}>^{(2)}`$ at $`bM/(qE)1`$. Also one can neglect both the term $`a^2M^2`$ in the brackets of the expression (132) for $`<T_{00}>^{(2)}`$ in case of strong background $`a^2M^2/\rho 1`$ and third term in the same expression if the magnetic field is not strong enough, $`qH/\rho 1.`$ The current density $`\stackrel{~}{j}_\mu ^{cr}=2<j_\mu >^{(2)}`$ does not depend on the asymptotic time. As $`b0`$ this expression coincides with the one for flat space, $`<\stackrel{~}{j}_\mu >^{cr}=2q\stackrel{~}{n}^{cr}\delta _\mu ^3.`$ The pressure component along the electric field direction $`\stackrel{~}{T}_{33}^{cr}=2<T_{33}>^{(2)}`$ is growing with time upon the action of the field. However, if $`qE/(bM)<<1`$ then the asymptotic condition for $`x_0`$ is consistent with the fact that the term $`(qEx_0)^2`$ in the expressions $`<T_{33}>^{(2)}`$ from (132) will not be dominant before large enough time instant $`x_0`$. In this case one can neglect the contribution which depends on the electric field, if the field is switched off at a fixed time.
The components of the pressure in the directions, which are perpendicular to the electric field, $`\stackrel{~}{T}_{11}^{cr}=2<T_{11}>^{(2)}`$ and $`\stackrel{~}{T}_{22}^{cr}=2<T_{22}>^{(2)}`$, decrease when the magnetic field increases. If an electromagnetic field is absent, the pressure is isotropic according to the symmetry of the space-time and to the corresponding symmetry of the vacuum. For $`x^0\mathrm{}`$ the remaining terms of the measured energy and pressure of the created particles taken per coordinate volume are
$`\stackrel{ˇ}{T}_0^{cr0}={\displaystyle \frac{\rho }{b}}2\stackrel{~}{n}^{cr},\stackrel{ˇ}{T}_3^{cr3}={\displaystyle \frac{(qE)^2}{\rho b}}2\stackrel{~}{n}^{cr},`$
and, if an electric field is absent the only remaining term is $`\stackrel{ˇ}{T}_0^{cr0}=2M\stackrel{~}{n}^{cr}.`$ The last represents the total rest mass per coordinate volume and coincides with the result of Ref..
Another problem is to study a back-reaction of particles created on the electromagnetic field and metrics, or to be more correct, it is better to say about back-reaction effects produced by both of particles created from a vacuum and polarization of an unstable vacuum. To this end one needs to use the expressions $`<j_\mu >^{in}`$ and $`<T_{\mu \nu }>^{in}`$ for all times $`x_0`$. The explicit proper-time expressions found above give us promising tool for accurate analysis of the total back-reaction related to the vacuum instability. In this approach we do not need to select phenomenologically parts come from real particles and from a vacuum polarization, that has been usually done in literature (see, for example,. review in ). Of course, to find self-consistent solution taking into account the back-reaction one needs numerical estimations (compare with pure electromagnetic case, ). Keeping in mind such an application, one needs to get preliminary information about the behavior of the expressions (106) and (107) in time and to select the leading components. To this end let us estimate the above mentioned expressions at characteristic large time, $`x_0^2>>\rho /(bM)^2,`$ and at characteristic small time, $`x_0^2<<\rho /(bM)^2.`$
Neglecting divergent terms in $`<T_{\mu \nu }>^c`$ according to the standard renormalization procedure, one can reduce Eq. (107) to the following finite form
$$<T_{\mu \nu }>_{fin}^{in}=Re<T_{\mu \nu }>_{fin}^c+<T_{\mu \nu }>^{(2)}+<T_{\mu \nu }>^{(3)}.$$
(133)
Here $`<T_{\mu \nu }>_{fin}^c`$ is a finite part of $`<T_{\mu \nu }>^c`$. An estimation of $`<T_{\mu \nu }>_{fin}^c`$ can be made only after renormalization. Since we are here interested in to reveal global features of the theory, details of the renormalization problem together with others, related to the renormalization, will be considered in the next paper. To select contributions related to the vacuum instability, we note that the functions $`\mathrm{\Delta }^c`$ (70) and $`\mathrm{\Delta }^{(3)}`$ (101) with $`x=x^{}`$ are even in $`x_0`$. Thus, the functions $`<T_{\mu \nu }>^c`$ and $`<T_{\mu \nu }>^{(3)}`$ are also even and do not vanish when $`x_00`$, whereas the functions $`<j_\mu >^c`$ and $`<j_\mu >^{(3)}`$ are odd and vanish in the limit. Moreover, $`<j_\mu >^c=<j_\mu >^{(3)}=0`$ for all $`x_0`$ with $`b=0`$. The proper-time integral $`\mathrm{\Delta }^{(2)}`$ (70) is odd in $`x_0`$ with $`x=x^{}`$ and vanishes when $`x_00`$. Thus, the expression $`<T_{\mu \nu }>^{(2)}`$ is also odd in $`x_0`$ and vanishes in this limit. The term $`<j_\mu >^{(2)}`$ is not zero if $`E0`$. It is even in $`x_0`$ and differs from zero when $`x_00`$.
The asymptotic behavior of $`<j_\mu >^{(3)}`$ and $`<T_{\mu \nu }>^{(3)}`$ is defined by the asymptotic expression for $`\mathrm{\Delta }^{(3)}.`$ Using the method presented in Appendix A of Ref. one can verify that asymptotically
$$\mathrm{\Delta }^{(3)}(x,x^{})=\text{sign}(x_0)\mathrm{\Delta }^{(2)}(x,x^{}),$$
(134)
and therefore,
$$<j_\mu >^{(3)}=\text{sign}(x_0)<j_\mu >^{(2)},<T_{\mu \nu }>^{(3)}=\text{sign}(x_0)<T_{\mu \nu }>^{(2)}.$$
(135)
The expression (113) does not need to be renormalized, thus, one can easily verify that the relation $`<j_3>^cx_0^10`$ holds asymptotically. Then, the asymptotics of $`<j_\mu >^{in}`$ is determined by the asymptotic behavior of $`<j_\mu >^{(2)}`$ and $`<j_\mu >^{(3)}.`$ If $`x_0`$ with the large modulus is negative, $`<j_\mu >^{in}=Re<j_\mu >^c`$ and $`<T_{\mu \nu }>_{fin}^{in}=Re<T_{\mu \nu }>_{fin}^c`$. Thus, in this regime the vacuum instability in the time-dependent background does not affect the expressions. If $`\beta =bM/(qE)<<1,`$ the expression $`Re<T_{\mu \nu }>_{fin}^c`$ remains finite with $`b0`$, while the terms $`<T_{\mu \nu }>^{(2)}`$ and $`<T_{\mu \nu }>^{(3)}`$ diverge in the limit since the global effect: the total density of excited vacuum states in an electric field grows infinitely and its mean momentum along the electric field increases. In this case the asymptotics of $`<T_{\mu \nu }>_{fin}^{in}`$ is determined by asymptotic behavior of $`<T_{\mu \nu }>^{(2)}`$ and $`<T_{\mu \nu }>^{(3)},`$ i.e., by the asymptotic global properties of the theory.
Let us estimate $`<j_\mu >^{in}`$ and $`<T_{\mu \nu }>_{fin}^{in}`$ for small $`x_0`$. The terms $`<j_\mu >^c`$ and $`<j_\mu >^{(3)}`$ are odd in $`x_0`$. They vanish at $`x_00.`$ That is why the leading term in (106) is $`<j_\mu >^{(2)}.`$ The leading contribution from the latter can be represent as
$`<j_\mu >^{(2)}=\alpha _\mu (s_2)\text{tr}\left\{\mathrm{\Delta }^{(2)}(x,x^{})\right\}|_{x=x^{}}.`$
Then, using the expressions (A2) from the Appendix A, one gets
$$<j_\mu >^{(2)}=qn^{(2)}\left[1+\left(c_2(bM)^2/(\rho qE)\right)^2\right]\delta _\mu ^3.$$
(136)
The expression (136) differs essentially from the asymptotic one (132). That remains true in the quasi-flat metrics ( $`\beta =bM/(qE)<<1`$), when the vacuum instability results only from the electric field. In the case of the small time limit the expression $`<j_\mu >^{(2)}=q\stackrel{~}{n}^{cr}\delta _\mu ^3`$ differs by a sign from the asymptotics (132). The factor $`\stackrel{~}{n}^{cr}`$ is known from the asymptotics. That is natural since value of the time instant $`x_0`$ is large enough ($`x_0+T/2>>(qE)^{1/2}`$ and $`x_0+T/2>>(aM)^2(qE)^{3/2}`$ according to (81) ) with respect to the initial time instant $`(x_0^{in}=T/2)`$. Then the field action time is large enough to obey the stabilization condition and the density of exited states has already the asymptotic form. However, such $`x_0`$ is still less than the asymptotic time for any $`\beta `$. Hence, at the time instant $`x_0`$ the vacuum differs essentially from $`|0,out>`$ vacuum. Technically that means that at the small time, $`x_0<<\sqrt{qE}/(bM)`$, one cannot obtain the normal form of operators by means of the $`\mathrm{\Delta }^{out}`$-function and, consequently, to use only the term $`<j_\mu >^{(2)}`$ for calculation of the mean current of quasiparticles, which are created up to that time. The small time expressions one can use only to calculate the back-reaction. Thus, in the external background under consideration a small time back-reaction to the mean electromagnetic field is not a screening (as it may be expected by analogy with a back-reaction of real particles created) but an induction. Since the small time considered is related to early stages of the Universe, such an observation may be important.
Since the function $`<T_{\mu \nu }>^{(2)}`$ is odd in $`x_0`$ and vanishes as $`x_00,`$ the leading contribution to (133) is determined by $`Re<T_{\mu \nu }>_{fin}^c`$ and $`<T_{\mu \nu }>^{(3)}.`$ However, if $`\beta =bM/(qE)<<1`$, the situation is more simple in the domain where $`\beta ^2>>qEx_0^2`$, and the leading contributions come from $`<T_{\mu \nu }>^{(3)}`$ only. That is related to the same global effect, mentioned above in the asymptotic case, since the expression $`Re<T_{\mu \nu }>_{fin}^c`$ remains finite with $`b0`$, while the term $`<T_{\mu \nu }>^{(3)}`$ diverges in the limit. The total density of excited vacuum states grows infinitely and the mean momentum of the states along the electric field increases, as well. Taking into account that the leading contributions in $`<T_{\mu \nu }>^{(3)}`$ result from the integral over a neighborhood of the singular point $`s_1`$, one gets
$`<`$ $`T_{00}>^{(3)}=<T_{33}>^{(3)}={\displaystyle \frac{\sqrt{qE}}{\pi \beta }}\stackrel{~}{n}^{cr},`$ (137)
$`<`$ $`T_{11}>^{(3)}=<T_{22}>^{(3)}=qH/\sqrt{qE}\mathrm{sinh}^1(2\pi H/E)\stackrel{~}{n}^{cr}\beta \mathrm{ln}\beta ^1.`$ (138)
We see that leading vacuum polarization effects are in main defined by the total density of excited vacuum states, which depend on the complete history of a state. Thus, it is a nonlocal contribution.
Considering the small time limit of $`<T_{\mu \nu }>^{(2)}`$ we get
$`<T_{\mu \mu }>^{(2)}=l_\mu (s_2)\text{tr}\left\{\mathrm{\Delta }^{(2)}(x,x^{})\right\}|_{x=x^{}},`$
where $`\text{tr}\mathrm{\Delta }^{(2)}`$ was found in (A2). If $`\beta <<1`$, the result has a simple form:
$`<`$ $`T_{00}>^{(2)}=<T_{33}>^{(2)}=qEx_0\stackrel{~}{n}^{cr},`$
$`<`$ $`T_{11}>^{(2)}=<T_{22}>^{(2)}=2qH\mathrm{sinh}^1(2\pi H/E)x_0\pi \beta ^2\stackrel{~}{n}^{cr}.`$
One can see that expressions for $`<T_{\mu \nu }>^{(2)}`$ at the asymptotic time and at the small time are quite different. The only components $`<T_{00}>^{(2)}=<T_{33}>^{(2)}`$ for $`\beta <<1`$ coincide.
The technics developed in the article allows one take into account global effects of both a real particles creation and a polarization on the same footing. The explicit form and the limit estimations for the mean vacuum values of EMT may be used for the calculations of the back reaction from the unstable vacuum to the gravitational background (in particular computing simulation similar to , see also and references therein). We see a behaviour of such components in time are quite different, and one needs to take into account characteristic polarization effect. The proper-time representations of the GF may be the necessary step in the study of chiral symmetry breaking in QED and the four-fermion models under the action of gravitational and electromagnetic fields . Such a study may have an immediate important application to EU, for example, through the construction of inflationary Universe where role of inflaton is played by the condensate $`<\overline{\mathrm{\Psi }}\mathrm{\Psi }>`$. One can also analyse symmetry breaking phenomenon under the combined action of these fields in the Standard Model (using also its gauged NJL form ), or GUT in the same way as it has been done in curved spacetime (without electromagnetic field) . Our methods maybe applied to the study of SUSY NJL model in curved spacetime introduced in . The chiral symmetry breaking under the action of weak gravitational field and constant electromagnetic field has been studied in . Having our results we may study the chiral symmetry breaking for such model in FRW Universe with constant electromagnetic field.
## V Acknowledgments
D.M.G. thanks Brazilian foundations CNPq and FAPESP for support. S.P.G. thanks Brazilian foundation CAPES for support, and the Department of Physics of UEL and Department of Mathematical Physics of USP for hospitality. We also thank Prof. S.D. Odintsov for useful discussions.
## A Small time expansion of $`\mathrm{\Delta }^{(2)}`$-function
To calculate small time expansion of $`\text{tr}\left\{\mathrm{\Delta }^{(2)}\right\}`$ given by Eq. (70), in the case $`x_0^2<<\rho /(bM)^2`$ and $`xx^{}`$, one can use the method presented in App. B of Ref. . The final result is
$`\text{tr}\left\{\mathrm{\Delta }^{(2)}(x,x^{})\right\}=e^{iq\mathrm{\Lambda }}\{[i(x_0+x_0^{})c_2(bM)^2/\rho qEy^3]n^{(2)}/(qE)\phi _0`$ (A1)
$`+i[(x_0+x_0^{})^3{\displaystyle \frac{(bM)^4}{12qE\rho ^2}}K(2)+(1/2)(x_0+x_0^{})(y_3)^2qEK(0)]\},`$ (A2)
$`\phi _0=\mathrm{exp}\left(i{\displaystyle \frac{qE}{2}}(x_0+x_0^{})y^3{\displaystyle \frac{\rho }{4c_2}}\left({\displaystyle \frac{qE}{bM}}\right)^2(x_0x_0^{})^2{\displaystyle \frac{\rho ^3(y_3)^2}{4c_2(bM)^2}}\right)`$ (A3)
$`\times \mathrm{exp}\left({\displaystyle \frac{1}{4}}y_{}qF\mathrm{cot}(\pi qF/\rho )y_{}\right),`$ (A4)
$`K(l)={\displaystyle \frac{c_2^{l+1}\left(\rho /(bM)\right)^2}{c_2^2+\left(\frac{qE}{bM}\right)^2+\left(\frac{qE}{bM}\right)^4}}\{a^2M^2/\rho c_2^1(l1/2)c_2^1\left({\displaystyle \frac{qE}{bM}}\right)^2`$ (A5)
$`c_2\left({\displaystyle \frac{bM}{qE}}\right)^2+2qH/\rho \mathrm{sinh}^1(2c_2qH/\rho )+2c_2^1{\displaystyle \frac{\left[1+\left(\frac{qE}{bM}\right)^2\right]\left[c_2^2+\left(\frac{qE}{bM}\right)^4\right]}{c_2^2+\left(\frac{qE}{bM}\right)^2+\left(\frac{qE}{bM}\right)^4}}\}n^{(2)},`$ (A6)
$`n^{(2)}=1/2\tau (s_2)n_{sc}^{(2)},`$ (A7)
$`n_{sc}^{(2)}={\displaystyle \frac{\sqrt{c_2^2+\left(\frac{qE}{bM}\right)^4}}{8\pi ^{3/2}\sqrt{c_2}\left[c_2^2+\left(\frac{qE}{bM}\right)^2+\left(\frac{qE}{bM}\right)^4\right]}}{\displaystyle \frac{q^2HE\rho ^{5/2}}{(bM)^3\mathrm{sin}(c_2qH/\rho )}}e^{c_2a^2M^2/\rho },`$ (A8)
If $`bM/qE<<1`$, the coefficients $`K(l)`$ and $`n^{(2)}`$ from (A2) have more simple form. Thus, in the case of an intensive electric field ($`a^2M^2/(qE)<1,|H/E|<1`$) we get
$$n^{(2)}=\stackrel{~}{n}^{cr},K(l)=\pi ^l\stackrel{~}{n}^{cr}.$$
(A9)
Since the final formula (B4) in App. B of Ref., which describes small time expansion of $`\mathrm{\Delta }_{sc}^{(2)}`$ in scalar case, was written with some misprints, we would like to represent its corrected form here:
$`\mathrm{\Delta }_{sc}^{(2)}(x,x^{})=e^{iq\mathrm{\Lambda }}\{[i(x_0+x_0^{})c_2(bM)^2/\rho qEy^3](1/2)n_{sc}^{(2)}/(qE)\phi _0`$
$`+i[(x_0+x_0^{})^3{\displaystyle \frac{(bM)^4}{12qE\rho ^2}}K_{sc}(2)+(1/2)(x_0+x_0^{})(y_3)^2qEK_{sc}(0)]\},`$
$`K_{sc}(l)={\displaystyle \frac{c_2^{l+1}\left(\rho /(bM)\right)^2}{c_2^2+\left(\frac{qE}{bM}\right)^2+\left(\frac{qE}{bM}\right)^4}}\{a^2M^2/\rho c_2^1(l1/2)c_2^1\left({\displaystyle \frac{qE}{bM}}\right)^2`$
$`+(qH/\rho )\mathrm{coth}(c_2qH/\rho )+2c_2^1{\displaystyle \frac{\left[1+\left(\frac{qE}{bM}\right)^2\right]\left[c_2^2+\left(\frac{qE}{bM}\right)^4\right]}{c_2^2+\left(\frac{qE}{bM}\right)^2+\left(\frac{qE}{bM}\right)^4}}\}n_{sc}^{(2)}.`$
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# SU(3) Symmetry and Scissors Mode Vibrations in Nuclei
## 1 Introduction
In the long history of the shell model, Elliott was the first to point out the advantage of using a deformed (intrinsic) many-body basis and developed the $`SU(3)`$ Shell Model for $`sd`$-shell nuclei. In this model, the classification of basis states and their projection onto good angular momentum can be carried out using group theory. It works nicely as long as the spin-orbit force is weak. In heavier nuclei, where the presence of a strong spin-orbit force is essential for the correct shell closures, the high-$`j`$ orbital in an $`N`$-shell is strongly pushed down by this interaction and intrudes into a lower $`N`$-shell. This results in a re-classification of basis states and the Elliott $`SU(3)`$ scheme is no longer valid for heavy nuclei.
However, in early work of Bhatt, Parikh and McGrory , it was shown that wave functions for yrast states of deformed nuclei in the $`1f`$-$`2p`$ shell that were obtained in shell model calculations within a severely truncated configuration space that cannot support microscopic $`SU(3)`$ symmetry of Elliott, have $`SU(3)`$-like structures when the wave function of a state with total angular momentum $`I`$ is expressed in terms of coupling of the collective neutron and proton angular momenta $`I_\nu `$ and $`I_\pi `$ within the truncated configuration space. On the basis of this analysis, they suggested that the rotational features of yrast states in heavy, deformed nuclei could be due to a similar manifestation of $`SU(3)`$ symmetry occurring at the macroscopic level, even if it is absent at the single particle level in the usual basis. Brink et al. suggested an alternative derivation without the use of symmetry arguments as well as the assumption of the motion of a particle in a static deformed field. Instead they approached the problem by diagonalising a rotationally invariant Hamiltonian using basis states of good angular momentum at all stages, which leads to the standard ‘symmetrised’ wavefunctions.
In our recent letter , the connection between the numerically calculated states of heavy nuclear system and such macroscopic $`SU(3)`$ symmetry states has been explored. The work was based on an extended Projected Shell Model (PSM). Starting from a deformed potential, instead of employing only one BCS vacuum as in the original Projected Shell Model , we use independent BCS condensates for neutrons and protons to construct separate bases for neutron and proton collective rotational states by exact angular momentum projection. The active nucleons in the model space were permitted to occupy three full major shells, which is a space far larger than a conventional shell model can handle. These rotational states were then coupled by diagonalizing a residual pairing plus quadrupole interaction. Many new bands emerge that are not contained in the original PSM, and it was shown that these bands exhibit a nearly perfect $`SU(3)`$ spectrum, even though there is no explicit dynamical symmetry in such a model.
Our work demonstrates a clear connection between a numerical microscopic method, which starts from certain set of intrinsic states, and an elegant group theoretical description of nuclear states. Thus the PSM may be viewed as a natural method to investigate the emergence of $`SU(3)`$ symmetry as a truly dynamical symmetry in heavy, deformed nuclear systems.
The present paper elaborates on our previous letter . The paper is organized as follows: In Section 2, we introduce the extended PSM with an emphasis on how to treat coupled rotations of neutron and proton systems. It is shown that this treatment generalizes the original PSM by allowing for relative motion between the neutron and proton systems. In Section 3, we give details of the calculation for an example, <sup>168</sup>Er. A much richer spectrum is obtained after diagonalization than from the usual PSM. With the help of group theory, the states obtained from the diagonalization are shown to possess a strong $`SU(3)`$ symmetry. We further demonstrate that these states can be classified as rotational bands built on spin-$`1\mathrm{}`$ phonon excitations, and we suggest that these correspond to a geometrical scissors mode and its generalizations. These states are interpreted in Section 4. In Section 5, we buttress our conclusions by comparing the angular momentum distribution in our deformed basis states and the intrinsic states of the $`SU(3)`$ model, and by evaluating overlaps of our total wave functions with the $`SU(3)`$ ones. Then we calculate $`B(E2)`$ values from the PSM wave functions to further confirm the state classifications. Some comments about the emergent $`SU(3)`$ symmetry and about a possible relation to scissors mode vibrations are made in Section 6 and 7, respectively. The lowest excitation bands, the first 1<sup>+</sup> and the first 2<sup>+</sup> band, as well as the second 0<sup>+</sup> band, are further discussed in Section 8. Finally, we summarize the paper in Section 9.
## 2 Extension of the Projected Shell Model
The shell model is the most fundamental way of describing many-nucleon systems fully quantum mechanically. However, using the shell model to study deformed heavy nuclei is a desirable but very difficult task because of large dimensionality and its related problems. Even with today’s computer power, the best standard shell model diagonalization can be done only in the full $`pf`$-shell space, as for example in the work of the Strasbourg–Madrid group , for which the dimension of the configuration space is well over one million. The PSM provides one possible solution for this difficulty. In this approach, one first truncates the configuration space with guidance from the deformed mean field by selecting only the BCS vacuum plus a few quasiparticle configurations in the Nilsson orbitals around the Fermi surface, performs angular momentum projection (and particle number projection if required) to obtain a set of laboratory-frame basis states, and finally diagonalizes a shell-model Hamiltonian in this space. Since the deformed mean field + BCS vacuum already incorporates strong particle-hole and particle-particle correlations, this truncation should be appropriate for the low-lying states dominated by quadrupole and pairing collectivity. Indeed, this approach has been very successful for ground band (g-band) properties and near-yrast quasiparticle excitations in high-spin physics for both normally deformed and superdeformed states .
However, in this formalism the basis vacuum is the usual BCS condensate of neutrons and protons in a fixed deformed potential (fixed shape in the geometric models). Without quasiparticle excitations, one obtains only the g-band after angular momentum projection. Thus, there is no room for studying any other collective excitations.
It is known that small perturbations of nuclear shapes and relative orientations around the equilibrium can give rise to physical states at low to moderate excitation energies. Classical examples of such motion include $`\beta `$\- and $`\gamma `$-vibrations , in which neutrons and protons undergo vibrations as a collective system. These small-amplitude motions are not built into the ground state for theories like Hartree–Fock–Bogoliubov or BCS. One may obtain the $`\beta `$\- and $`\gamma `$-vibrations by building additional correlations into the ground state, as for example in the Random Phase Approximation (see, for example, Ref. ), or by enriching the intrinsic degrees of freedom in the deformed potential .
Another possible mode of collective motion results if we consider neutron and proton fields of the fixed deformation, but permit small oscillations in the relative orientations of these fields . This geometric picture may be related to the two-rotor model . Because of the strong restoring force (with its physical origin in the neutron-proton interaction), this oscillation is confined to a small angle between the protons and neutrons (thus it is termed small-amplitude scissors motion). This vibration, and the $`\beta `$\- and $`\gamma `$-vibrations, can be classified using group theoretical methods; they belong to the lowest collective excitations of the ground state, as pointed out by Iachello .
As introduced in our previous letter , in order to study microscopically the relative motion between neutron and proton intrinsic states with a fixed deformation, we have extended the original PSM in the following manner: Instead of a single BCS vacuum, the angular momentum projection is now performed for separate neutron and proton deformed BCS vacua. Although the introduction of two separately projected BCS vacua seems to treat neutrons and protons as two independent systems, the equal deformation used in the Nilsson calculation of the basis already contains strong correlation between the two systems . Further proton-neutron correlations are introduced by explicitly diagonalizing the total Hamiltonian in the basis formed by the angular momentum projected neutron and proton states. This procedure gives the usual g-band corresponding to the coherently coupled BCS condensate of neutrons and protons, but also leads to a new set of states built on a more complex vacuum that incorporates fluctuations in the relative orientation of the neutron and proton fields of constant deformation.
The neutron and proton valence spaces employed in the present work are those of Ref. . Our single particle space contains three major shells ($`N`$ = 4, 5, and 6) for neutrons and ($`N`$ = 3, 4 and 5) for protons; this has been shown to be sufficient for a quantitative description of rare-earth g-bands and bands built on a few quasiparticle excitations . Separable forces of pairing plus quadrupole type are used in our Hamiltonian. We note that a recent study by Dufour and Zuker has shown explicitly that the residual part of the realistic force is strongly dominated by pairing and quadrupole interactions.
The Hamiltonian can be expressed as $`\widehat{H}=\widehat{H}_\nu +\widehat{H}_\pi +\widehat{H}_{\nu \pi }`$, where $`H_\tau `$ $`(\tau =\nu ,\pi )`$ is the like-particle pairing plus quadrupole Hamiltonian, with inclusion of quadrupole-pairing,
$$\widehat{H}_\tau =\widehat{H}_\tau ^0\frac{\chi _{\tau \tau }}{2}\underset{\mu }{}\widehat{Q}_\tau ^\mu \widehat{Q}_\tau ^\mu G_\text{M}^\tau \widehat{P}_\tau ^{}\widehat{P}_\tau G_\text{Q}^\tau \underset{\mu }{}\widehat{P}_\tau ^\mu \widehat{P}_\tau ^\mu ,$$
(1)
and $`\widehat{H}_{\nu \pi }`$ is the n–p quadrupole–quadrupole residual interaction
$$\widehat{H}_{\nu \pi }=\chi _{\nu \pi }\underset{\mu }{}\widehat{Q}_\nu ^\mu \widehat{Q}_\pi ^\mu .$$
(2)
In Eq. (1), the four terms are, respectively, the spherical single-particle energy which contains a proper spin-orbit force , the quadrupole-quadrupole interaction, the monopole-pairing, and the quadrupole-pairing interaction. The interaction strengths $`\chi _{\tau \tau }`$ ($`\tau =\nu `$ or $`\pi `$) are related self-consistently to the deformation $`\epsilon `$ by
$$\chi _{\tau \tau }=\frac{\frac{2}{3}\epsilon (\mathrm{}\omega _\tau )^2}{\mathrm{}\omega _\nu \widehat{Q}_0_\nu +\mathrm{}\omega _\pi \widehat{Q}_0_\pi }.$$
(3)
Obviously, neutrons and protons are coupled by the self-consistency condition. Following Ref. , the strength $`\chi _{\nu \pi }`$ is assumed to be $`\chi _{\nu \pi }=(\chi _{\nu \nu }\chi _{\pi \pi })^{1/2}`$. Similar parameterizations were used in earlier work .
In this paper, we concentrate on collective motion only and do not consider quasiparticle excitations; however, we shall discuss possible physical implication of including quasiparticle states in a later section. Particle number projection is not included, so particle number is conserved only on average in the BCS calculation. It has been shown by a systematic calculation that particle number projection does not alter the unprojected results very much for those well-deformed heavy nuclei. We do not expect that particle number projection will change conclusions of this paper with respect to the emergence of $`SU(3)`$ symmetry.
The intrinsic state $`|0`$ of an axially deformed nucleus is, in this approximation, taken to be the product of the Nilsson-BCS quasiparticle vacuum states $`|0_\nu `$ and $`|0_\pi `$
$$|0=|0_\nu |0_\pi .$$
(4)
The original PSM ground band with angular momentum $`|I`$ is obtained by angular momentum projection onto the vacuum:
$$|I=\mathrm{}^I\widehat{P}^I|0,$$
(5)
where $`\widehat{P}^I`$ is the (one-dimensional) angular momentum projection operator
$$\widehat{P}^I=(I+\frac{1}{2})_0^\pi 𝑑\beta sin\beta d^I(\beta )\widehat{R}(\beta ),$$
(6)
with $`d^I(\beta )`$ the small $`d`$-function and $`\widehat{R}(\beta )`$ the one-dimensional rotational operator and $`\mathrm{}^I`$ is the normalization constant,
$$\mathrm{}^I=0|\widehat{P}^I|0^{\frac{1}{2}}.$$
(7)
In the extended PSM, we first project out the neutron (proton) states $`|I_\nu `$ ($`|I_\pi `$) with angular momentum $`I_\nu `$ ($`I_\pi `$) from the intrinsic state $`|0_\nu `$ ($`|0_\pi `$). The projected states $`|I_\nu `$ and $`|I_\pi `$ are coupled to form the basis states $`|(I_\nu I_\pi )I`$ for different total angular momentum $`I`$. These basis states are used to construct the matrix of the total Hamiltonian of Eqs. (1) and (2),
$`(I_\nu I_\pi )I|`$ $`\widehat{H}`$ $`|(I_\nu ^{}I_\pi ^{})I`$ (8)
$`=`$ $`\left[I_\nu |\widehat{H}_\nu |I_\nu ^{}+I_\pi |\widehat{H}_\pi |I_\pi ^{}\right]\delta _{I_\nu I_\nu ^{}}\delta _{I_\pi I_\pi ^{}}`$
$``$ $`\chi _{\nu \pi }(I_\nu I_\pi )I|\widehat{Q}_\nu ^{}\widehat{Q}_\pi |(I_\nu ^{}I_\pi ^{})I.`$
The last term in Eq. (8) can be written explicitly as
$`(I_\nu `$ $``$ $`I_\pi )I|\widehat{Q}_\nu ^{}\widehat{Q}_\pi |(I_\nu ^{}I_\pi ^{})I`$ (9)
$`=`$ $`𝒲(I_\pi 2II_\nu ^{};I_\pi ^{}I_\nu )I_\nu \widehat{Q}_\nu I_\nu ^{}`$
$`\times `$ $`I_\pi \widehat{Q}_\pi I_\pi ^{}/\sqrt{(2I_\nu +1)(2I_\pi ^{}+1)},`$
where $`𝒲`$ is the 6-$`j`$ symbol.
The term $`I_\nu |\widehat{H}_\nu |I_\nu `$ ($`I_\pi |\widehat{H}_\pi |I_\pi `$) is the energy of the state $`|I_\nu `$ ($`|I_\pi `$) projected from the intrinsic state $`|0_\nu `$ ($`|0_\pi `$) with the neutron (proton) part of Hamiltonian $`\widehat{H}_\nu `$ ($`\widehat{H}_\pi `$) given by Eq. (1). The Hamiltonian matrix of Eq. (8) is diagonalized and the resulting PSM eigenstates $`|\alpha ,I_{\text{PSM}}`$ are expressed as a linear combination of the basis states $`|(I_\nu I_\pi )I`$:
$$|\alpha ,I_{\text{PSM}}=\underset{I_\nu I_\pi }{}f_{\text{PSM}}^\alpha (I_\nu I_\pi ;I)|(I_\nu I_\pi )I,$$
(10)
where $`\alpha `$ labels different eigenstates having the same angular momentum. We shall show in the following sections that the amplitudes $`f_{\text{PSM}}^\alpha (I_\nu I_\pi ;I)`$ carry information about $`SU(3)`$ symmetry emerging dynamically from the extended PSM.
## 3 An Example: The Spectrum of <sup>168</sup>Er
Let us take as a typical example the rotational nucleus, <sup>168</sup>Er. The deformed basis is constructed at deformation $`\epsilon _2=0.273`$ for this nucleus . The numerical values of the parameters in the Hamiltonian (Eqs. (1) and (2)) appropriate for this nucleus are (in MeV): $`\chi _{\nu \nu }=0.0206`$, $`\chi _{\pi \pi }=0.0160`$, $`\chi _{\nu \pi }=0.0182`$, $`G_M^\nu =0.1049`$, and $`G_M^\pi =0.1346`$. The ratio of the strength of the quadrupole-pairing interaction to that of the monopole-pairing interaction is 0.16. All these values are the same as those employed in the early paper of Hara and Iwasaki and as those in the review article of Hara and Sun . The states $`|I_\nu `$ ($`|I_\pi `$) with angular momenta $`I_\nu =0,2,\mathrm{},32`$ ($`I_\pi =0,2,\mathrm{},16`$) were projected from the Nilsson-BCS vacuum state $`|0_\nu `$ ($`|0_\pi `$). We have omitted in this calculation the neutron states with $`I_\nu >32`$ and the proton states with $`I_\pi >16`$ because their probabilities in the intrinsic states are very small ($`<`$ 0.001).
The energy spectra $`E(I_\tau )=I_\tau |\widehat{H}_\tau |I_\tau `$ obtained from the projected states $`|I_\tau `$ are almost rotational ( $`E(I_\tau )=A_\tau I_\tau (I_\tau +1)`$ ), with the moment of inertia parameters $`A_\nu =`$ 0.019 MeV and $`A_\pi =`$ 0.048 MeV.
For each spin $`I`$, we diagonalize Eq. (8) in the basis $`|(I_\nu I_\pi )I`$. The resulting spectrum with energy up to 20 MeV and angular momenta up to 12$`\mathrm{}`$ is shown in Fig. 1. It shows strikingly regular structure.
The lowest band is the g-band, which is nearly identical to that of earlier PSM calculations with the same Hamiltonian where no separation and re-coupling of neutron and proton components was considered. The calculated g-band in Fig. 1 reproduces the experimental g-band very well, as we shall show later in Fig. 4. In addition to the g-band, many new excited bands emerge that are not found in the earlier PSM calculations. These bands exhibit a curvature similar to the g-band, suggesting that their moments of inertia are similar to the g-band.
## 4 Emergence of SU(3) Symmetry
The strikingly regular pattern shown in Fig. 1 can be understood as manifestation of a nearly perfect $`SU(3)`$ symmetry: all bands can be well classified as a spectrum with $`SU(3)`$ symmetry if the projected neutron and proton BCS vacuum states are considered to be two independent $`SU(3)`$ representations coupled through the $`Q_\nu `$-$`Q_\pi `$ interaction.
This may be demonstrated by considering a model Hamiltonian with $`SU(3)^\nu SU(3)^\pi SU(3)^{\nu +\pi }`$ dynamical symmetry expressed in the form (see Eq. (3.107) in Ref. ):
$$\widehat{H}=\chi _\nu ^{\text{eff}}\widehat{C}_{su3}^\nu +\chi _\pi ^{\text{eff}}\widehat{C}_{su3}^\pi \chi _{\nu \pi }^{\text{eff}}\widehat{C}_{su3}^{\nu +\pi }+A\widehat{I}^2,$$
(11)
where $`\widehat{C}_{su3}^\tau `$ are the $`SU(3)^\tau `$ ($`\tau =\nu ,\pi `$) Casimir operators for neutrons, protons, and the n–p coupled symmetry $`SU(3)^{\nu +\pi }`$ ($`\tau =\nu +\pi `$). The eigenvalue of the lowest-order $`SU(3)`$ Casimir operator for a given representation $`(\lambda ,\mu )`$ is $`C(\lambda ,\mu )=\frac{1}{2}(\lambda ^2+\mu ^2+\lambda \mu +3\lambda +3\mu )`$. We assume that the two BCS vacua correspond to the $`SU(3)`$ symmetric representations $`(\lambda _\nu ,0)`$ and $`(\lambda _\pi ,0)`$, respectively, and that the permissible irreps $`(\lambda ,\mu )`$ of $`SU(3)^{\nu +\pi }`$ are given by the Littlewood rule, $`(\lambda ,\mu )=(\lambda _m2\mu ,\mu )`$, where $`\lambda _m=\lambda _\nu +\lambda _\pi `$ is the maximum value of $`\lambda `$ and $`\mu =0,1,2,\mathrm{},\lambda _\pi `$, if $`\lambda _\pi \lambda _\nu `$.
The eigenvalue spectrum $`E[(\lambda ,\mu )I]`$ of the Hamiltonian of (11) obtained from the states with angular momentum $`I`$ belonging to the coupled representation $`\{[\lambda _\nu ,0][\lambda _\pi ,0]\}(\lambda ,\mu )I`$ is
$`E[(\lambda ,\mu )I]E_{\text{g.s.}}`$ $`=`$ $`\chi _{\nu \pi }^{\text{eff}}\left[C(\lambda _m,0)C(\lambda =\lambda _m2\mu ,\mu )\right]+AI(I+1)`$ (12)
$`=`$ $`\mu \mathrm{}\omega _{\mathrm{}}\left[1{\displaystyle \frac{\mu 1}{\lambda _m}}\right]+AI(I+1),`$
with
$$\mathrm{}\omega _{\mathrm{}}=\frac{3}{2}\lambda _m\chi _{\nu \pi }^{\text{eff}}.$$
(13)
The second term in Eq. (12) gives a rotational band with a moment of inertia parameter $`A=\mathrm{}^2/2\mathrm{}`$. $`E_{\text{g.s.}}`$ is the ground state energy corresponding to the coupled $`SU(3)`$ representation $`[\lambda _m,0]`$
$$E_{\text{g.s.}}=\chi _\nu ^{\text{eff}}C(\lambda _\nu ,0)+\chi _\pi ^{\text{eff}}C(\lambda _\pi ,0)\chi _{\nu \pi }^{\text{eff}}C(\lambda _m,0).$$
(14)
The parameter $`\lambda _m`$ fixes the $`SU(3)`$ representation $`[\lambda _m,0]`$ of the g-band. The parameter $`\mu `$ then fixes the representation $`[\lambda =\lambda _m2\mu ,\mu ]`$ of the excited bands through the Littlewood rule.
A more physical meaning of the parameter $`\mu `$ emerges from the first term of Eq. (12), which represents a vibration-like spectrum which tends to be harmonic for $`\lambda _m\mu `$. In this limit, $`\mu `$ can be interpreted as the number of phonons of this “vibration” and $`\mathrm{}\omega _{\mathrm{}}`$ as the one-phonon excitation energy as $`\lambda _m\mathrm{}`$, with its value being, approximately, the energy of the first excited $`I=1^+`$ state relative to the ground state. The allowed angular momenta $`I`$ belonging to a representation $`[\lambda ,\mu ]`$ are determined from the usual $`SU(3)`$ subgroup reduction rules . For example, the representation $`(\lambda ,\mu )=(40,4)`$ can have a $`K=0`$ band with $`I=0,2,4,\mathrm{},44`$, a $`K=2`$ band with $`I=2,3,4,\mathrm{},43`$ and a $`K=4`$ band with $`I=4,5,6,\mathrm{},41`$ .
The PSM calculated spectrum shown in Fig. 1 can then be described by the $`SU(3)`$ symmetry spectrum of Eq. (12). The parameters $`\lambda _m`$, $`\mathrm{}\omega _{\mathrm{}}`$ and $`A`$ of Eq. (12), which fit best the PSM spectrum of Fig. 1, are $`\lambda _m=48`$, $`\mathrm{}\omega _{\mathrm{}}=2.9`$ MeV ($`\chi _{\nu \pi }^{\text{eff}}=0.0403`$), and $`A=0.013`$ MeV. The $`SU(3)`$ band structure obtained with these parameters is shown in Fig. 1 as the dashed lines and reproduces remarkably well the PSM spectrum obtained from numerical projection and diagonalization. The fact that the parameter $`\lambda _m=48`$ fits the spectrum in Fig. 1 suggests that the g-band of PSM may be associated with the $`SU(3)`$ representation $`[\lambda _m,0]=[48,0]`$. All the states shown in Fig. 1 can be labeled by the $`SU(3)`$ irrep labels $`(\lambda ,\mu )`$ and the bandhead of each rotational band within the irrep is labeled by the $`SU(3)`$ quantum number $`K`$. Degeneracy at each spin may be deduced by counting one for each $`K`$ band, except only even spins are present for $`K=0`$ bands. We list the degeneracy of the $`(40,4)`$ representation as an example in Fig. 1. Not all PSM states can be seen clearly in the plot because of the high degeneracy, but there is a one-to-one correspondence between all predicted $`SU(3)`$ states and those observed in the PSM calculation.
One can see from Eq. (12) and the $`SU(3)`$ reduction rules that the whole spectrum of Fig. 1 can be viewed as a set of rotational bands built on different multi-phonon excitation states with a phonon energy $`\mathrm{}\omega =\mathrm{}\omega _{\mathrm{}}\left[1\frac{\mu 1}{\lambda _m}\right]`$ and phonon spin $`1\mathrm{}`$. For example, a 3-phonon system could have two allowed states with energy $`3\mathrm{}\omega `$ and total spin 1$`\mathrm{}`$ and 3$`\mathrm{}`$; a 4-phonon system could have three states with energy $`4\mathrm{}\omega `$ and total spin 0$`\mathrm{}`$, 2$`\mathrm{}`$, and 4$`\mathrm{}`$; and so on. This provides an alternative explanation of the degeneracy of the bands obtained by the PSM diagonalization. Comparing the $`SU(3)`$ and phonon classifications, one can see clearly that the $`SU(3)`$ quantum numbers $`\mu `$ and $`K`$ in Fig. 1 may be interpreted as the number of phonons and their allowed total spins, respectively.
As long as $`\mathrm{}\omega _{\mathrm{}}`$ is held constant, the spectrum is sensitive to the effective $`SU(3)`$ quantum number $`\lambda _m`$ only through a small anharmonicity. For the present example, the phonon energy decreases smoothly from 2.9 MeV to 2.5 MeV as the phonon number increases from 1 to 7. If $`\lambda _m\mathrm{}`$, then $`\mathrm{}\omega \mathrm{}\omega _{\mathrm{}}`$ and the vibration becomes harmonic. Thus, the anharmonicity originates in the finite number of particles for the nuclear system.
## 5 Further Inspection on the Wave Functions
The remarkably quantitative agreement of all calculated excited states with the $`SU(3)`$ spectrum demonstrated above makes it essentially certain that an $`SU(3)`$ symmetry emerges from the purely numerical PSM calculation. To remove all doubt, we may study the wave functions. In this section, we shall do this in three steps: we first compare angular momentum components in the PSM and $`SU(3)`$ intrinsic states. Then, we directly calculate overlaps of the actual wave functions from the PSM calculation with $`SU(3)`$ wave functions. Finally, we compute $`B(E2)`$ values from PSM calculations that connect states in each band and compare with those of the $`SU(3)`$ model.
### 5.1 SU(3) Structure of the Intrinsic States
Let us start by recalling that an intrinsic state $`|0_\tau `$, $`(\tau =\nu ,\pi `$), is related to a projected state $`|I_\tau `$ by
$$|0_\tau =\underset{I_\tau }{}C_{I_\tau }|I_\tau .$$
(15)
In Eq. (15), $`C_I`$ is the norm matrix element in the PSM described in Eq. (2.33) of Ref.
$$|C_I|^2=0|\widehat{P}^I|0.$$
(16)
It gives the probability of finding angular momentum $`I`$ in the intrinsic state $`|0`$ (with $`_I|C_I|^2=1`$, a sum rule obtained in Eq. (A80) of Ref. ). In other words, the quantity $`|C_I|^2`$ describes the angular momentum distribution in the intrinsic state $`|0`$. For the present calculation, we projected the states with angular momenta $`I_\nu =0,2,\mathrm{},32`$ and $`I_\pi =0,2,\mathrm{},16`$ from the Nilsson-BCS intrinsic states $`|0_\tau `$ of <sup>168</sup>Er. The probabilities $`|C_{I_\tau }|^2`$ for $`I_\nu >32`$ and $`I_\pi >16`$ were very small and hence such states were not included in the basis state $`|[I_\nu I_\pi ]I`$ used in the present PSM diagonalization. The values $`|C_{I_\tau }|^2`$ (for the nucleus <sup>168</sup>Er) calculated for the first nine values $`I_\tau =0,2,\mathrm{},16`$ contained in the intrinsic state $`|0_\tau `$ are listed in Table 1. For comparison, the values $`|C_I|^2`$ for the state $`|I`$ in the total intrinsic vacuum state $`|0_\nu |0_\pi `$ are also presented in the same table.
In order to make a connection with the emergent $`SU(3)`$ symmetry at the intrinsic state level, we compare the probabilities $`C_{I_\tau }^2`$(PSM) of the PSM intrinsic states $`|0_\tau _{\text{PSM}}`$ with the corresponding probabilities $`C_{I_\tau }^2`$($`SU(3)`$) in the intrinsic states $`|0_\tau _{SU(3)}`$ belonging to axially symmetric $`SU(3)`$ representation $`[\lambda _\tau ,0]`$. The latter probabilities are given by Elliott’s $`a((\lambda ,0)I)`$ coefficients . Analytical formulas for the $`a((\lambda ,\mu )I)`$ are given in Table 2A of the Vergados’s paper . The effective $`SU(3)`$ representation $`[\lambda _\tau ,0]`$ for protons and neutrons is determined by equating the probabilities $`C_{I_\tau }^2`$(PSM) = $`C_{I_\tau }^2`$($`SU(3)`$) = $`a((\lambda _\tau ,0)I_\tau )`$ for a given $`I_\tau `$ (say $`I_\tau =0`$). For our present case of <sup>168</sup>Er, the values of $`\lambda _\tau `$ obtained by this procedure are $`\lambda _\nu `$ 32 and $`\lambda _\pi `$ 16, to the nearest even integer.
The coefficients $`a((\lambda _\tau ,0)I)`$ are then calculated for all values of $`I_\tau =0,2,\mathrm{},\lambda _\tau `$, allowed by the effective $`\lambda _\tau `$. (We denote these values as $`a_I`$ in the following discussion.) The $`a_{I_\tau }`$ values are compared with the corresponding $`C_{I_\tau }`$ values of the PSM in Fig. 2. The excellent agreement between these values shows that the distribution of angular momenta of neutrons and protons in our Nilsson–BCS vacuum states are very similar to those in the intrinsic states of the $`SU(3)`$-representations $`[\lambda _\nu =32,0]`$ and $`[\lambda _\pi =16,0]`$.
The total Nilsson–BCS vacuum state $`|0_{\text{PSM}}=|0_\pi _{\text{PSM}}|0_\nu _{\text{PSM}}`$ should then have an effective $`SU(3)`$-representation $`[\lambda =\lambda _\pi +\lambda _\nu ,0]=[48,0]`$ and the probabilities $`C_I^2`$ of finding angular momentum $`I`$ in the total PSM intrinsic state should be similar to the probabilities $`a((\lambda =48,0)I)`$ corresponding to the $`SU(3)`$ intrinsic state of representation . The agreement between these two probabilities is also shown in Fig. 2.
These PSM effective $`SU(3)`$ representations for the Nilsson–BCS intrinsic states of neutrons and protons for <sup>168</sup>Er in the three-shell valence space should be compared with the effective representations $`[\lambda _\nu ^{\text{eff}},\mu _\nu ^{\text{eff}}]=[40,0],[\lambda _\pi ^{\text{eff}},\mu _\pi ^{\text{eff}}]=[24,0]`$, and $`[\lambda ^{\text{eff}},\mu ^{\text{eff}}]=[64,0]`$ obtained by Kahane et al. for just the Nilsson intrinsic state of <sup>168</sup>Er in the $`\nu (82126)`$ and $`\pi (5082)`$ single major shell valence space. These representations are larger than the effective representations $`[\lambda _\nu =32,0]`$ and $`[\lambda _\pi =16,0]`$ mainly because Kahane et al. have included all the projected states with probabilities $`|C_I|^2<0.001`$, which have been omitted in the present PSM calculations. For comparison, we also note that in the one major shell valence space, the pseudo-$`SU(3)`$ representations for the normal parity nucleons in $`{}_{68}{}^{}{}_{}{}^{168}`$Er<sub>100</sub> are $`[\lambda _\nu ^{\text{pseudo}},\mu _\nu ^{\text{pseudo}}]=[20,4]`$, $`[\lambda _\pi ^{\text{pseudo}},\mu _\pi ^{\text{pseudo}}]=[10,4]`$, and $`[\lambda _n^{\text{pseudo}}=30,\mu _n^{\text{pseudo}}=8]`$. These representations are strongly triaxial but their effective axially symmetric representations would be $`[\lambda _\nu ^{\text{eff}}=24,0]`$, $`[\lambda _\pi ^{\text{eff}}=14,0]`$ and $`[\lambda ^{\text{eff}}=38,0]`$. In the FDSM , the $`SU(3)`$ representation for 10 neutrons in the normal parity states of the (82 – 126) shell is $`[\lambda _\nu ^{\text{FDSM}},\mu _\nu ^{\text{FDSM}}]=[10,0]`$. There are also 10 protons in normal parity states of the (50 – 82) shell, and they do not show an $`SU(3)`$, but an $`SO(6)`$ symmetry that corresponds to a $`\gamma `$-soft rotor. However, when the protons are strongly coupled to the neutrons, they will be synchronized with neutrons and stabilized at $`\gamma 0`$, so that the proton system behaves like an $`SU(3)`$ rotor. Therefore, an effective $`SU(3)`$ representation $`[\lambda _\pi ^{\text{FDSM}},\mu _\pi ^{\text{FDSM}}]=[10,0]`$ could be assigned to the protons. Quantitatively, the FDSM with one major shell can only fit to the first three bands in Fig. 1. Larger deviations are expected for the higher bands, although the qualitative band structure remains correct. If the $`SU(3)`$ symmetry which appears to emerge from the PSM calculation presented here is identified as the FDSM symmetry in the extended three-shell valence space with the effective neutron and proton numbers $`n_\nu ^{\text{eff}}=32`$ and $`n_\pi ^{\text{eff}}=16`$, then agreement with the PSM results is perfect, as shown in Fig. 1.
### 5.2 SU(3) Symmetry in the Wave Functions
Next we shall verify that the structure of the PSM wave function $`|I_{\text{PSM}}`$ belonging to different bands obtained by diagonalization is similar to that of the $`SU(3)`$ wave function $`|(\lambda ,\mu ),I_{SU(3)}`$ belonging to the $`SU(3)`$ representation $`[\lambda ,\mu ]`$ obtained by coupling the $`[\lambda _\pi ,0]`$ and $`[\lambda _\nu ,0]`$. Here we consider only the g-band. The PSM wave functions for the g-band obtained by diagonalizing the Hamiltonian in the basis $`\{(I_\nu I_\pi )I\}`$ can be written as
$$|I_{\text{PSM}}=\underset{I_\nu ,I_\pi }{}f_{\text{PSM}}(I_\nu I_\pi ;I)|(I_\nu I_\pi )I,$$
(17)
and the total $`SU(3)`$ wave function for the g-band representation $`[\lambda ,0]`$ can be written as
$`|(\lambda ,\mu ),I_{SU(3)}`$ $`=`$ $`\mathrm{}_I{\displaystyle \underset{I_\nu ,I_\pi }{}}a_{I_\nu }a_{I_\pi }I_\nu 0I_\pi 0|I0|[I_\nu I_\pi ]I`$ (18)
$`=`$ $`{\displaystyle \underset{I_\nu ,I_\pi }{}}f_{SU(3)}(I_\nu I_\pi ;I)|(I_\nu I_\pi )I,`$
where $`\mathrm{}_I`$ is the normalization constant
$$|\mathrm{}_I|^2=\frac{1}{_{I_\nu ,I_\pi }(a_{I_\nu }a_{I_\pi }I_\nu 0I_\pi 0|I0)^2}.$$
(19)
The overlap of these two wave functions is
$`O_I`$ $`=`$ $`{}_{SU(3)}{}^{}(\lambda ,\mu ),I|I_{\text{PSM}}^{}`$ (20)
$`=`$ $`{\displaystyle \underset{I_\nu ,I_\pi }{}}f_{SU(3)}(I_\nu I_\pi ;I)f_{\text{PSM}}(I_\nu I_\pi ;I).`$
Knowing the $`SU(3)`$ coefficients $`f_{SU(3)}`$ from $`a_{I_\nu }`$ and $`a_{I_\pi }`$, and the PSM coefficients $`f_{\text{PSM}}`$ from numerical diagonalization, we are able to compute their overlap. The calculated overlaps $`O_I`$ are listed in Table 2 <sup>1</sup><sup>1</sup>1A similar calculation for the overlaps for the PSM wave functions $`|I_{PSM}`$ belonging to higher bands with the $`SU(3)`$ wave functions of the representations , , etc. requires calculation of $`SU(3)`$ Clebsch–Gordon coefficients. The proof should be straightforward and is not given here.. All the numbers are very close to 1, indicating a strong overlap of the PSM wave functions with the $`SU(3)`$ wave functions.
### 5.3 SU(3)-like Systematics of $`B(E2,II2)`$ values
The $`B(E2)`$ values are useful to determine the collectivity of a band. In the discussion of the last section, the numerically calculated PSM states were classified using group theoretical methods. We shall now verify this classification with $`B(E2)`$ calculations using PSM wave function, which will allow us to classify PSM states into collective bands. Since we have proved the similarity of the PSM wave functions to the $`SU(3)`$ ones, we expect to obtain similar $`B(E2)`$ values if we make a corresponding calculation using $`SU(3)`$ wave functions.
The $`B(E2)`$ transitions are calculated as
$$B(E2,I_iI_f)=\frac{|I_f𝒯^2I_i|^2}{2I_i+1},$$
(21)
where the operator $`𝒯^2=𝒯_\nu ^2+𝒯_\pi ^2`$ is related to the quadrupole operators by
$`𝒯_\nu ^2`$ $`=`$ $`e_\nu ^{\text{eff}}\sqrt{{\displaystyle \frac{5}{16\pi }}}Q_\nu ^2,`$
$`𝒯_\pi ^2`$ $`=`$ $`e_\pi ^{\text{eff}}\sqrt{{\displaystyle \frac{5}{16\pi }}}Q_\pi ^2.`$ (22)
In the calculations, we have used the usual effective charges $`e_\nu ^{\text{eff}}=0.5e`$ and $`e_\pi ^{\text{eff}}=1.5e`$. By employing the PSM wave functions of Eq. (17), the reduced matrix element appearing in Eq. (21) can be evaluated as
$`I_f𝒯^2I_i`$ $`=`$ $`{\displaystyle \underset{I_f(I_{f,\nu },I_{f,\pi })}{}}{\displaystyle \underset{I_i(I_{i,\nu },I_{i,\pi })}{}}f_f(I_{f,\nu }I_{f,\pi };I_f)f_i(I_{i,\nu }I_{i,\pi };I_i)`$ (23)
$`(I_{f,\nu }I_{f,\pi })I_f𝒯_\nu ^2+𝒯_\pi ^2(I_{i,\nu }I_{i,\pi })I_i,`$
in which the matrix elements for a coupled system can be explicitly expressed as
$`[I_{f,\nu }I_{f,\pi }]I_f`$ $`𝒯_\nu ^2[I_{i,\nu }I_{i,\pi }]I_i=\sqrt{(2I_i+1)(2I_f+1)}`$
$`()^{I_{f,\nu }+I_{f,\pi }+I_i}𝒲(I_{i,\nu }I_{i,\pi }2I_f;I_iI_{f,\nu })I_{f,\nu }𝒯_\nu ^2I_{i,\nu }\delta _{I_{f,\pi }I_{i,\pi }},`$
$`[I_{f,\nu }I_{f,\pi }]I_f`$ $`𝒯_\pi ^2[I_{i,\nu }I_{i,\pi }]I_i=\sqrt{(2I_i+1)(2I_f+1)}`$ (24)
$`()^{I_{i,\nu }+I_{i,\pi }+I_f}𝒲(I_{i,\nu }I_{i,\pi }I_f2;I_iI_{f,\pi })I_{f,\pi }𝒯_\pi ^2I_{i,\pi }\delta _{I_{f,\nu }I_{i,\nu }}.`$
In the $`SU(3)`$ model, assuming the transition operator $`𝒯^2=\alpha P^2`$, where $`P^2`$ is the generator of the coupled $`SU(3)^{\nu +\pi }`$, the $`B(E2)`$ formula for the intra-band transition is
$$B(E2,I_iI_f)=\frac{2I_f+1}{2I_i+1}\alpha ^2\left|(\lambda \mu )KI_i,(11)2(\lambda \mu )KI_f\right|^2C(\lambda ,\mu ),$$
(25)
where $`(\lambda \mu )KI_i,(11)2(\lambda \mu )KI_f`$ is the $`SU(3)R(3)`$ Wigner coefficient, and $`\alpha `$ is a parameter of the effective transition operator in the $`SU(3)`$ model determined by fitting data. The inter-band transition rate is zero. In the symmetry limit, if the transition operator is not a generator of $`SU(3)^{\nu +\pi }`$ or there exists small symmetry breaking in the wavefunction, the inter-band transition rate would be expected to be small but not zero.
Generally, the shell-model $`E2`$ operator has different effective charges for neutrons and protons and as such is not a generator of the coupled $`SU(3)^{\nu +\pi }`$. We note, however, that Eq. (25) is still a good approximation for this more general $`E2`$ operator through the inclusion of a factor $`(e_\nu \lambda _\nu +e_\pi \lambda _\pi )/(\lambda _\nu +\lambda _\pi )`$ in the proportionality constant $`\alpha `$. Different neutron and proton effective charges (or gyromagnetic ratios for $`M1`$’s) will also give rise to inter-band transitions, which is an essential mechanism to obtain the strong $`M1`$’s from scissors states to ground band.
The $`B(E2)`$ values calculated by the PSM are shown in Fig. 3, in comparison with the $`SU(3)`$ results and the experimental data . For the g-band (labeled as the first $`0^+`$ band in the tables), good agreement is found for all states except for the transition $`6^+4^+`$. Besides the g-band, we have calculated the $`B(E2)`$’s for the first $`1^+`$ band, the second $`0^+`$ band, and the first $`2^+`$ band (see Fig. 1). The values displayed in Table 3 confirm strong collectivity in each of these intra-transition bands and is in good agreement with what found in the $`SU(3)`$ model. Linking $`B(E2)`$ values between any of two bands are also calculated, with the results displayed in Table 4. These inter-band $`B(E2)`$ values are typically two orders of magnitude smaller than the values for the intra-band transitions. From these calculations, the sets of nearly degenerate bands seen in the PSM calculations of Fig. 1 (for example, the second $`0^+`$ band and the first $`2^+`$ band) can be separated easily. All these results strongly support the conclusion that the structure of well deformed nuclei resulting from PSM calculations indeed possesses a very strong $`SU(3)`$ symmetry.
## 6 Comments on the Emergent SU(3) Symmetry
The connection between $`SU(3)`$ symmetry and nuclear rotation has a long history. It was first explored by Elliott in his classic 1958 papers for $`sd`$-shell nuclei. This idea was later extended to heavy systems using the pseudo-$`SU(3)`$ model , which showed that it is possible to represent approximately the ground band of a deformed nucleus by one leading $`SU(3)`$ representation belonging only to the normal parity nucleons in a single major shell. However this model has difficulty in reproducing $`\beta `$-bands of heavy nuclei, and quantitative results generally do not appear in simple symmetry limits and therefore require numerical calculations. Later, the interacting boson model (IBM) demonstrated that nuclear rotational motion including $`\beta `$\- and $`\gamma `$-bands can be described by an $`SU(3)`$ dynamical symmetry based on $`s`$-$`d`$ or $`s`$-$`d`$-$`g`$ bosons. The fermion dynamical symmetry model (FDSM), based on the $`S`$-$`D`$ fermion pairs of normal parity nucleons in a major shell, demonstrated analytically the equivalence between the particle–rotor model and the $`SU(3)`$ dynamical symmetry limit of the FDSM when particle number $`n\mathrm{}`$; if the Pauli effect is neglected (by assuming the shell pair-degeneracy $`\mathrm{\Omega }\mathrm{}`$) the FDSM reduces to the IBM (see Section 3 and 4 in Ref. ).
However, all the above-mentioned models are algebraic and the possibility for symmetries arises naturally in them. The extended PSM described here has not built an explicit $`SU(3)`$ symmetry into the problem, and no free parameters have been adjusted. Nevertheless, the spectra and wavefunctions obtained from the PSM calculations can again be well classified using the representation theory of the $`SU(3)`$ group. The $`SU(3)`$ symmetry just emerges naturally at the macroscopic level from a shell model diagonalization in the basis obtained by angular momentum projection from a deformed intrinsic state. This is a non-trivial result because our study is not confined in the g-band, but extended up to high spins and high excitations, and because the correspondence has been demonstrated not only at the spectral level, but in the wavefunctions and electromagnetic transition rates too. This strongly indicates that a well deformed nucleus is in fact a very good $`SU(3)`$ rotor.
Now we may ask an important question: what is the nature of the $`SU(3)`$ symmetry that emerges from the PSM? There are many $`SU(3)`$ models. Mathematically, they are equivalent if their representations are the equivalent. Physically, the models differ in the microscopic basis on which the symmetry is built, and have different permissible $`SU(3)`$ irreps, and thus different band structures. The emergent $`SU(3)`$ symmetry we have found here is in principle consistent with any $`SU(3)`$ symmetry model that has allowed representations consistent with the observed states and band structure. We emphasize that this is a physical rather than mathematical criterion, for it is the physical input to an $`SU(3)`$ model that determines its allowed irreps. Therefore, to answer the question of which $`SU(3)`$ model is consistent with the symmetry of the well deformed rotor that we have explored in this paper, one should look at the entire band structure and not just the ground-state rotational band.
As we have already demonstrated, the $`SU(3)`$ symmetry that emerges from the PSM for <sup>168</sup>Er corresponds to effective $`\{[\lambda _\nu =32,0][\lambda _\pi =16,0]\}(\lambda ,\mu )I`$ irreps for neutrons and protons, respectively. Among the $`SU(3)`$ models being employed in nuclear structure physics, the FDSM (or equivalently the IBM that results if Pauli effects are neglected in the FDSM) can naturally accommodate such a representation if the shell degeneracy $`\mathrm{\Omega }`$ is large enough . As can be easily checked, no other currently existing $`SU(3)`$ fermion model permits naturally such an $`SU(3)`$ representation. For example, the leading $`SU(3)`$ irreps for <sup>168</sup>Er given by the pseudo-$`SU(3)`$ model would be $`\{[\lambda _\nu ^{\text{pseudo}}=20,\mu _\nu ^{\text{pseudo}}=4][\lambda _\pi ^{\text{pseudo}}=10,\mu _\pi ^{\text{pseudo}}=4]\}(\lambda ,\mu )I`$, which will give a band structure very different from what we have seen in Fig. 1. Thus, we conclude that in a model such as the pseudo-$`SU(3)`$, the $`SU(3)`$ spectrum found in the numerical PSM calculation cannot be reproduced by any simple symmetry limit and could emerge (if at all) only from a large mixing of irreps through numerical diagonalization.
However, it should be noted that although we have demonstrated that the $`SU(3)`$ symmetry that emerges from the PSM is of FDSM-type (or the IBM-type if Pauli effects are omitted), this does not necessarily mean that the simplest single-major shell FDSM can accommodate such an $`SU(3)`$ symmetry microscopically. If we assume that the $`SU(3)`$ symmetry which emerges for <sup>168</sup>Er in the present PSM calculation is the fermion dynamical $`SU(3)`$ symmetry of Ref. based on $`S`$-$`D`$ fermion pairs, the $`SU(3)`$ quantum numbers $`\lambda _\nu `$ and $`\lambda _\pi `$ have a very simple interpretation. They are the effective number of neutrons and protons in the valence space, $`\lambda _\nu =n_\nu ^{\text{eff}}`$ and $`\lambda _\pi =n_\pi ^{\text{eff}}`$, which form the coherent $`S`$-$`D`$ pairs and are responsible for the collective rotation. In the simplest implementations of the FDSM, the model space is restricted to one major shell for protons and neutrons. The effective neutron (proton) number $`n_\nu ^{\text{eff}}`$ ($`n_\pi ^{\text{eff}}`$) is just the number of neutrons (protons) in the normal-parity levels of a single major shell, which is obviously too small to satisfy the requirement of $`n_\nu ^{\text{eff}}32`$ and $`n_\pi ^{\text{eff}}16`$ for the <sup>168</sup>Er case. In order to accommodate the $`\{[\lambda _\nu =32,0][\lambda _\pi =16,0]\}(\lambda ,\mu )I`$ $`SU(3)`$ symmetry, the one-major shell FDSM has to be extended to a multi-major shell FDSM. Namely, the coherent $`S`$-$`D`$ pairs should be redefined in a multi-major shell space. This has already been discussed extensively in conjunction with the extension of the FDSM to the description of superdeformation, and much of that discussion applies to the present case.
One should note in this regard that there is a conceptual distinction between the $`SU(3)`$ symmetry that emerged from the PSM diagonalized in multiple shells and that which arises in the one-major shell FDSM. In the PSM calculations the $`SU(3)`$ symmetry arises from the explicit dynamical participation of both normal and abnormal-parity nucleons, while in the symmetry limit of the single one-major shell FDSM the $`SU(3)`$ symmetry itself arises entirely from the normal-parity nucleons and the abnormal parity orbital enters only implicitly through the Pauli effect and by renormalizing the $`SU(3)`$ parameters. Thus, the multi-major shell FDSM is also necessary to resolve this conceptual ambiguity. The reason that the one-major shell FDSM treats the contribution of the abnormal-parity component to the collective motion differently from normal parity orbitals is because in a single major shell there is only one abnormal parity level. It has been shown that a single-$`j`$ shell does not have enough quadrupole collectivity (the possibility to form a $`D`$ pair in the $`i_{\frac{13}{2}}`$ intruder level compared to that in the corresponding normal-parity levels is only $`\frac{6}{46}`$; see Ref. ). When the FDSM is extended to a three-shell valence space as in the PSM, the situation will be changed. A bunch of abnormal-parity levels, which are located just below the normal-parity levels, will open up. Abnormal-parity nucleons thus have enough collectivity to form coherent $`D`$ pairs and participate in collective motion. This idea has already been developed in the extension of the FDSM to superdeformation.
Although the basic theme of the present work is that the $`SU(3)`$ symmetry-like features emerge for deformed nuclei from a realistic Hamiltonian which has no such symmetry at a microscopic level, it is tempting to look for the microscopic basis underlying the emergent $`SU(3)`$ symmetry. The preceding arguments give us strong reason to speculate that this emergent $`SU(3)`$ symmetry is just a manifestation of the FDSM operating over a three-shell valence space.
## 7 Comments on the Scissors Mode Vibrations
A second important consequence of our study is that the rotational states that we have just described in terms of an $`SU(3)`$ symmetry can be regrouped and interpreted as phonon vibrations (even though no explicit vibrational information has been given to the calculation). A standard signature of ideal harmonic vibrational motion is the appearance of an equally spaced spectrum of energy levels with a characteristic degeneracy pattern. A similar connection between $`SU(3)`$ symmetry and vibration in rotating nuclei was pointed out by Wu et al. for $`\beta `$\- and $`\gamma `$-vibrations . In the present case, we have clearly shown that the band heads of the different $`SU(3)`$ representations $`[\lambda ,\mu ]`$ obtained by coupling the proton–neutron representations $`[\lambda _\pi ,0]`$ and $`[\lambda _\mu ,0]`$ tend to be equally spaced. We have also shown that the spectra of the states obtained by PSM diagonalization is dominated by nearly equally spaced bands (apart from a small anharmonicity). Therefore it is appropriate to call these “vibrational bands”. Physically the only vibrational mode allowed for protons and neutrons that retain their individual $`SU(3)`$ irreps $`[\lambda _\tau ,0]`$ (this algebraic restriction is related to a geometrical requirement of fixed deformation for both neutrons and protons) is the “scissors mode”. The scissors mode corresponds geometrically to an oscillation in the relative orientation angle of the quadrupole-deformed neutron and proton potentials.
Traditionally, collective vibrations appear as a consequence of small oscillations around the equilibrium in the deformed potential, as we know from the physics of $`\beta `$ and $`\gamma `$vibrations . Theories of the RPA type allow superposition of states displaced around the equilibrium, thus describing the physics of small oscillations, as required for $`\beta `$ and $`\gamma `$vibrations. However, the collective vibrations considered here originate from relative motion between proton and neutron systems at fixed deformation. As long as one does not decompose the neutron and proton degrees of freedom, the physics discussed in our work is not contained in such theories.
On the other hand, it is clear that the present results do not describe $`\beta `$\- and $`\gamma `$-spectra because the $`\beta `$\- and $`\gamma `$-vibrational modes do not belong to the present configuration space. This can be seen clearly with the aid of the preceding $`SU(3)`$ classification of the PSM states. The $`\beta `$\- and $`\gamma `$-vibrations (the quadrupole phonon or spin-2$`\mathrm{}`$ excitations), correspond to a coherent superposition of the $`[(\lambda _\nu 2\mu ,\mu )(\lambda _\pi ,0)]`$ and $`[(\lambda _\nu ,0)(\lambda _\pi 2\mu ,\mu )]`$ representations with $`\mu `$ equal to a non-zero even integer. In other words, either the neutron or the proton core (or both), must be excited. It is easy to show from Eq. (11) that the dominant $`SU(3)`$ representation for the lowest $`\beta `$\- and $`\gamma `$-bands is $`[(\lambda _\nu 4,2)(\lambda _\pi ,0)](\lambda 4,2)`$ if $`\lambda _\nu >\lambda _\pi `$. The $`\beta `$-band head energy is given by
$$E_{\beta 0}E_{\text{g.s.}}=3\lambda _\nu \chi _\nu ^{\text{eff}}+2\mathrm{}\omega _{\mathrm{}},\mathrm{}\omega _{\mathrm{}}=\frac{3}{2}\lambda _m\chi _{\nu \pi }^{\text{eff}}$$
(26)
Thus the $`\beta `$-vibrational band head energy is significantly lower than the energy $`E_{\text{2nd}0^+}=2\mathrm{}\omega _{\mathrm{}}`$ of the lowest $`K=0^+`$ band of the scissors mode excitation shown in Fig. 1 which corresponds to the $`[(\lambda _\nu ,0)(\lambda _\pi ,0)](\lambda _m4,2)`$ representation. In fact, $`E_{\beta 0}`$ will be lower than even the energy $`E_{\text{1st}1^+}`$ of the band head of the first 1<sup>+</sup> scissors mode $`E_{\text{1st}1^+}=\mathrm{}\omega _{\mathrm{}}+2A`$, if the condition $`\chi _\nu ^{\text{eff}}>\frac{\lambda _m}{2\lambda _\nu }\chi _{\nu \pi }^{\text{eff}}`$ is fulfilled. In the case of <sup>168</sup>Er discussed here, this condition becomes $`\chi _\nu ^{\text{eff}}>0.75\chi _{\nu \pi }^{\text{eff}}`$.
This means that, unlike the scissors mode, the $`\beta `$\- and $`\gamma `$-vibrations do not correspond to relative motion between fixed neutron and proton fields, but to an internal collective excitation of neutrons or/and protons. This picture is consistent with the conventional shape vibration picture . Therefore, to obtain the classical $`\beta `$\- and $`\gamma `$-bands within the present framework one must build a richer set of correlations into the vacuum. In principle, this could be achieved by mixing a large set of multi-quasiparticle states, but in practice, the configuration space may be too large to handle if there is substantial collectivity in these modes. A more efficient way to introduce these states is to build those degrees of freedom into the intrinsic basis. Very recently, $`\gamma `$-band and multi-phonon $`\gamma `$-vibrational states in rare-earth nuclei have been obtained by Sun et al. by using three-dimensional angular momentum projection on triaxially deformed potential .
## 8 The First 1<sup>+</sup>, 2<sup>+</sup>, and the Second 0<sup>+</sup> Bands
Several lowest excited bands appearing in Fig. 1 warrant further discussion. In Fig. 4 we plot these bands, together with the g-band from the PSM calculations and from experiment . There are other observed low-lying collective states in this nucleus (for example, the 2<sup>+</sup> $`\gamma `$-band starting at 0.821 MeV , which was extensively studied in the PSM framework in Ref. ). Discussion of these states is beyond the scope of the present paper, and therefore, we omit plotting them in Fig. 4.
The first 1<sup>+</sup> band corresponds to a 1-phonon excitation with the excitation energy relative to the ground state depending on the interaction strengths used in the calculation. The good agreement of the calculated g-band with data (and similar results for many other calculations in this mass region ) suggests that the strengths we use here are realistic. These excitations are due to an $`SU(3)`$ coupling $`(n_\nu ,0)(n_\pi ,0)`$ in which both neutron and proton intrinsic systems remain in the ground states; thus they are related physically to relative motion between neutrons and protons. We conclude that this may be the 1<sup>+</sup> scissors mode band previously suggested in other models . Our obtained bandhead of this 1<sup>+</sup> band is at an excitation of about 3 MeV, reasonably within the energy range of experimental observations for scissors mode .
The present results indicate that the PSM provides a microscopic framework in which collective modes that may be closely identified with those proposed in earlier geometrical and algebraic descriptions emerge as the lowest excitations. Furthermore, it is already well established that the PSM describes structures built on quasiparticle excitation very well . Therefore, the extension of the present calculations to a larger basis including 2 and possibly 4 quasiparticle excitations of the n–p coupled vacuum will provide a microscopic formalism in which collective and quasiparticle degrees of freedom enter on an equal footing. Such calculations are possible and are presently being explored. We may expect that the long-debated question of whether the observed $`1^+`$ states are collective or two-quasiparticle in nature may then be resolved through such quantitative calculations of this sort.
The other excited bands are generalizations of the scissors mode corresponding to multi-phonon excitations. In Fig. 4, the next two bands at 6.5 MeV excitation energy are nearly degenerate 2-phonon states, corresponding to the coupling of two 1$`\mathrm{}`$-phonons to total spins 0$`\mathrm{}`$ and 2$`\mathrm{}`$. These are theoretically predicted multi-phonon excitations of the scissors mode, and have not to our knowledge been seen experimentally. Although the level density is expected to be high at that excitation energy and symmetry breaking will fractionate the strength, these predicted states might be detectable in the new generation of modern detectors.
## 9 Summary
We have found many new collective states in a shell model diagonalization based on separately projected neutron and proton Nilsson + BCS vacuum states. We have shown that these states exhibit an almost perfect $`SU(3)`$ symmetry, both in their spectra and their wavefunctions. We have shown also that these states can be classified systematically in a phonon spectrum with weak anharmonicity. Among these states, the lowest 1<sup>+</sup> band at about 3 MeV corresponds to the scissors mode predicted in a classical geometrical picture. The PSM is a shell model diagonalization without explicit $`SU(3)`$ symmetries. However, the quantitative agreement with an $`SU(3)`$ model provides an algebraic fermion classification scheme for the states obtained from the PSM diagonalization, and suggests that the projected BCS vacuum for a well-deformed system is very close to the $`SU(3)`$ dynamical symmetry limit of an $`S`$-$`D`$ pair fermion system. This in turn implies a good boson algebraic symmetry if Pauli effects may be ignored. Finally, we have proposed that the extension of the present calculations to include quasiparticle excitations provides a quantitative framework to determine whether “scissors mode” 1<sup>+</sup> states are more properly viewed as collective excitations or as quasiparticle states.
The present paper deals only with even-even nuclear system. However, our preliminary results have shown that similar conclusions can also be drawn for odd- and odd-odd nuclei.
Y.S. thanks Professor Gui-Lu Long of Department of Physics of Tsinghua University for warm hospitality, where the final version of this paper was completed. K.H.B. thanks Dr. S. Raman of Oak Ridge National Laboratory for partial support.
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# Analytic Theory of fractal growth patterns in 2 dimensions
Diffusion Limited Aggregation (DLA) is a model of fractal growth that was introduced in 1981 by Witten and Sander . It had since attained a paradigmatic status due to its simplicity and its underlying role for a variety of pattern forming processes including dielectric breakdown , viscous fingering, electro-chemical deposition , dendritic and snowflake growth , chemical dissolution , geological phenomena and certain biological phenomena such as bacterial growth and viscous fingering through gastric mucin . The algorithm begins with fixing one particle at the center of coordinates in $`d`$-dimensions, and follows the creation of a cluster by releasing random walkers from infinity, allowing them to walk around until they hit any particle belonging to the cluster. The fundamental difficulty of such growth processes is that their mathematical description calls for solving equations with boundary conditions on a complex, evolving interface. Despite tremendous efforts this and related difficulties defied all attempts to develop analytic theory of DLA. In fact, even the numerical estimates of the fractal dimension $`D`$ of DLA turned out to converge extremely slowly with the number of particles $`n`$ of the cluster, leading even to speculations that asymptotically the clusters were plane filling (i.e. $`D=2`$). In this Letter we offer a theory for fractal growth patterns in 2-d, including DLA as a particular case. In this theory the fractal dimension of the asymptotic cluster manifests itself as a renormalization exponent observable already at very early growth stages. Using early stage dynamics we compute $`1.6896<D<1.7135`$, and explain why traditional numerical estimates converged so slowly. The present theory is equally applicable to other fractal growth processes in 2-dimensions, and we discuss similar computations for such growth models as well.
For continuous time processes in 2-dimensions the above mentioned difficulty was efficiently dealt with in the past by considering the conformal map from the exterior of the unit circle in the complex plane to the exterior of the (simply connected) growing pattern. In this way the “interface” in the mathematical plane remains forever simple, and the complexity of the evolving interface is delegated to the dynamics of the conformal map. For discrete particle growth such a language was developed recently , showing that a large variety of fractal clusters in two dimensions can be grown by iterating conformal maps. In this Letter we employ this language to show that the fractal dimension of the asymptotically large clusters has a surprising and useful role as a renormalization exponent in a rescaling theory of these clusters in their early growth phases. This finding allows us to compute the fractal dimension to desired accuracy.
Once a fractal object is well developed, it is extremely difficult to find a conformal map from a smooth region to its boundary, simply because the conformal map is terribly singular on the tips of a fractal shape. The derivative of the inverse map is the growth probability for a random walker to hit the interface (known as the “harmonic measure”) which has been shown to be a multifractal measure characterized by infinitely many exponents . Accordingly, in the present approach one grows the cluster by iteratively constructing the conformal map starting from a smooth initial interface. Consider $`\mathrm{\Phi }^{(n)}(w)`$ which conformally maps the exterior of the unit circle $`e^{i\theta }`$ in the mathematical $`w`$–plane onto the complement of the (simply-connected) cluster of $`n`$ particles in the physical $`z`$–plane . The unit circle is mapped to the boundary of the cluster. The map $`\mathrm{\Phi }^{(n)}(w)`$ is made from compositions of elementary maps $`\varphi _{\lambda ,\theta }`$,
$$\mathrm{\Phi }^{(n)}(w)=\mathrm{\Phi }^{(n1)}(\varphi _{\lambda _n,\theta _n}(w)),$$
(1)
where the elementary map $`\varphi _{\lambda ,\theta }`$ transforms the unit circle to a circle with a “bump” of linear size $`\sqrt{\lambda }`$ around the point $`w=e^{i\theta }`$. An example of a good elementary map $`\varphi _{\lambda ,\theta }`$ was proposed in , endowed with a parameter $`a`$ in the range $`0<a<1`$, determining the shape of the bump. We employ $`a=2/3`$ which is consistent with semicircular bumps. Accordingly the map $`\mathrm{\Phi }^{(n)}(w)`$ adds on a new bump to the image of the unit circle under $`\mathrm{\Phi }^{(n1)}(w)`$. The bumps in the $`z`$-plane simulate the accreted particles in the physical space formulation of the growth process. Since we want to have fixed size bumps in the physical space, say of fixed area $`\lambda _0`$, we choose in the $`n`$th step
$$\lambda _n=\frac{\lambda _0}{|\mathrm{\Phi }_{}^{(n1)}{}_{}{}^{}(e^{i\theta _n})|^2}.$$
(2)
The recursive dynamics can be represented as iterations of the map $`\varphi _{\lambda _n,\theta _n}(w)`$,
$$\mathrm{\Phi }^{(n)}(w)=\varphi _{\lambda _1,\theta _1}\varphi _{\lambda _2,\theta _2}\mathrm{}\varphi _{\lambda _n,\theta _n}(\omega ).$$
(3)
The difference between various growth models will manifest itself in the different itineraries $`\{\theta _1\mathrm{}\theta _n\}`$. To grow a DLA we have to choose random positions $`\theta _n`$. This way we accrete fixed size bumps in the physical plane according to the harmonic measure (which is transformed into a uniform measure by the analytic inverse of $`\mathrm{\Phi }^{(n)}`$). The DLA cluster is fully determined by the stochastic itinerary $`\{\theta _k\}_{k=1}^n`$. In Fig. 1 we present a typical DLA cluster grown by this method to size $`n=`$100 000.
Other fractal clusters can be obtained by choosing a non-random itinerary . A beautiful family of growth patterns is obtained from quasi-periodic itineraries,
$$\theta _{k+1}=\theta _k+2\pi W,$$
(4)
where $`W`$ is a quadratic irrational number. An example is shown in Fig. 2, in which $`W`$ is the golden mean $`(\sqrt{5}+1)/2`$. In it was argued that itineraries obtained from (4) using for $`W`$ other values of quadratic irrationals lead to clusters of different appearance but the same dimension, which was estimated numerically to be $`D=1.86\pm 0.03`$. In the same paper other deterministic itineraries (not obtained from circle maps) where shown to lead to clusters with different dimensions. One (trivial) example that is nevertheless useful for our consideration below is the itinerary $`\theta _k=0`$ for all $`k`$. Such an itinerary grows a 1 dimensional wire of width $`\sqrt{\lambda _0}`$.
The great advantage of the availability of a conformal map is that it affords us analytic power that is not obtainable otherwise. To understand this consider the Laurent expansion of $`\mathrm{\Phi }^{(n)}(w)`$ :
$$\mathrm{\Phi }^{(n)}(w)=F_1^{(n)}w+F_0^{(n)}+F_1^{(n)}w^1+F_2^{(n)}w^2+\mathrm{}$$
(5)
The recursion equations for the Laurent coefficients of $`\mathrm{\Phi }^{(n)}(w)`$ can be obtained analytically, and in particular one shows that
$$F_1^{(n)}=\underset{k=1}{\overset{n}{}}[1+\lambda _k]^a.$$
(6)
The first Laurent coefficient $`F_1^{(n)}`$ has a distinguished role in determining the fractal dimension of the cluster, being identical to the Laplace radius which is the radius of a charged disk having the same field far away as the charged cluster . Moreover, defining $`R_n`$ as the minimal radius of all circles in $`z`$ that contain the $`n`$-cluster, one can prove that
$$R_n4F_1^{(n)}.$$
(7)
Accordingly one expects that for sufficiently large clusters (to be made precise below)
$$F_1^{(n)}n^{1/D}\sqrt{\lambda _0},n\mathrm{},$$
(8)
as $`\sqrt{\lambda _0}`$ remains the only length scale in the problem when the radius of the cluster is much larger than the radius of the initial smooth interface (which we take as the unit circle in this discussion, $`\mathrm{\Phi }^{(0)}(\omega )=\omega )`$).
These observations lead now to the central development of this Letter, and to the most important result. Consider a renormalization process in which we fix the initial smooth interface, but change $`\lambda _0`$, and then rescale $`n`$ such as to get the “same” cluster. Of course we need to specify what do we mean by the “same” cluster, and a natural requirement is that the electrostatic field on coarse scales (i.e. far from the cluster) will remain invariant. In other words, we should require the invariance of the Laplace radius $`F_1^{(n)}`$ (and possibly of additional low order Laurent coefficients) under renormalization. Clearly, for a given itinerary $`\{\theta _k\}_{k=1}^n`$, $`F_1^{(n)}`$ is a function of $`n`$ and $`\lambda _0`$ only. Accordingly, considering Eq.(8), we note that such a renormalization process can reach a fixed point if and only if $`F_1^{(n)}(\lambda _0)`$ attains a nontrivial fixed point function $`F_1^{}`$ of the single “scaling” variable $`x=\sqrt{\lambda _0}n^{1/D}`$. Obviously in the asymptotic limit $`x1`$ $`F_1^{(n)}(\lambda _0)`$ must converge to $`F_1^{}`$ which is linear in $`x`$ in this regime. The main new findings of this Letter are that $`F_1^{}`$ exists as a nonlinear function of $`x`$, and that $`F_1^{(n)}(\lambda _0)`$ converges (within every universality class) to its fixed point function $`F_1^{}`$ already for $`x1`$.
In principle one can demonstrate the convergence to $`F_1^{}`$ analytically. This is easy to do in the case of the degenerate itinerary growing a wire. In this case we can demonstrate convergence after the addition of 2-3 bumps, even in the limit $`\lambda _00`$, see Fig. 3 . For $`x0`$ $`F_1^{}(x)1x^2`$ in this case.
For other nontrivial itineraries it becomes increasingly cumbersome to demonstrate the convergence by hand. With the assistance of the machine we can demonstrate the convergence in all the other cases. In Fig. 4 we present $`F_1^{(n)}(\lambda _0)1`$ as a function of $`x`$ for a typical DLA itinerary and for values of $`\lambda _0`$ ranging between $`10^8`$ to $`10^3`$. We note that for $`\lambda _00`$ the convergence to the fixed point function is obtained infinitesimally close to the initial circle for which $`F_1^{(n=0)}=1`$. In fact, data collapse (with $`D`$ chosen right) for this itinerary, as well as for all other nontrivial itineraries, is obtained for $`nn_c`$ where $`n_c2\pi /\sqrt{\lambda _0}`$. This is the number of bumps required to obtain one-layer coverage of the original circular interface. Obviously $`n_c^{1/D}\sqrt{\lambda _0}0`$ for $`\lambda _00`$, demonstrating the convergence to $`F_1^{}`$ for $`x1`$. In Fig.5 we exhibit the convergence for the Golden Mean itinerary. Note that the fixed point functions are different, and they both differ from the wire case. The main point of this analysis is that convergence to $`F_1^{}`$ can be obtained for $`x`$ arbitrarily small by going to the limit $`\lambda _00`$.
The existence of a fixed point function translates immediately to a calculational scheme. Consider a given itinerary $`\{\theta _k\}_{k=1}^N`$ of one of the above classes, and calculate $`F_1^{(n)}(\lambda _0)`$ for $`N>n>n_c(\lambda _0)`$. Rescale now $`\lambda _0\lambda _0/s`$, and calculate $`F_1^{(n^{})}(\lambda _0/s)`$ for $`N>n^{}>n_c(\lambda _0/s)`$. We can compute $`D`$ from finding the value $`n^{}`$ which preserves the Laplace radius under rescaling of $`\lambda _0`$ by $`s`$:
$$\left(\frac{n^{}}{n}\right)^{1/D}=\sqrt{s}.$$
(9)
Since $`F_1^{}`$ is monotonic (as is immediately seen from Eq.(6)), there is only one solution to this equation,
$$D=\frac{2[\mathrm{log}n^{}\mathrm{log}n]}{\mathrm{log}s}.$$
(10)
As a first example consider the wire case. Computing $`F_1^{(10)}(10^6)`$ we find that for $`\lambda _0=10^5`$ the same value of $`F_1^{(n)}`$ is obtained for $`n`$ between 3 and 4. Eq.(10) with $`s=10`$ then predicts $`0.796<D<1.045`$. Repeating for $`F_1^{(100)}(10^6)`$ we find the same value of $`F_1^{(n)}(10^5)`$ for $`n`$ between 31 and 32. From Eq.(10) $`0.9897<D<1.0173`$. We stress that this precision is obtained when $`F_1^{(n)}1`$ is still 0.0133! Lastly, using $`F_1^{(500)}(10^6)=1.18254`$ we compute $`0.9951<D<1.0006`$, and any desired accuracy can be achieved by increasing $`n`$. The reader should note that the fractal dimension of the asymptotic cluster ($`D=1`$ in this case) can be extracted from growth events infinitesimally close to the initial unit circle by decreasing $`\lambda _0`$ (even without increasing $`n`$ in this specific case). Secondly we consider the deterministic itinerary (4) with Golden Mean winding number. Using values of $`F_1^{(n)}(10^6)`$ and $`F_1^{(n)}(10^5)`$ between 1.10 and 1.20 we can bound the dimension of the cluster to $`1.8305<D<1.8380`$. Note that in this case we need to have at least one layer covering which is obtained only when $`n_c2\pi /\sqrt{\lambda _0}`$. Finally, we consider DLA. Here the itineraries are stochastic and one c ould imagine that only under extensive ensemble averaging one would obtain tight bounds on $`D`$. In fact we find that using values of $`F_1^{(n)}(10^8)`$ and $`F_1^{(n)}(2\times 10^8)`$ between 1.002 and 1.01 we can bound $`D`$ as tightly as $`1.6896<D<1.7135`$. Note that to achieve this accuracy we did not need to go to high values of $`F_1^{(n)}`$, but rather used small values of $`\lambda _0`$ to reach convergence very early. This demonstrates again the unexpected fact that the asymptotic dimension appears as a renormalization exponent right after one or a few layers of particles cover the circle, and very much before $`F_1^{(n)}n^{1/D}`$.
At this point we need to address a few questions:
(i) Why classical numerical estimates of the fractal dimension of DLA converge so slowly?
In standard numerical experiments the radius of gyration of the grown cluster was plotted in log-log coordinates against the number of particles, with $`D`$ estimated from the slope. Examining our fixed point functions $`F_1^{}`$ (see Figs.3-5) we note the slow crossover to linear behaviour, which is not fully achieved even for extremely high values of $`n`$. In this respect we understand from Eq.(7) that a reliable estimate of $`D`$ from radius of gyration calculation requires inhuman effort, as was indeed experienced by workers in the field . In the present formulation the appearance of the asymptotic $`D`$ as a renormalization exponent already at early stages of the growth allows a convergent calculation.
(ii) Is a typical DLA cluster self-averaging?
It was demonstrated in that an ensemble of DLA clusters exhibits statistics of $`F_1^{(n)}`$ with standard deviation that shrinks to zero when $`n\mathrm{}`$. Nevertheless in it was argued that there may remain residual fluctuations of the dimension $`D`$ as extracted from (8). The numerics presented above is not sufficient to resolve this question, but additional highly accurate numerics in the limit $`\lambda _00,n\mathrm{}`$ should do.
(iii) Is the problem co-dimension 1?
The multi-scaling properties of the harmonic measure have left the impression that computing the fractal dimension of DLA will require a simultaneous control of the host of exponents characterizing the measure. For example the scaling relation $`D_3=D/2`$ (with $`D_3`$ beeing a generalized dimension in the Hentschel-Procaccia sense ) that was derived first by Halsey strengthened this impression. The approach presented here indicates that an appropriate fixed point structure can be obtained with only one relevant exponent, i.e. $`1/D`$, and this exponent appears in the dynamics much before the measure becomes multiscaling.
###### Acknowledgements.
This work has been supported in part by the European Commission under the TMR program and the Naftali and Anna Backenroth-Bronicki Fund for Research in Chaos and Complexity.
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# Anomalous Proximity Effect in Underdoped YBa2Cu3O6+x Josephson Junctions
## Abstract
Josephson junctions were photogenerated in underdoped thin films of the YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub> family using a near-field scanning optical microscope. The observation of the Josephson effect for separations as large as 100 nm between two wires indicates the existence of an anomalously large proximity effect and show that the underdoped insulating material in the gap of the junction is readily perturbed into the superconducting state. The critical current of the junctions was found to be consistent with the conventional Josephson relationship. This result constrains the applicability of SO(5) theory to explain the phase diagram of high critical temperature superconductors.
Despite being among the most intensely studied condensed-matter systems, high temperature superconductors (HTS) have resisted a microscopic understanding. They are strongly anisotropic highly correlated electronic systems, showing anomalous characteristics in both the superconducting and non-superconducting phases. In particular, the nature of the transition between the low carrier concentration insulating antiferromagnetic (AF) phase and the high carrier concentration metallic and superconducting (SC) phase is not known and it is believed to be key to uncovering the superconducting mechanism in these materials. The quest to gain a better understanding of these issues is reflected in extensive experimental and theoretical work. Recently an elegant theory, based on SO(5) group symmetry, proposed a basic framework to explain the HTS phase diagram. A five component superspin was introduced with two of its components associated with the order parameter in a d-wave SC state and the other three identified with the order parameter of the AF phase. In this theory, the quantum phase transition between the AF and SC phases corresponds to a change in the orientation of the superspin in this five-dimensional space. A different approach for describing the rich phase diagram of the HTS materials postulates superconducting pairing at a temperature $`T^{}`$ well above the superconducting critical temperature $`T_c`$. The low stiffness of the superconducting order parameter, due to their large penetration length, leads to fluctuations in the phase of the order parameter between $`T^{}`$ and $`T_c`$. It is not until phase coherence is achieved at $`T_c`$ that superconductivity is established in the material. The overall behavior of the superconductor between $`T^{}`$ and $`T_c`$ resembles that of a Kosterlitz-Thouless transition in conventional two-dimensional superconductors.
The electronic nature of the underdoped system near the superconducting state is significantly different than the normal state found in conventional superconductors. Consequently, the experimental manifestation of superconductivity may also be expected to differ. Within the framework of the SO(5) theory, Demler et al. have predicted that the current-phase relationship in the coupling between two HTS separated by a thin AF layer (a SAS junction) is modified from the Josephson relation $`I_J=I_osin\varphi `$, with $`\varphi `$ the superconducting phase difference across the junction. If the thickness of the junction $`d<d_c=\pi \xi _A`$ the AF material becomes a superconductor for $`\varphi <\varphi _c`$ and “conventional” Josephson effect occurs only for $`\varphi >\varphi _c`$, where $`\xi _A`$ represents a new superconducting correlation length in the AF phase and $`\varphi _c\pi \sqrt{1(\frac{d}{d_c})^2}`$.
This difference in the current-phase relationship of a SAS junction leads to different behavior of Josephson junctions, as investigated by den Hertog and co-workers. In particular, they predicted that the critical current of the junctions $`I_c(H)`$ would show a linear decrease with an applied magnetic field around $`H=0`$; i. e., a cusp in the $`I_c(H)`$ dependence, instead of the quadratic low field behavior found in conventional junctions. We are not aware of an analogous prediction within the framework of the fluctuating-phase model.
Both the fluctuating-phase approach and the SO(5) theory share the possibility of a large proximity effect when the material separating the HTS wires is the insulating HTS precursor. These considerations suggest that Josephson junctions may be made by separating HTS with a relatively thick layer of the precursor material.
We exploit the capability of locally photodoping an insulating RBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub> material (with R a rare earth), to induce superconducting wires separated by a non-superconducting region. The flexibility provided by our near-field scanning optical microscope (NSOM) allows us to vary the gap between the $`w150`$ nm wide superconducting wires.
The samples under consideration are c-axis oriented thin films. One is a 180 nm thick GdBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub> film grown on (100) MgO substrates by dc-magnetron sputtering. The other sample is a 120 nm thick YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub> (YBCO) film deposited on a SrTiO<sub>3</sub> substrate by laser ablation. The as-grown films show good physical properties with linear temperature dependence of the dc resistivity. Their critical temperatures, determined by ac-susceptibility, were $`T_c=89.4`$ K and $`89.2`$ K, respectively. To place the samples in the insulating side of the Metal-Insulator transition (MI), their oxygen content was adjusted to $`x0.4`$. After reduction, the resistivity of the samples at 4 K was found to be $`\rho 6`$ m$`\mathrm{\Omega }`$cm. The results obtained in both samples are very similar and we will concentrate on the data obtained in the YBCO film.
The film was mounted on the insert of a continuous-flow He cryostat and photodoped with an Al-coated 60 nm aperture NSOM probe. Light from either a 3 mW He-Ne or a 1 mW, $`\lambda `$ = 1.55 $`\mu `$m, InGaAsP laser was coupled into the inputs of a 50/50 2 $`\times `$ 2 optical fiber coupler. One of the outputs of the coupler was connected to the NSOM probe, while the other was used to monitor the laser stability. Photogeneration was accomplished by illuminating with the 1.96 eV light from the He-Ne laser, which is close to the maximum of the photodoping efficiency. The photoinduced changes in the sample were detected by imaging the reflectance variations at $`\lambda `$ = 1.55 $`\mu `$m with the InGaAsP laser. The reflected light was collected in the far field with conventional optics. The InGaAsP laser was chosen because $`\lambda `$ = 1.55 $`\mu `$m radiation provides the maximum change in reflectivity when crossing the MI transition and it does not induce any further photogeneration. Photon fluxes per unit time were estimated to be $`Q_{exc}8.6\times 10^{20}`$ photons/(cm<sup>2</sup> s) and $`Q_{ref}3.2\times 10^{20}`$ photons/(cm<sup>2</sup>s) for photoexcitation and reflectivity measurements, respectively. All the measurements involving the NSOM were performed at room temperature. When necessary to prevent the superconducting wires from decaying by e-h recombination, the NSOM head was removed, the cryostat closed and pumped to 10<sup>-6</sup> Torr and the sample cooled to 200 K in less than 15 min. It is well documented that below 250 K the photoinduced state is metastable.
Typical NSOM reflectance scans are shown in Fig. 1. The wires were defined and the scans obtained as described in Ref. . Fig. 1 shows that the reflectance and $`T_c`$ of the wires increase with the duration of the photogeneration. These results are explained by the photoinduced local increase of free holes in the CuO planes of YBCO.
Josephson junctions were defined by photogenerating a wire and leaving an unilluminated gap along its length, as described in Ref. . A typical example of these junctions is illustrated in Fig. 2a. We determine the gap $`d`$ between the superconducting wires from the reflectance data shown in Fig. 1. We define $`d`$ as the range where the reflectance is lower than that corresponding to the wire in Fig. 1b, which has a $`T_c`$ 4 K. From this definition of $`d`$ the part of the junction between points b of Fig. 2b is insulating in character. The separation between these points is $`90`$ nm. Because of the finite resolution of the NSOM this procedure gives a lower limit of the thickness of the barrier between the wires. After defining an identical junction and cooling the system to 4 K, $`IV`$ characteristics were obtained, as shown in the inset of Fig. 2a. The zero dissipation region shows the existence of Josephson effect between the two wires. The rounding of the I-V curves is understood in terms of thermal fluctuations in the Josephson junction.
The observation of the Josephson effect for a separation between superconducting wires much greater than the coherence length in the superconducting state ($`\xi _o1`$ nm) is one of the main results of this paper. This “colossal” proximity effect is inexplicable even if it is assumed that the material in the gap is metallic. In this case, the coherence length for the metal in the clean limit is $`\zeta \frac{\mathrm{}v_f}{2\pi k_BT}`$, where $`v_f`$ is its Fermi velocity. For reasonable values of $`v_F`$, $`\zeta `$ turns out to be the same order of magnitude as $`\xi _o`$. From these considerations we conclude that conventional proximity effect cannot explain our results and that the insulating material in the junctions exhibits an anomalously large proximity effect. Further evidence that a superconducting state is induced in the gap material is provided by the value of the critical current of the junction. For the $`d=45`$ nm junction, $`I_c=2.6\mu `$A, very close to the measured value of $`I_c=11.5\mu `$A for the wire of Fig. 1a.
As can be seen in the inset of Fig. 2b, $`I_cR_N`$ ($`R_N=\frac{dV}{dI}|_{I=7\mu \mathrm{A}}`$ is the shunt resistance of the junction) is nearly independent of $`d`$ for $`d<110`$ nm. This fact has implications for both SO(5) theory and the fluctuating-phase models: When the existence of free vortices is used to explain the lack of phase coherence in the gap material, $`I_cR_N\mathrm{exp}(d/2\zeta _g)`$ is expected, where $`\zeta _g\xi _o\mathrm{exp}(T_\mathrm{\Theta }/T)`$ is the correlation length for phase fluctuations, and $`T_\mathrm{\Theta }`$ is the phase stiffness expressed in temperature units. From the observed weak $`d`$ dependence of $`I_cR_N`$ when $`d\stackrel{<}{}`$ 100 nm, it follows that $`\zeta _g90`$ nm and $`T_\mathrm{\Theta }\stackrel{>}{}30`$ K. As the separation between wires increases, however, a collapse in the phase coherence is observed and no Josephson effect is observed. In the SO(5) theory $`I_c\mathrm{exp}(d/\xi _A)`$ implying that the correlation length satisfies $`\xi _Ad`$. It is possible to have strong $`d`$ dependence in both $`I_c`$ and $`R_N`$, that cancel each other, as in the case in a tunnel Josephson junction in conventional superconductors. This possibility is ruled out, however, by the measured temperature dependence of the resistance of the non-superconducting material. $`R_N`$ is temperature independent if the Josephson effect is observed. Once the separation between the wires is large enough that the Josephson effect is absent, an insulating-like response is observed for the unphotodoped material. The difference found in the temperature dependence of the non-superconducting material for different junctions is attributable to their difference in length, as shown in Fig. 3. The minor deviations between the data sets might be associated with the specific geometric details of the path followed by the transport current in each case, together with the definition of $`d`$.
The observed behavior of these Josephson junctions may be expected if the gap material is thermodynamically very close to the superconducting state and the presence of the superconducting leads quenches the superconducting phase fluctuations in the gap material. The induction of superconductivity in the non-superconducting material is a feature of both the SO(5) theory and the fluctuating phase models. More generally, as discussed in Ref. , the only ingredient necessary to obtain such a large correlation length is the close proximity to a second order quantum phase transition.
One observation of this experiment appears to be in contradiction with the SO(5) theory. Since the observed behavior of $`I_cR_N`$ implies that $`\xi _Ad`$ our experiments correspond to the case analyzed in Ref. , if the material between the superconducting wires is in the AF phase. In this scenario, $`I_c(H)`$ should show a cusp for $`H0`$. The cusp is expected even for situations where the critical current is not uniform along the width of the junction, as expected for junctions made by our technique. The results obtained for small magnetic fields are shown in Fig. 4. The figure shows that for $`H0`$, $`I_c(H)\propto ̸|H|`$, but is better described by a quadratic dependence. Also shown in the figure is the calculated $`I_c(H)`$ using the model from Ref. . The experimental value of $`H=135`$ Oe for the first flux-quantum trapped in the junction, obtained from the Fraunhoffer-like dependence of $`I_c(H)`$, was used in the calculation. As shown in the figure, the agreement between the experimental data and the calculation is reasonable only if $`dd_c`$, in clear contradiction with the $`\xi _Ad`$ conclusion obtained from Fig. 2b.
The effects observed in the junctions are independent of the sequence in which the magnetic field is applied. The same results are observed if the sample is cooled through the superconducting transition in zero magnetic field, and then the field is brought up to the desired value (zero-field cooled experiment); or if the sample is cooled in the desired field (field-cooled experiment). Although the lower critical field in these underdoped HTS materials is expected to be very small, no difference is found since $`w<\lambda _p500`$ nm, where $`\lambda _p`$ is the penetration length in the superconductor. In this condition size effects are very important and it is energetically unfavorable to induce vortices in the superconducting wires.
The results described in this report show that the local combination of superconducting and insulating materials, obtained by photodoping insulating YBCO with a NSOM probe, provides an ideal opportunity to examine the validity of the different models of high critical temperature superconductors. The observation of the Josephson effect for separations $`d`$ 100 nm cannot be explained by the conventional proximity effect. We conclude that the material between the superconducting wires, although insulating, must be very close to a superconducting phase transition since superconductivity with a large phase coherence length can be induced in it. We have also been able to restrict the margin of applicability of the SO(5) theory to explain the phase diagram of HTS. In view of the observed results and with the available models, only the existence of an AF and SC phases separated by a quantum disordered phase remains compatible with the theory. Although our results seem to be in qualitative agreement with fluctuating-phase models, more theoretical work is required to obtain the characteristics of the Josephson effect within this picture.
We would like to thank C. Lobb, M. P. A. Fisher, and D. Prober for useful discussions. We are also indebted to A. J. Millis, S. A. Kivelson, and E. Demler for a critical reading of the manuscript. This work was partially supported by the National Science Foundation’s MRSEC program through the University of Maryland/NSF grant DMR-963252. Work at the Centro Atómico Bariloche and Instituto Balseiro was suported by grant ANPCYT PICT97-03-00061-01117 and grants by Consejo Nacional de Investigaciones Científicas y Tecnológicas (CONICET), Fundación Antorchas, and Fundación Balseiro of Argentina. J. G. and E. O. are members of CONICET, Argentina.
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# Orlicz-Pettis polynomials on Banach spaces
## 1. Introduction
In the study of the isomorphic properties of Banach spaces, some classes of (bounded linear) operators have been introduced which include the isomorphisms and preserve certain properties of the spaces. These are the semigroups of operators, such as the semi-Fredholm operators, associated to the ideal of compact operators , the tauberian operators, associated to the weakly compact operators , the Orlicz-Pettis operators, related to the unconditionally converging operators , etc. The semigroups and the operator ideals are somehow opposite notions: for every Banach space $`X`$, the identity map $`I_X`$ belongs to all the semigroups, while the null operator belongs to all the ideals .
These semigroups do not have an exact analogue within the class of polynomials between Banach spaces. Nevertheless, in the present paper we introduce the class of Orlicz-Pettis polynomials, related to the unconditionally converging polynomials, and show that they share certain properties with the Orlicz-Pettis (linear) operators (see Section 2), and do not satisfy some others (Section 3). In order to obtain the results of Section 3, we are led to give a number of counterexamples, mainly of vector valued polynomials on $`c_0`$, which is the key space when we deal with weakly unconditionally Cauchy series. These counterexamples are of independent interest and can give new insight into the differences between linear operators and polynomials.
Throughout the paper, $`X`$, $`Y`$ and $`Z`$ denote Banach spaces, $`X^{}`$ is the dual of $`X`$, $`B_X`$ is its closed unit ball, $`(X,Y)`$ stands for the space of operators from $`X`$ into $`Y`$, $`𝒫(^kX,Y)`$ represents the space of all $`k`$-homogeneous (continuous) polynomials from $`X`$ into $`Y`$, $`(^kX,Y)`$ is the space of all $`k`$-linear (continuous) mappings from $`X^k`$ into $`Y`$. When the range space $`Y`$ is omitted, it is supposed to be the scalar field (real or complex). We denote by $`X\widehat{}_\pi Y`$ the projective tensor product of $`X`$ and $`Y`$; the product of $`k`$ spaces is represented by $`\widehat{}_\pi ^kX:=X\widehat{}_\pi \mathrm{}\widehat{}_\pi X`$. We use the notation $`x^{(k)}:=x\stackrel{(k)}{\mathrm{}}x`$, where $`xX`$. The set of natural numbers is denoted by $``$, and $`(e_n)`$ is the unit vector basis of the space $`c_0`$. The coordinates of a vector $`xc_0`$ are denoted by $`x(i)`$ $`(i=1,2,\mathrm{})`$.
A formal series $`x_n`$ in $`X`$ is weakly unconditionally Cauchy (w.u.C., for short) if, for every $`\varphi X^{}`$, we have $`|\varphi (x_n)|<+\mathrm{}`$. Equivalent definitions may be seen in \[2, Theorem V.6\]. The series is unconditionally convergent (u.c., for short) if every subseries converges. Equivalent definitions may be seen in \[3, Theorem 1.9\].
It may be helpful to recall that every polynomial between Banach spaces takes w.u.C. (resp. u.c.) series into w.u.C. (resp. u.c.) series \[7, Theorem 2\]. The following simple fact will also be useful:
###### Proposition 1.1.
Given a polynomial $`P𝒫(^kX,Y)`$ and a w.u.C. series $`x_n`$ in $`X`$, the sequence $`\left(P\left(_{i=1}^nx_i\right)\right)_n`$ is weak Cauchy.
Proof. Let $`j(c_0,X)`$ be the operator given by $`j(e_n)=x_n`$. Then $`Pj𝒫(^kc_0,Y)`$. Since $`c_0`$ has the Dunford-Pettis property, $`Pj`$ takes weak Cauchy sequences into weak Cauchy sequences . So, the sequence
$$\left(P\left(\underset{i=1}{\overset{n}{}}x_i\right)\right)_n=\left(Pj\left(\underset{i=1}{\overset{n}{}}e_i\right)\right)_n$$
is weak Cauchy. $`\mathrm{}`$
A polynomial $`P𝒫(^kX,Y)`$ $`(k1)`$ is unconditionally converging if, for each w.u.C. series $`x_n`$ in $`X`$, the sequence $`\left(P\left(_{i=1}^nx_i\right)\right)_n`$ is convergent in $`Y`$. The space of all unconditionally converging polynomials is denoted by $`𝒫_{uc}(^kX,Y)`$ (or $`_{uc}(X,Y)`$ if $`k=1`$). This class of polynomials has been very useful for obtaining polynomial characterizations of Banach space properties (see ). Easily, $`T_{uc}(X,Y)`$ if and only if for each w.u.C. series $`x_n`$ in $`X`$, the series $`T(x_n)`$ is u.c. in $`Y`$. The polynomial $`P𝒫(^kX,Y)`$ is (weakly) compact if $`P(B_X)`$ is relatively (weakly) compact in $`Y`$. The space of compact polynomials from $`X`$ into $`Y`$ is denoted by $`𝒫_{co}(^kX,Y)`$. Every weakly compact polynomial is unconditionally converging, and every unconditionally converging polynomial on $`c_0`$ is compact (see or ).
The standard notations and definitions in Banach space theory may be seen in . For the basics in the theory of polynomials, we refer to .
## 2. Positive results
In this Section, we introduce the Orlicz-Pettis polynomials, as those satisfying the following main result. We give some other properties, and a first example of a polynomial in this class.
###### Theorem 2.1.
Given $`k`$ and $`P𝒫(^kX,Y)`$, the following assertions are equivalent:
(A) Given a w.u.C. series $`x_n`$ in $`X`$, if the set $`\left\{P\left(_{i=1}^{\mathrm{}}a_n(i)x_i\right)\right\}_n`$ is relatively weakly compact for every bounded sequence $`(a_n)c_0`$, then $`x_n`$ is u.c.
(B) Given a w.u.C. series $`x_n`$ in $`X`$, if the set $`\left\{P\left(_{i=1}^{\mathrm{}}a_n(i)x_i\right)\right\}_n`$ is relatively compact for every bounded sequence $`(a_n)c_0`$, then $`x_n`$ is u.c.
(C) If the sequence $`(x_n)X`$ is equivalent to the $`c_0`$-basis, then there is a bounded sequence $`(a_n)c_0`$ such that the set $`\left\{P\left(_{i=1}^{\mathrm{}}a_n(i)x_i\right)\right\}_n`$ is not relatively weakly compact.
(D) If the sequence $`(x_n)X`$ is equivalent to the $`c_0`$-basis, then there is a bounded sequence $`(a_n)c_0`$ such that the set $`\left\{P\left(_{i=1}^{\mathrm{}}a_n(i)x_i\right)\right\}_n`$ is not relatively compact.
(E) For every operator $`T(Z,X)`$, if $`PT𝒫_{uc}(^kZ,Y)`$, then $`T`$ is unconditionally converging.
(F) For every operator $`T(c_0,X)`$, if $`PT𝒫_{co}(^kc_0,Y)`$, then $`T`$ is compact.
(G) For every subspace $`MX`$ containing a copy of $`c_0`$, the polynomial $`Pj_M`$ is not unconditionally converging, where $`j_M`$ denotes the embedding of $`M`$ into $`X`$.
(H) For every subspace $`MX`$ isomorphic to $`c_0`$, the polynomial $`Pj_M`$ is not compact, where $`j_M`$ denotes the embedding of $`M`$ into $`X`$.
Proof. (A) $``$ (B) and (C) $``$ (D) are obvious.
(A) $``$ (C): If $`(x_n)`$ is equivalent to the $`c_0`$-basis, then the series $`x_n`$ is w.u.C., not u.c. So it is enough to apply (A).
(B) $``$ (D): By the same argument.
(D) $``$ (E): Assume $`T(Z,X)`$ is not unconditionally converging. Then we can find a sequence $`(z_n)Z`$ such that $`(z_n)`$ and $`(T(z_n))`$ are equivalent to the $`c_0`$-basis. By (D), there is a bounded sequence $`(a_n)c_0`$ such that the set $`\left\{PT\left(_{i=1}^{\mathrm{}}a_n(i)z_i\right)\right\}_n`$ is not relatively compact. Hence, letting $`M`$ be the closed linear span of $`\left\{z_n\right\}`$, there is a bounded sequence $`(x_n)M`$ such that $`\left\{PT(x_n)\right\}`$ is not relatively compact. Therefore, $`PTj_M`$ is not compact. Since $`M`$ is isomorphic to $`c_0`$, this implies that $`PTj_M`$ is not unconditionally converging, and we conclude that $`PT`$ is not unconditionally converging.
(E) $``$ (F): Take a noncompact operator $`T(c_0,X)`$. Then $`T`$ is not unconditionally converging. By (E), $`PT`$ is not unconditionally converging.
(F) $``$ (G): Suppose there is a subspace $`MX`$ containing $`c_0`$ such that $`Pj_M`$ is unconditionally converging. Then, there is a subspace $`NM`$ isomorphic to $`c_0`$ so that $`Pj_N`$ is unconditionally converging and so compact. However, $`j_N`$ is not compact.
(G) $``$ (H): This is clear, since every unconditionally converging polynomial on $`c_0`$ is compact.
(H) $``$ (A): Assume there is a w.u.C. series $`x_n`$ in $`X`$, not u.c., such that the set $`\left\{P\left(_{i=1}^{\mathrm{}}a_n(i)x_i\right)\right\}_n`$ is relatively weakly compact for every bounded sequence $`(a_n)c_0`$. Taking blocks, we can assume that $`(x_n)`$ is equivalent to the $`c_0`$-basis. Let $`M`$ be the closed linear span of $`\left\{x_n\right\}`$. Then, $`Pj_M`$ takes bounded sequences into relatively weakly compact sequences, and so $`Pj_M`$ is compact. $`\mathrm{}`$
###### Definition 2.2.
We say that $`P𝒫(^kX,Y)`$ is an Orlicz-Pettis polynomial if it satisfies the equivalent assertions of Theorem 2.1. We denote by $`𝒫_{\mathrm{uc}+}(^kX,Y)`$ the space of all $`k`$-homogeneous Orlicz-Pettis polynomials from $`X`$ into $`Y`$.
The choice of the name is due to the relationship with the u.c. series, which were studied by Orlicz and Pettis \[2, Chapter IV\].
The classes $`𝒫_{uc}(^kX,Y)`$ and $`𝒫_{\mathrm{uc}+}(^kX,Y)`$ may be described by means of a family of sets. We say that a subset $`AX`$ is a WUC-set if there is an operator $`T(c_0,X)`$ such that $`A=T\left(B_{c_0}\right)`$. We need a previous lemma.
###### Lemma 2.3.
Given a polynomial $`P𝒫(^kX,Y)`$ which is not unconditionally converging, there is an embedding $`j:c_0X`$ such that $`Pj𝒫_{uc}(^kc_0,Y)`$.
Proof. If $`P𝒫_{uc}(^kX,Y)`$, we can find a w.u.C. series $`x_n`$ in $`X`$ such that the sequence $`\left(P(x_1+\mathrm{}+x_n)\right)_n`$ is not convergent. Let $`j:c_0X`$ be given by $`j(e_n)=x_n`$. Then the sequence
$$\left(Pj(e_1+\mathrm{}+e_n)\right)_n=\left(P(x_1+\mathrm{}+x_n)\right)_n$$
is not convergent, so $`Pj𝒫_{uc}(^kc_0,Y)`$. $`\mathrm{}`$
The next Proposition highlights the opposition between the classes $`𝒫_{uc}(^kX,Y)`$ and $`𝒫_{\mathrm{uc}+}(^kX,Y)`$.
###### Proposition 2.4.
For a polynomial $`P𝒫(^kX,Y)`$, we have:
(a) $`P𝒫_{uc}(^kX,Y)`$ if and only if $`P`$ takes WUC-sets into relatively (weakly) compact sets;
(b) $`P𝒫_{\mathrm{uc}+}(^kX,Y)`$ if and only if, for every WUC-set $`AX`$, if $`P(A)`$ is relatively (weakly) compact, then so is $`A`$.
Proof. (a) Let $`P𝒫_{uc}(^kX,Y)`$ and $`A=T\left(B_{c_0}\right)`$. Then $`PT𝒫_{uc}(^kc_0,Y)`$, and so $`P(A)=PT\left(B_{c_0}\right)`$ is relatively compact. Conversely, suppose $`P𝒫_{uc}(^kX,Y)`$. By the Lemma, there is $`T(c_0,X)`$ such that $`PT𝒫_{uc}(^kc_0,Y)`$. Therefore, $`T\left(B_{c_0}\right)`$ is a WUC-set so that $`PT\left(B_{c_0}\right)`$ is not relatively weakly compact.
(b) Let $`P𝒫_{\mathrm{uc}+}(^kX,Y)`$ and choose a WUC-set $`AX`$ so that $`P(A)`$ is relatively weakly compact. Take $`T(c_0,X)`$ with $`A=T\left(B_{c_0}\right)`$. Then $`PT\left(B_{c_0}\right)`$ is relatively weakly compact, so $`PT`$ is compact. Since $`P𝒫_{\mathrm{uc}+}(^kX,Y)`$, $`T`$ is compact and $`A`$ is relatively compact. Conversely, let $`T(c_0,X)`$ with $`PT`$ compact. Since $`PT\left(B_{c_0}\right)`$ is relatively compact, we have that $`T\left(B_{c_0}\right)`$ is relatively compact, so $`T`$ is compact. Therefore, $`P𝒫_{\mathrm{uc}+}(^kX,Y)`$. $`\mathrm{}`$
The following result gives a polynomial satisfying the assertions of Theorem 2.1. Other examples are shown in Section 3.
###### Proposition 2.5.
For every Banach space $`X`$, the polynomial $`\gamma _k:X\widehat{}_\pi ^kX`$ given by
$$\gamma _k(x)=x^{(k)}$$
is an Orlicz-Pettis polynomial.
Proof. Take a subspace $`MX`$ isomorphic to $`c_0`$, and let $`j:c_0M`$ be a surjective isomorphism. From
$$\gamma _kj_Mj(e_n)=j_Mj(e_n)^k⟶̸0,$$
we get that $`\gamma _kj_M`$ is not compact. Then apply (H) of Theorem 2.1. $`\mathrm{}`$
## 3. Counterexamples
In this Section, we first give some more properties of the polynomial $`\gamma _k`$ considered in Proposition 2.5, some of which are used to establish a theorem about polynomials on spaces containing $`c_0`$. All these previous results are applied in the main theorem of the Section that provides sufficient conditions for a polynomial to be Orlicz-Pettis. A number of counterexamples are given to show that these conditions are not necessary.
Our first theorem gives a property of polynomials on spaces containing a copy of $`c_0`$. We need two previous results.
###### Lemma 3.1.
Given $`k`$, the sequence $`\left(e_n^{(k)}\right)_n`$ in $`\widehat{}_\pi ^kc_0`$ is equivalent to the unit vector basis in $`c_0`$.
Proof. By induction on $`k`$, we show that
$$a_1e_1^{(k)}+\mathrm{}+a_ne_n^{(k)}=a_1e_1+\mathrm{}+a_ne_n$$
where $`(a_i)_{i=1}^n`$ is a finite sequence of scalars.
For $`k=1`$ there is nothing to prove. Suppose the result holds for $`k1`$, and let $`r_n(t)=sign\mathrm{sin}2^n\pi t`$ for $`t[0,1]`$. Then, assuming $`\mathrm{max}|a_i|=1`$, we have
$`a_1e_1^{(k)}+\mathrm{}+a_ne_n^{(k)}=`$
$`=`$ $`{\displaystyle _0^1}\left[a_1r_1(t)e_1^{(k1)}+\mathrm{}+a_nr_n(t)e_n^{(k1)}\right]\left[a_1r_1(t)e_1+\mathrm{}+a_nr_n(t)e_n\right]𝑑t`$
$`=`$ $`2^n{\displaystyle \underset{i=1}{\overset{2^n}{}}}\left[a_1ϵ_1(i)e_1^{(k1)}+\mathrm{}+a_nϵ_n(i)e_n^{(k1)}\right]\left[a_1ϵ_1(i)e_1+\mathrm{}+a_nϵ_n(i)e_n\right]`$
where $`ϵ_j(i)`$ is the value of $`r_j(t)`$ on the interval
$$(\frac{i1}{2^n},\frac{i}{2^n})\text{for }1i2^n.$$
By the induction hypothesis, we get
$$a_1e_1^{(k)}+\mathrm{}+a_ne_n^{(k)}1.$$
On the other hand, there is $`i_0\{1,\mathrm{},k\}`$ so that
$$a_1e_1+\mathrm{}+a_ne_n=\left|a_{i_0}\right|.$$
Considering $`e_{i_0}`$ as a vector of $`\mathrm{}_1`$, take
$$e_{i_0}^{(k)}\left(\widehat{}_\pi ^kc_0\right)^{}(^kc_0).$$
Clearly,
$$e_{i_0}^{(k)}=1,\text{and}e_{i_0}^{(k)},a_1e_1^{(k)}+\mathrm{}+a_ne_n^{(k)}=a_{i_0},$$
so the result follows. $`\mathrm{}`$
###### Proposition 3.2.
There is a w.u.C. series $`y_i`$ in $`c_0`$, not u.c., such that $`_i\gamma _k(y_i)`$ is u.c. in $`\widehat{}_\pi ^kc_0`$ for each $`k2`$.
Proof. For simplicity, consider the case $`k=2`$. Take the vectors
$$y_i=\frac{e_n}{n}\text{for}n,\frac{n(n1)}{2}<i\frac{n(n+1)}{2}.$$
Clearly, the series $`y_i`$ is w.u.C. Since
$$\underset{i=1+n(n+1)/2}{\overset{m(m+1)/2}{}}y_i=1\text{for}m>n,$$
the series is not u.c. Moreover, for every finite sequence
$$\frac{n(n+1)}{2}<i_1<\mathrm{}<i_l,$$
we have, from Lemma 3.1,
$$\underset{j=1}{\overset{l}{}}y_{i_j}y_{i_j}<\frac{1}{n}.$$
Therefore, $`y_iy_i`$ is u.c. in $`c_0\widehat{}_\pi c_0`$. $`\mathrm{}`$
###### Theorem 3.3.
Given Banach spaces $`X`$ and $`Y`$, with $`X`$ containing a copy of $`c_0`$, an integer $`k>1`$ and a polynomial $`P𝒫(^kX,Y)`$, we can find a w.u.C. series $`x_i`$ in $`X`$, not u.c., such that $`P(x_i)`$ is u.c. in $`Y`$.
Proof. Let $`j:c_0X`$ be an embedding. Consider the commutative diagram
$$\begin{array}{ccc}X& \stackrel{\gamma _k}{}& \widehat{}_\pi ^kX\\ j& & ^kj& & \\ c_0& \underset{\gamma _k}{}& \widehat{}_\pi ^kc_0\end{array}$$
Let $`y_i`$ be the series constructed in the proof of Proposition 3.2. Then, the series $`j(y_i)`$ is w.u.C., not u.c., in $`X`$, and the series $`_i\gamma _kj(y_i)=_i\left(^kj\right)\gamma _k(y_i)`$ is u.c. Let $`\stackrel{~}{P}:\widehat{}_\pi ^kXY`$ be the operator defined by $`\stackrel{~}{P}(x_1\mathrm{}x_k):=\widehat{P}(x_1,\mathrm{},x_k)`$, where $`\widehat{P}(^kX,Y)`$ is the symmetric $`k`$-linear mapping associated to $`P`$. Then the series $`_iPj(y_i)=_i\stackrel{~}{P}\gamma _kj(y_i)`$ is u.c. in $`Y`$. $`\mathrm{}`$
The next theorem shows that the polynomial $`\gamma _k:X\widehat{}_\pi ^kX`$ takes sequences equivalent to the $`c_0`$-basis into sequences equivalent to the $`c_0`$-basis. Again, we need two preparatory results.
###### Lemma 3.4.
The polynomial $`\gamma _k:c_0\widehat{}_\pi ^kc_0`$ takes sequences equivalent to the $`c_0`$-basis into sequences equivalent to the $`c_0`$-basis.
Proof. Letting $`j:c_0c_0`$ be an isomorphism, consider the commutative diagram
$$\begin{array}{ccc}c_0& \stackrel{\gamma _k}{}& \widehat{}_\pi ^kc_0\\ j& & ^kj& & \\ c_0& \underset{\gamma _k}{}& \widehat{}_\pi ^kc_0\end{array}$$
Since $`\left(\gamma _k(e_n)\right)_n`$ is equivalent to the $`c_0`$-basis (Lemma 3.1), it is enough to show that $`^kj`$ is an injective isomorphism. Since $`j(c_0)`$ is complemented in $`c_0`$, there is an operator $`S:c_0c_0`$ such that $`Sj=I_{c_0}`$. Then, $`\left(^kS\right)\left(^kj\right)`$ is the identity map on $`\widehat{}_\pi ^kc_0`$. Hence, $`^kj`$ is an injective isomorphism. $`\mathrm{}`$
###### Proposition 3.5.
Let $`j:c_0X`$ be an injective isomorphism. Then the operator
$$j^{(k)}:\widehat{}_\pi ^kc_0\widehat{}_\pi ^kX$$
is an injective isomorphism.
Proof. Take $`z\widehat{}_\pi ^kc_0`$ with $`z=_{i=1}^{\mathrm{}}x_i^1\mathrm{}x_i^k`$, and $`A\left(\widehat{}_\pi ^kc_0\right)^{}(^kc_0)`$, with $`A=1`$ and $`A,z=z`$. There is an extension $`\stackrel{~}{A}(^k\mathrm{}_{\mathrm{}})`$ of $`A`$ with $`\stackrel{~}{A}=A`$ . Consider the second adjoint $`j^{}:\mathrm{}_{\mathrm{}}X^{}`$ of $`j`$. By the injectivity of $`\mathrm{}_{\mathrm{}}`$, the operator $`(j^{})^1:j^{}(\mathrm{}_{\mathrm{}})\mathrm{}_{\mathrm{}}`$ has an extension to an operator $`\pi :X^{}\mathrm{}_{\mathrm{}}`$; clearly, $`\pi j^{}=I_{\mathrm{}_{\mathrm{}}}`$. Let $`B:=\stackrel{~}{A}(\pi J_X)^k(^kX)`$, where $`J_X:XX^{}`$ is the canonical embedding. Then, $`B\stackrel{~}{A}\pi ^k`$ and
$$Bj^k=\stackrel{~}{A}(\pi J_Xj)^k=\stackrel{~}{A}\left(\pi j^{}J_{c_0}\right)^k=\stackrel{~}{A}J_{c_0}^k=A.$$
Therefore,
$`\pi ^kAj^{(k)}(z)`$ $``$ $`|B,j^{(k)}(z)|`$
$`=`$ $`\left|{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}B(j\left(x_i^1\right),\mathrm{},j\left(x_i^k\right))\right|`$
$`=`$ $`\left|{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}A(x_i^1,\mathrm{},x_i^k)\right|`$
$`=`$ $`|A,z|`$
$`=`$ $`z,`$
and this finishes the proof. $`\mathrm{}`$
###### Theorem 3.6.
The polynomial $`\gamma _k:X\widehat{}_\pi ^kX`$ takes sequences equivalent to the $`c_0`$-basis into sequences equivalent to the $`c_0`$-basis.
Proof. If $`X`$ contains no copy of $`c_0`$, the result is trivially true. If $`X=c_0`$, see Lemma 3.4. If $`X`$ contains a copy of $`c_0`$, the last Proposition reduces the problem to the case $`X=c_0`$. $`\mathrm{}`$
The following result gives an example of a polynomial on $`c_0`$ which will be useful.
###### Proposition 3.7.
There is a polynomial $`P𝒫(^2c_0,c_0)`$ such that $`P(e_n)=0`$ for all $`n`$, but $`P`$ is not compact on any infinite dimensional subspace.
Proof. Consider a bijection
$$(\alpha ,\beta ):\{(n,m)\times :nm\}.$$
Define
$$P(x):=\left(x(\alpha (i))x(\beta (i))\right)_{i=1}^{\mathrm{}}\text{for}x=(x(i))c_0.$$
Then $`P(e_n)=0`$ for all $`n`$. If $`Mc_0`$ is an infinite dimensional subspace, we can find a norm one sequence $`(x_n)c_0`$ disjointly supported such that $`dist(x_n,M)<2^n`$. For each $`n`$, let $`k_n`$ satisfy $`|x_n(k_n)|=1`$. Then
$$P(x_{2n}+x_{2n+1})|x_{2n}(k_{2n})x_{2n+1}(k_{2n+1})|=1,$$
which implies that $`P`$ is not compact on $`M`$. $`\mathrm{}`$
We can now state the main result of the Section.
###### Theorem 3.8.
Let $`P𝒫(^kX,Y)`$ be a polynomial, with $`k2`$. Consider the following assertions:
(A) If $`x_n`$ is w.u.C. in $`X`$, and $`P(x_n)`$ is u.c. in $`Y`$, then $`x_n`$ is u.c.
(B) Every sequence $`(x_n)X`$ equivalent to the $`c_0`$-basis has a subsequence $`\left(x_{n_i}\right)`$ such that $`\left(P\left(x_{n_i}\right)\right)`$ is equivalent to the $`c_0`$-basis.
(C) If $`x_n`$ is w.u.C. in $`X`$, and $`\left(P\left(_{i=1}^nx_i\right)\right)_n`$ is convergent, then $`x_n`$ is u.c.
(D) If the sequence $`(x_n)X`$ is equivalent to the $`c_0`$-basis, then $`\left(P\left(_{i=1}^nx_i\right)\right)_n`$ is not relatively compact.
(E) If the sequence $`(x_n)X`$ is equivalent to the $`c_0`$-basis, then $`limP(x_n)\to ̸0`$.
(F) $`P`$ is an Orlicz-Pettis polynomial.
Then the following and only the following implications hold:
$$\begin{array}{ccccc}\text{(A)}& & \text{(B)}& & \text{(E)}\\ & & & & \\ \text{(C)}& & \text{(D)}& & \text{(F)}\end{array}$$
Proof. (A) $``$ (B) and (A) $``$ (C): If $`P`$ satisfies (A), then Theorem 3.3 implies that $`X`$ contains no copy of $`c_0`$. So (B) and (C) are satisfied in a trivial way.
(B) $``$ (E) is obvious.
(E) $``$ (B): Let $`(x_n)X`$ be equivalent to the $`c_0`$-basis. Then $`x_n`$ is w.u.C., so $`P(x_n)`$ is also w.u.C. . In particular, $`(P(x_n))`$ is weakly null. By (E), passing to a subsequence, we can assume that $`(P(x_n))`$ is seminormalized and basic, so it is equivalent to the $`c_0`$-basis \[2, Corollary V.7\].
(B) $`\Rightarrow ̸`$ (A): Consider the polynomial $`\gamma _k:c_0\widehat{}_\pi ^kc_0`$ (see Proposition 3.2 and Lemma 3.4).
(D) $`\Rightarrow ̸`$ (A): Consider the polynomial defined in Proposition 3.7.
(C) $``$ (D): Assume $`P`$ does not satisfy (D). Then there is a sequence $`(x_n)X`$ equivalent to the $`c_0`$-basis, such that $`\left(P\left(_{i=1}^nx_i\right)\right)_n`$ is relatively compact. We can find an increasing sequence of indices $`(m_i)`$ so that $`\left(P\left(_{i=1}^ny_i\right)\right)_n`$ is convergent, where
(1)
$$y_i=\underset{j=m_i+1}{\overset{m_{i+1}}{}}x_j.$$
Since $`(y_i)`$ is equivalent to the $`c_0`$-basis, $`P`$ does not satisfy (C).
(D) $``$ (C): Assume $`P`$ does not satisfy (C). Then there is a w.u.C. series $`x_n`$ in $`X`$, not u.c., so that $`\left(P\left(_{i=1}^nx_i\right)\right)_n`$ is convergent. Take an increasing sequence of indices $`(m_i)`$ so that $`(y_i)`$ is equivalent to the $`c_0`$-basis, where $`y_i`$ is defined as in (1). Then $`\left(P\left(_{i=1}^ny_i\right)\right)_n`$ is a subsequence of $`\left(P\left(_{i=1}^nx_i\right)\right)_n`$ and so it converges, in contradiction with (D).
(E) $``$ (F): Assume $`T(Z,X)`$ is not unconditionally converging. Then we can find a sequence $`(z_n)Z`$ such that $`(z_n)`$ and $`(T(z_n))`$ are equivalent to the $`c_0`$-basis. If $`P`$ satisfies (E), we have $`PT(z_n)\to ̸0`$, which implies that $`PT𝒫_{uc}(^kZ,Y)`$. So, by Theorem 2.1(E), $`P𝒫_{\mathrm{uc}+}(^kX,Y)`$.
(F) $`\Rightarrow ̸`$ (E): The polynomial $`P`$ of Proposition 3.7 does not satisfy (E). To see that it does satisfy (F), take an operator $`T(Z,c_0)`$ not unconditionally converging. There is an operator $`j:c_0Z`$ such that $`(Tj(e_n))`$ is equivalent to $`(e_n)`$. Passing to a perturbed subsequence, we can assume that $`(Tj(e_n))`$ is disjointly supported. The series $`_n\left(j(e_{2n})+j(e_{2n+1})\right)`$ is w.u.C. However, $`PT\left(j(e_{2n})+j(e_{2n+1})\right)`$ is bounded away from $`0`$, so $`PT`$ is not unconditionally converging. By Theorem 2.1(E), $`P𝒫_{\mathrm{uc}+}(^kZ,c_0)`$.
(D) $``$ (F): Assume $`T(Z,X)`$ is not unconditionally converging. Then we can find a sequence $`(z_n)Z`$ such that $`(z_n)`$ and $`(T(z_n))`$ are equivalent to the $`c_0`$-basis. If $`P`$ satisfies (D), the sequence $`\left(PT\left(_{i=1}^nz_i\right)\right)_n`$ is not relatively compact. Hence, $`PT𝒫_{uc}(^kX,Y)`$. By Theorem 2.1(E), $`P𝒫_{\mathrm{uc}+}(^kX,Y)`$.
(D) $`\Rightarrow ̸`$ (E): Let $`P`$ be the polynomial defined in Proposition 3.7, and $`(x_n)c_0`$ a sequence equivalent to the $`c_0`$-basis. Denote $`y_n:=x_1+\mathrm{}+x_n`$, and $`z:=_{n=1}^{\mathrm{}}x_n\mathrm{}_{\mathrm{}}\backslash c_0`$. Let
$$3\delta :=\underset{i}{lim\; sup}z(i)>0.$$
If the $`lim\; sup`$ were not positive, then we would take the $`lim\; inf`$. Choose $`i_1`$ with $`|z(i_1)3\delta |<\delta `$. There is $`n_1`$ so that $`|y_n(i_1)3\delta |<\delta `$ for all $`nn_1`$. Choose now $`i_2`$ $`(i_2>i_1)`$ so that $`\left|y_{n_1}(i)\right|<\delta /2`$ for all $`ii_2`$ and $`|z(i_2)3\delta |<\delta `$. There is $`n_2`$ $`(n_2>n_1)`$ so that $`\left|y_n(i_2)3\delta \right|<\delta `$ for all $`nn_2`$.
Proceeding in this way, we obtain two increasing sequences of integers $`(i_j)`$, $`(n_j)`$ so that, for $`j<l`$,
$`P\left(y_{n_j}\right)P\left(y_{n_l}\right)`$ $``$ $`\left|y_{n_j}(i_1)y_{n_j}(i_l)y_{n_l}(i_1)y_{n_l}(i_l)\right|`$
$``$ $`\left|y_{n_l}(i_1)\right|\left|y_{n_l}(i_l)\right|\left|y_{n_j}(i_1)\right|\left|y_{n_j}(i_l)\right|`$
$`>`$ $`2\delta ^2.`$
Therefore, $`P`$ satisfies (D). Clearly, $`P`$ does not satisfy (E).
(E) $`\Rightarrow ̸`$ (D): Let $`P𝒫(^2c_0,c_0)`$ be given by
$$P(x)=\underset{j=2}{\overset{\mathrm{}}{}}\underset{i=1}{\overset{j1}{}}\left(x(j)x(i)\right)x(j)e_{j^2+i}\text{for }x=(x(i))_{i=1}^{\mathrm{}}c_0.$$
Clearly, $`P`$ satisfies (E). Since $`P(e_1+\mathrm{}+e_n)=0`$ for all $`n`$, $`P`$ does not satisfy (D). $`\mathrm{}`$
###### Remark 3.9.
In the linear case $`(k=1)`$, all the assertions of Theorem 3.8 are equivalent . So, our choice for the definition of the Orlicz-Pettis polynomials provides the widest possible class.
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# On the Law of Addition of Random Matrices
## 1 Introduction
The paper deals with the eigenvalue distribution of the sum of two $`n\times n`$ Hermitian or real symmetric random matrices as $`n\mathrm{}`$. Namely we express the limiting normalized counting measure of eigenvalues of the sum via the same measures of its two terms, assuming that latter exist and that terms are randomly rotated one with respect another by an unitary or an orthogonal random matrix uniformly distributed over the group $`U(n)`$ or $`O(n)`$ respectively.
One may mention several motivations of the problem. First, it can be regarded in the context of general problem to describe the eigenvalues of the sum of two matrices in terms of eigenvalues of two terms of the sum. The latter problem dates back at least to the paper of H. Weyl , was treated in a number of papers, including the recent paper , and related to interesting questions of combinatorics, geometry, algebra etc. (see e.g. for recent results and references). The problem is also of considerable interest for mathematical physics because of its evident links with spectral theory and quantum mechanics (perturbation theory in particular).
It is clear that one cannot expect in general a simple and closed expression for eigenvalues of the sum of two given matrices via eigenvalues of terms. Hence, it is natural to look for a “generic” asymptotic answer, studying a randomized version of the problem in which at least one of the two terms is random and both behave rather regularly as $`n\mathrm{}`$. Particular results of this type were given in where it was proved that under certain conditions the divided by $`n`$ eigenvalue counting measure of the sum converges in probability to the nonrandom limit that can be found as a unique solution of a certain functional equation. Thus, a randomized version of the problem admits a rather constructive and explicit solution. These results were developed in several directions (see e.g. \- and the recent work ). Similar problems arose recently in operator algebras studies, known now as the free (non-commutative) probability (see for results and references). In particular, the notion of the $`R`$-transform and the free convolution of measures were introduced by Voiculescu and allowed the limiting eigenvalue distributions of the sum to be given in a rather general and simple form. From the point of view of the random matrix theory the problem that we are going to consider is a version of the problem of the deformation (see e.g. for this term) of a given random matrix (that can be a non-random matrix in particular) by another random matrix in the case when ”randomness” of the latter includes as an independent part the random choice of the basis in which this matrix is diagonal. We will discuss this topic in more details in Section 2.
In this paper we present a simple method of deriving functional equations for the limiting eigenvalue distribution in a rather general situation. The method is based on certain differential identities for expectations of smooth matrix functions with respect to the normalized Haar measure of $`U(n)`$ ( or $`O(n)`$ ) and on elementary matrix identities, the resolvent identity first of all. The basic idea is the same as in : to study not the moments of the counting measure, as it was proposed in the pioneering paper by Wigner , but rather its Stieltjes (called also the Cauchy or the Borel) transform, playing the role of appropriate generating (or characteristic) function of the moments (the measure). However, the technical implementation of the idea in this paper is different and simpler then in (see Remark 1 after Theorem 2.1).
The paper is organized as follows. In Section 2 we present our main results (Theorem 2.1) and give their discussion. In Section 3 we prove Theorems 3.1 and 3.2 giving the solution of the problem under the conditions of the uniform in $`n`$ boundedness of the forth moments of the normalized counting measure of the terms. These conditions are more restrictive than those for our principle result, given in Theorem 2.1. Their advantage is that they allow us to use the main ingredients of our approach in more transparent and free of technicalities form. In Section 4 we prove Theorem 2.1, whose main condition is the uniform boundedness of the first absolute moment of the normalized counting measure of one of the two terms of the sum. In Section 5 we study certain properties of solutions of the functional equation and of the limiting counting measure. In Section 6 we discuss topics related to our main result and our technique.
## 2 Model and Main Result.
We consider the ensemble of $`n`$-dimensional Hermitian (or real symmetric) random matrices $`H_n`$ of the form
$$H_n=H_{1,n}+H_{2,n},$$
(2.1)
where
$$H_{1,n}=V_n^{}A_nV_n,H_{2,n}=U_n^{}B_nU_n.$$
We assume that $`A_n`$ and $`B_n`$ are random Hermitian (or real symmetric) matrices having arbitrary distributions, $`V_n`$ and $`U_n`$ are unitary (or orthogonal) random matrices uniformly distributed over the unitary group $`U(n)`$ (or over the orthogonal group $`O(n)`$) with respect to the Haar measure, and $`A_n`$, $`B_n`$, $`V_n`$ and $`U_n`$ are mutually independent. For the sake of definiteness we will restrict ourself to the case of Hermitian matrices and the group $`U(n)`$ respectively. The results for symmetric matrices and for the group $`O(n)`$have the same form, although their proof is more involved technically (see Section 6).
We are interested in the asymptotic behavior as $`n\mathrm{}`$ of the normalized eigenvalue counting measure (NCM) $`N_n`$ of the ensemble (2.1), defined for any Borel set $`\mathrm{\Delta }`$ by the formula
$$N_n(\lambda )=\frac{\mathrm{\#}\{\lambda _i\mathrm{\Delta }\}}{n},$$
(2.2)
where $`\lambda _i,i=1,\mathrm{},n`$ are the eigenvalues of $`H_n`$.
The problem was studied recently in the context of free (non-commutative) probability. In particular, it follows from results of that if the matrices $`A_n`$ and $`B_n`$ are non-random, their norms are uniformly bounded in $`n`$, i.e. their NCM $`N_{1,n}`$ and $`N_{2,n}`$ have uniformly in $`n`$ compact supports and if these measures have weak limits as $`n\mathrm{}`$
$$N_{1,n}N_1,N_{2,n}N_2,$$
(2.3)
then the NCM (2.2) of random matrix (2.1) converges weakly with probability 1 to a non-random measure $`N`$. Besides, if
$$f(z)=_{\mathrm{}}^{\mathrm{}}\frac{N(\mathrm{d}\lambda )}{\lambda z},\mathrm{Im}z>0,$$
(2.4)
is the Stieltjes transform of this limiting measure and
$$f_r(z)=_{\mathrm{}}^{\mathrm{}}\frac{N_r(\mathrm{d}\lambda )}{\lambda z},r=1,2,$$
(2.5)
are the Stieltjes transforms of $`N_r,`$ $`r=1,2`$ of (2.3), then according to $`f(z)`$ satisfies the functional equation
$$f(z)=f_1(z+R_2(f(z))),$$
(2.6)
where $`R_2(f)`$ is defined by the relation
$$z=\frac{1}{f_2(z)}R_2(f_2(z)))$$
(2.7)
and is known as $`R`$-transform of the measure $`N_2`$ of (2.3) (see Remark 3 after Theorem 2.1 and for the definition and properties of this transform taking into account that our definition (2.7) differs from that of by the sign). The proof of this result in was based on the asymptotic analysis of the expectations $`m_k^{(n)}`$ of moments of measure (2.2). Since, according to the spectral theorem and the definition (2.2),
$$m_k^{(n)}=𝐄\{M_k^{(n)}\},M_k^{(n)}=n^1\mathrm{Tr}H_n^k,$$
(2.8)
one can study the averaged moments $`m_k^{(n)}`$ by computing asymptotically the expectations of the divided by $`n`$ traces of the $`k`$-th powers of (2.1), i.e. of corresponding multiple sums. This direct method dates back to the classic paper by Wigner and requires a considerable amount of combinatorial analysis, existence of all moments measures $`N_{1,2}^{(n)}`$ and their rather regular behavior as $`n\mathrm{}`$ to obtain the convergence of expectations (2.8) for all integer $`k`$ and to guarantee that limiting moments determine uniquely corresponding measure. By using this method it was proved in that the expectation of $`N_n`$ converges to the limit, determined by (2.6) \- (2.7) and in that the variance $`\mathrm{Var}\{M_k^{(n)}\}=`$ $`𝐄\{(M_k^{(n)})^2\}𝐄^2\{M_k^{(n)}\}`$ admits the bound
$$\mathrm{Var}\{M_k^{(n)}\}\frac{C_k}{n^2},$$
(2.9)
where $`C_k`$ is independent of $`n`$. This bound yields evidently the convergence of all moments with probability 1, thereby the weak convergence with probability 1 of random measures (2.2) to the non-random limit, determined by (2.6) - (2.7). The convergence with probability 1 here and below is understood as that in the natural probability space
$$\mathrm{\Omega }=\underset{n}{}\mathrm{\Omega }_n,$$
(2.10)
where $`\mathrm{\Omega }_n`$ is the probability space of matrices (2.1) that is the product of respective spaces of $`A_n`$ and $`B_n`$ and two copies of the group $`U(n)`$ for $`U_n`$ and $`V_n`$.
In this paper we obtain the analogous result under weaker assumptions and by using a method, that does not involve combinatorics. This is because we work with the Stieltjes transforms of measures (2.2) and (2.3) and derive directly the functional equations for their limits and the bound analogous to (2.9) for the rate of their convergence (rather well known in the random matrix theory, see e.g. ) by using certain simple identities for expectations of matrix functions with respect to the Haar measure (Proposition 3.2 below) and elementary facts on resolvents of Hermitian matrices.
The Stieltjes transform was first used in studies of the eigenvalue distribution of random matrices in paper and proved to be an efficient tool in the field (see e.g. ). We list the properties of the Stieltjes transform that we will need below (see e.g.).
###### Proposition 2.1
Let $`m`$ be a non-negative and normalized to unity measure and
$$s(z)=\frac{m(d\lambda )}{\lambda z},\mathrm{Im}z0$$
(2.11)
be the Stieltjes transform of $`m`$ (here and below integrals without limits denote the integrals over the whole axis). Then:
1. $`s(z)`$ is analytic in $``$ and
$$|s(z)||\mathrm{Im}z|^1;$$
(2.12)
2. $$\mathrm{Im}s(z)\mathrm{Im}z>0,\mathrm{Im}z0;$$
(2.13)
3. $$\underset{y\mathrm{}}{lim}y|s(iy)|=1;$$
(2.14)
4. for any continuous function $`\phi `$ with a compact support we have the inversion (Frobenius-Perron) formula
$$\varphi (\lambda )N(\mathrm{d}\lambda )=\underset{\epsilon 0}{lim}\frac{1}{\pi }\varphi (\lambda )\mathrm{Im}s(\lambda +i\epsilon );$$
(2.15)
5. conversely, any function verifying (2.12) - (2.14) is the Stieltjes transform of a non-negative and normalized to unity measure and this one-to-one correspondence between measures and their Stieltjes transforms is continuous if one will use the topology of weak convergence for measures and the topology of convergence on compact sets of $``$ for their Stieltjes transforms.
We formulate now our main result. Since eigenvalues of a Hermitian matrix are unitary invariant we can replace matrices (2.1) by
$$H_n=A_n+U_n^{}B_nU_n,$$
(2.16)
where $`A_n`$, $`B_n`$ and $`U_n`$ are the same as in (2.1). However, it is useful to keep in mind that the problem is symmetric in $`A_n`$ and $`B_n`$. We prove
###### Theorem 2.1
Let $`H_n`$ be the random $`n\times n`$ matrix of the form (2.1). Assume that the normalized eigenvalue counting measures $`N_{r,n},r=1,2`$ of matrices $`A_n`$ and $`B_n`$ converge weakly in probability as $`n\mathrm{}`$ to the non-random nonnegative and normalized to 1 measures $`N_r,r=1,2`$ respectively and that
$$\underset{n}{sup}|\lambda |𝐄N_{r,n}^{}(\mathrm{d}\lambda )m_1<\mathrm{},$$
(2.17)
where $`N_{r,n}^{}`$ is one of the measures $`N_{1,n}`$ or $`N_{2,n}`$. Then the normalized eigenvalue counting measure $`N_n`$of $`H_n`$ converges in probability to a non-random nonnegative and normalized to 1 measure $`N`$ whose Stieltjes transform (2.4) is a unique solution of the system
$$f(z)=f_1\left(z\frac{\mathrm{\Delta }_2(z)}{f(z)}\right)$$
$$f(z)=f_2\left(z\frac{\mathrm{\Delta }_1(z)}{f(z)}\right)$$
(2.18)
$$f(z)=\frac{1\mathrm{\Delta }_1(z)\mathrm{\Delta }_2(z)}{z}$$
in the class of functions $`f(z)`$ satisfying (2.12) - (2.14) and functions $`\mathrm{\Delta }_r(z),r=1,2`$ analytic for $`\mathrm{Im}z0`$ and satisfying conditions
$$\mathrm{\Delta }_{1,2}(z)0\mathrm{as}\mathrm{Im}z\mathrm{},$$
(2.19)
where $`f_r(z),r=1,2`$ are Stieltjes transforms (2.5) of the measures $`N_r,r=1,2`$ and $`𝐄\{\}`$ denotes the expectation with respect to the probability measure, generated by $`A_n`$, $`B_n`$, $`U_n`$ and $`V_n`$
The theorem will be proved in Section 4. Here we make several remarks related to the theorem (see also Section 5).
###### Remark 1
The historically first example of a random matrix ensemble representable in the form (2.16) was proposed in and has the form
$$H_{m,n}=H_{0,n}+\underset{i=1}{\overset{m}{}}\tau _iP_{q_i},$$
(2.20)
where $`H_{0,n}`$ is a non-random $`n\times n`$ Hermitian matrix such that its normalized eigenvalue counting measure converges weakly to a limiting nonnegative and normalized to 1 measure $`N_0`$, $`\tau _i`$, $`i=1,..m`$ are i.i.d. random variables and $`P_{q_i}`$ are orthogonal projections on unit vectors $`q_i`$, $`i=1,..m`$ uniformly and independently of one another and of $`\{\tau _i\}_{i=1}^m`$ distributed over the unit sphere in $`^n`$ <sup>1</sup><sup>1</sup>1In fact, in a more general class of independent random vectors was considered, but we restrict ourself here to the unit vectors, in order to have an example of an ensemble of form (2.1).. It is clear that the matrix
$$\underset{i=1}{\overset{m}{}}\tau _iP_{q_i}$$
(2.21)
can be written in the form $`U_n^{}B_nU_n`$ of the second term of (2.1) or (2.16). According to the NCM of random matrix (2.21) converges in probability as $`n\mathrm{}`$, $`m\mathrm{}`$, $`m/nc0`$ to a non-random nonnegative and normalized to 1 measure whose Stieltjes transform $`f_{MP}(z)`$ satisfies the equation
$$f_{MP}(z)=\left(z+c\frac{\tau \sigma (\mathrm{d}\tau )}{1+\tau f_{MP}(z)}\right)^1,$$
(2.22)
where $`\sigma `$ is the probability law of $`\tau _i`$ in (2.20). Assume that $`\sigma `$ has the finite first moment
$$|\tau |\sigma (\mathrm{d}\tau )<\mathrm{}$$
(2.23)
Then taking (2.21) as the second term of (2.1) we get, in view of inequality
$$𝐄\{|\lambda |N_{2,n}(\mathrm{d}\lambda )\}n^1\underset{i=1}{\overset{m}{}}𝐄\{|\tau _i|\}=\frac{m}{n}𝐄\{|\tau |\}<\mathrm{},$$
the condition (2.17) of Theorem 2.1. Applying then Theorem 2.1 in which $`f_2(z)`$ is given by (2.22), we obtain from the two last equations of the system (2.18) that
$$\frac{\mathrm{\Delta }_1(z)}{f(z)}=c\frac{\tau \sigma (\mathrm{d}\tau )}{1+\tau f_{MP}(z)}.$$
This and the first equation of (2.18) yield the functional equation for the Stieltjes transform of the limiting eigenvalue distribution of ensemble (2.20)
$$f(z)=f_0\left(zc\frac{\tau \sigma (\mathrm{d}\tau )}{1+\tau f(z)}\right)$$
(2.24)
where $`f_0(z)`$ is the Stieltjes transform of the limiting NCM $`N_0`$ of the non-random matrix $`H_{0,n}`$. This equation was obtained in by another method, whose main ingredient was careful analysis of changes of the resolvent of matrices (2.20) induced by addition of the $`(m+1)`$-th term, i.e. by a rank-one perturbation. This allowed the authors to prove that the sequence $`g_{i,n}(z)=n^1\mathrm{Tr}(H_{i,n}z)^1,i=1,\mathrm{},m`$ converges in probability to the non-random limit $`f(z,t),z\backslash 𝐑,t[0,1]`$, as $`n\mathrm{},m\mathrm{},i\mathrm{},m/nc,i/mt`$, and that the limiting function $`f(z,t)`$ satisfies the quasilinear PDE
$$\frac{f}{t}+c\frac{\tau (t)}{1+\tau (t)f}\frac{f}{z},f(z,0)=f_0(z),$$
(2.25)
where $`\tau (t)`$ is the inverse of the probability distribution $`\sigma (\tau )=𝐏\{\tau _i\tau \}`$. It can be shown that the solution of (2.25) at $`t=1`$ coincides with (2.20) . Equation (2.25) with $`\tau (t)\mathrm{const}`$ is a particular case of the so-called complex Burgers equation appeared in the free probability , where the random matrices (2.20) provide an analytic model for the stationary processes with free increments, like in the conventional probability the heat equation and sums of i.i.d. random variables comprise an important ingredient of the theory of random processes with independent increments.
###### Remark 2
Consider the ensemble known as the deformed Gaussian ensemble :
$$H_n=H_{0,n}+M_n,$$
(2.26)
where $`H_{0,n}`$ is a non-random matrix such that its normalized eigenvalue counting measure converges weakly to the limit $`N_0`$ and $`M_n=\{M_{jk}\}_{j,k=1}^n`$ is a random Hermitian matrix whose matrix elements $`M_{jk}`$ are complex Gaussian random variables satisfying conditions:
$$\overline{M_{jk}}=M_{kj},𝐄\{M_{jk}\}=0,𝐄\{M_{j_1k_1}\overline{M_{j_2k_2}}\}=\frac{2w^2}{n}\delta _{j_1j_2}\delta _{k_1k_2}.$$
(2.27)
In other words, the ensemble is defined by the distribution
$$𝐏(\mathrm{d}M)=Z_n^1exp\left\{\frac{n}{4w^2}\mathrm{Tr}M^2\right\}\mathrm{d}M,$$
(2.28)
$$\mathrm{d}M=\underset{j=1}{\overset{n}{}}\mathrm{d}M_{jj}\underset{1j<kn}{}\mathrm{dRe}M_{jk}\mathrm{dIm}M_{jk},$$
where $`Z_n`$ is the normalization constant. The distribution defines the Gaussian Unitary Ensemble (GUE) . This is why ensemble (2.26) is called the deformed GUE . It is known that $`M_n`$ can be written in the form
$$M_n=U_n^{}\mathrm{\Lambda }_nU_n,$$
(2.29)
where $`U_n`$ are unitary matrices whose probability law is the Haar measure on $`U(n)`$ and $`\mathrm{\Lambda }_n`$ is independent of $`U_n`$ diagonal random matrix whose normalized eigenvalue counting measure converges with probability 1 to the semicircle law. The Stieltjes transform $`f_{sc}(z)`$ of the latter satisfies the simple functional equation
$$f_{sc}(z)=(z+2w^2f_{sc}(z)),$$
(2.30)
whose solution yields the semicircle law by Wigner
$$N_{sc}(d\lambda )=(4\pi w^2)^1\sqrt{8w^2\lambda ^2}\chi _{[2\sqrt{2}w,2\sqrt{2}w]}(\lambda )d\lambda ,$$
(2.31)
where $`\chi _{[a,b]}(\lambda )`$ is the indicator of the interval $`[a,b]`$. It is easy to see that
$$𝐄\{n^1\mathrm{Tr}M_n^2\}=2w^2<\mathrm{}.$$
Denoting $`N_{sc,n}`$ the NCM of the random matrices defined by (2.28) we can rewrite this inequality in the form
$$_{\mathrm{}}^{\mathrm{}}\lambda ^2𝐄\{N_{sc,n}(d\lambda )\}<\mathrm{}.$$
(2.32)
Thus, if we use (2.29) as the second term in (2.16), it will satisfy condition (2.1). Taking $`f_{sc}(z)`$ as $`f_2(z)`$ in (2.18) we find from the two last equations of the system that $`\mathrm{\Delta }_2(z)/f(z)=2w^2f(z)`$ and then the first equation of (2.18) takes the form
$$f(z)=f_0(z+2w^2f(z)),$$
(2.33)
where $`f_0(z)`$ is the Stieltjes transform of the limiting counting measure of matrices $`H_{0,n}`$. This functional equation determining the limiting eigenvalue distribution of the deformed GUE was found by another method in (see also ) for random matrices (2.26) in which $`M_n`$ has independent (modulo the Hermitian symmetry conditions) entries, for (2.28) in particular.
###### Remark 3
Consider now a probability measure $`m(\mathrm{d}\lambda )`$ and assume that its second moment $`m_2`$ is finite. In this case we can write the Stieltjes transform $`s(z)`$ of $`m`$ in the form
$$s(z)=(z+\mathrm{\Sigma }(z))^1,$$
where $`\mathrm{\Sigma }(z)`$ is the Stieltjes transform of a non-negative measure whose total mass is $`m_2`$ (to prove this fact one can use, for example, the general integral representation for functions satisfying (2.13) ). Since $`s^{}(z)=z^2(1+o(1))`$, $`z\mathrm{},`$ then, according to the local inversion theorem, there exists a unique functional inverse $`z(s)`$ of $`s(z)`$ defined and analytic in a neighborhood of zero and assuming its values in a neighborhood of infinity. Denote
$$\mathrm{\Sigma }(z(s))=R_m(s)$$
(2.34)
and following Voiculescu call $`R_m(s)`$ the $`R`$-transform of the probability measure $`m`$. By using the $`R`$-transforms $`R_{1,2}`$ of measures $`N_{1,2}`$ we can rewrite the first two equations of system (2.18) in the form
$$\frac{\mathrm{\Delta }_{1,2}}{f(z)}=\frac{1}{f(z)}+z+R_{2,1}(f(z))=R(f(z))+R_{2,1}(f(z)),$$
(2.35)
where $`R`$ denotes the $`R`$-transform of the limiting normalized counting measure $`N`$ of the ensemble (2.1) (the measure whose Stieltjes transform is $`f`$). These relations and the third equation of system (2.18) lead to the remarkably simple expression of $`R`$ via $`R_1`$ and $`R_2`$
$$R(f)=R_1(f)+R_2(f),$$
(2.36)
that ”linearizes” the rather complex system (2.18). The relation was obtained by Voiculescu in the context of $`C^{}`$-algebra studies (see for results and references). Thus, one can regard the system (2.18) as a version of the binary operation on measures defined by (2.36) and known as the non-commutative convolution. A simple precursor of relation (2.36) containing the functional inverses of $`f`$ and $`f_{1,2}`$ for real $`z`$ lying outside of the support of $`N_0`$ in (2.24) was used in (see also ) to locate the support of $`N`$ in terms of the support of $`N_0`$ in the case of ensemble (2.20). The simplest form of the relation (2.36) for the case when both measures are semicircle measures (2.31), i.e. when $`R_{1,2}=2w_{1,2}^2f`$, was indicated in . Formal derivation of relation (2.36) for the case then matrices $`H_1`$ and $`H_2`$ distributed both according to the laws
$$P_{1,2}^{(n)}(\mathrm{d}H)=Z_{1,2}^{(n)}exp\{nV_{1,2}(H)\}\mathrm{d}H.$$
(2.37)
where $`V_{1,2}:_+`$ are polynomials of an even degree was given in . The derivation is based on the perturbation theory with respect to the non-quadratic part of $`V_{1,2}`$ and the $`R`$-transform is related to the sum of irreducible diagrams of the formal perturbation series. Existence of the limiting eigenvalue counting measure for the random matrix ensemble (2.37) was rigorously proved in for a rather broad class of functions $`V`$ (not necessary polynomials). It was also proved that the normalized counting measure (2.2) converges in probability to the limiting measure. The form (2.29) of matrices of ensemble (2.37) can be deduced from known results on the ensemble (2.37) (see e.g.) in the same way as for the GUE (2.28), where $`V(\lambda )=\lambda ^2/4w^2`$ (see ). Condition (2.17) follows from results of . Thus we can apply Theorem 2.1 to obtain rigorously relation (2.36) in the case when matrices $`H_r,r=1,2`$ in (2.1) are distributed according to (2.37).
###### Remark 4
The problem of addition of random Hermitian (real symmetric matrices) has a natural multiplicative analogues in the case of positive defined Hermitian (real symmetric) or unitary (orthogonal) matrices. Namely, assuming that $`A_n`$ and $`B_n`$ are positive defined matrices and $`U_n`$ is the unitary (orthogonal) Haar distributed random matrix we can consider the positive defined random matrix
$$H_n=A_n^{1/2}U_n^{}B_nU_nA_n^{1/2}.$$
(2.38)
Likewise, if $`S_n`$ and $`T_n`$ are unitary (orthogonal) matrices and and $`U_n`$ is as above we can consider the random matrices
$$V_n=S_nU_n^{}T_nU_n.$$
(2.39)
In this case the normalized eigenvalue counting measure is defined as $`n^1`$ times the number of eigenvalues belonging to a Borel set of the unit circle.
In both cases (2.38) and (2.39) one can study the limiting properties of the NCM’s of respective random matrices provided that the ”input” matrices $`A_n,B_n,S_n`$ and $`T_n`$ have limiting eigenvalue distributions. The first examples of ensembles of the above forms as multiplicative analogues of the ensemble (2.20) were proposed in , where the respective functional equations analogous to (2.24) were derived. A general class of the random matrix ensembles of these forms were studied in free probability , where the notions of the $`S`$ \- transform and the free multplicative convolution of measures were proposed and used to give a general form of the limiting eigenvalue distributions of products (2.38) and (2.39). It will be shown in the subsequent paper that a version of the method of this paper leads to results, analogous to those given in Theorem 2.1 above.
## 3 Convergence with Probability 1 for non-Random $`A_n`$ and $`B_n`$.
As the first step of the proof of Theorem 2.1 we prove the following
###### Theorem 3.1
Let $`H_n`$ be the random $`n\times n`$ matrix of the form (2.1) in which $`A_n`$ and $`B_n`$ are non-random Hermitian matrices, $`U_n`$ and $`V_n`$ are random independent unitary matrices distributed each according to the normalized to unity Haar measure on $`U(n)`$. Assume that the normalized counting measures $`N_{r,n},r=1,2`$ of matrices $`A_n`$ and $`B_n`$ converge weakly as $`n\mathrm{}`$ to nonnegative and normalized to 1 measures $`N_r,r=1,2`$ respectively and that
$$\underset{n}{sup}\lambda ^4N_{r,n}(\mathrm{d}\lambda )=m_4<\mathrm{},r=1,2.$$
(3.1)
Then the normalized eigenvalue counting measure (2.2) of $`H_n`$ converges with probability 1 to a non-random and normalized to 1 measure whose Stieltjes transform (2.4) is a unique solution of the system (2.18) in the class of functions $`f(z)`$, $`\mathrm{\Delta }_r(z),r=1,2`$ analytic for $`\mathrm{Im}z0`$ and satisfying conditions (2.12)-(2.14) and (2.19) respectively.
###### Remark 1
The theorem generalizes the results of proved under the condition that supports of the NCM $`N_{r,n},r=1,2`$ of $`A_n`$ and $`B_n`$ are uniformly bounded in $`n.`$
###### Remark 2
By mimicking the proof of the Glivenko - Cantelli theorem (see e.g. ), one can prove that the random distribution functions $`N_n(\lambda )=N_n(]\mathrm{},\lambda [)`$ corresponding to measures (2.2) converge uniformly with probability 1 to the distribution function $`N(\lambda )=N(]\mathrm{},\lambda [)`$ corresponding to measure $`N`$:
$$𝐏\{\underset{n\mathrm{}}{lim}\underset{\lambda }{sup}|N_n(\lambda )N(\lambda )|=0\}=1.$$
We present now our technical means. First is a collection of elementary facts of linear algebra.
###### Proposition 3.1
Let $`𝐌_n`$ be the algebra of linear transformations of $`^n`$ in itself ($`n\times n`$ complex matrices) equipped with the norm, induced by the Euclidean norm of $`^n`$.
We have :
1. if $`M𝐌_n`$ and $`\{M_{jk}\}_{j,k=1}^n`$ is the matrix of $`M`$ in any orthonormalized basis of $`^n`$, then
$$|M_{jk}|M;$$
(3.2)
2. if $`\mathrm{Tr}M=\underset{j=1}{\overset{n}{}}M_{jj}`$, then
$$|\mathrm{Tr}M_1M_2|(\mathrm{Tr}M_1M_1^{})^{1/2}(\mathrm{Tr}M_2M_2^{})^{1/2},$$
(3.3)
where $`M^{}`$ is the Hermitian conjugate of $`M`$, and if $`P`$ is a positive defined transformation, then
$$|\mathrm{Tr}MP|M\mathrm{Tr}P;$$
(3.4)
3. for any Hermitian transformation $`M`$ its resolvent
$$G(z)=(Mz)^1$$
(3.5)
is defined for all non-real $`z`$, $`\mathrm{Im}z0`$,
$$G(z)|\mathrm{Im}z|^1$$
(3.6)
and if $`\{G_{jk}(z)\}_{j,k=1}^n`$ is the matrix of $`G(z)`$ in any orthonormalized basis of $`^n`$ then
$$|G_{jk}(z)||\mathrm{Im}z|^1;$$
(3.7)
4. if $`M_1`$ and $`M_2`$ are two Hermitian transformations and $`G_r(z),r=1,2`$ are their resolvents, then
$$G_2(z)=G_1(z)G_1(z)(M_2M_1)G_2(z)$$
(3.8)
(the resolvent identity);
5. if $`G(z)=(Mz)^1`$ is regarded as a function of $`M,`$ then the derivative $`G^{}(z)`$ of $`G(z)`$ with respect to $`M`$ verifies the relation
$$G^{}(z)X=G(z)XG(z)$$
(3.9)
for any Hermitian $`X𝐌_n`$, and, in particular,
$$G^{}(z)G(z)^2|\mathrm{Im}z|^2$$
(3.10)
Now is our main technical tool.
###### Proposition 3.2
Let $`\mathrm{\Phi }:𝐌_n`$ be a continuously differentiable function. Then the following relation holds for any $`M𝐌_n`$ and any Hermitian element $`X𝐌_n`$:
$$\underset{U(n)}{}\mathrm{\Phi }^{}(U^{}MU)[X,U^{}MU]dU=0,$$
(3.11)
where
$$[M_\mathit{1},M_\mathit{2}]=M_\mathit{1}M_\mathit{2}M_\mathit{1}M_\mathit{2}$$
(3.12)
is the commutator of $`M_\mathit{1}`$ and $`M_\mathit{2}`$and the symbol
$$\underset{U(n)}{}\mathrm{}\mathrm{𝑑𝑈}$$
(3.13)
denotes the integration over $`U(n)`$ with respect to the normalized Haar measure $`\mathrm{𝑑𝑈}`$.
Proof. To prove (3.11) we use the right shift invariance of the Haar measure: $`\mathrm{d}U=\mathrm{d}(UU_0)`$, $`U_0U(n)`$ according to which the integral
$$\underset{U(n)}{}\mathrm{\Phi }\left(e^{i\epsilon X}U^{}MUe^{i\epsilon X}\right)dU$$
is independent of $`\epsilon `$ for any Hermitian $`X𝐌_n`$. Thus its derivative with respect to $`\epsilon `$ at $`\epsilon =0`$ is zero. This derivative is the l.h.s. of (3.11).$`\mathrm{}`$
###### Proposition 3.3
System (2.18) has a unique solution in the class of functions $`f(z)`$, $`\mathrm{\Delta }_{1,2}(z)`$ analytic for $`\mathrm{Im}z0`$ and satisfying conditions (2.12)–(2.14) and (2.19).
Proof. Assume that there exist two solutions $`(f^{^{}},\mathrm{\Delta }_{1,2}^{^{}})`$ and $`(f^{^{\prime \prime }},\mathrm{\Delta }_{1,2}^{^{\prime \prime }})`$ of the system. Denote $`\delta f=f^{^{}}f^{^{\prime \prime }}`$, $`\delta \mathrm{\Delta }_{1,2}=\mathrm{\Delta }_{1,2}^{^{}}\mathrm{\Delta }_{1,2}^{^{\prime \prime }}`$. Then, by using (2.18) and the integral representation (2.5) for $`f_{1,2},`$ we obtain the linear system for $`\delta \varphi =z\delta f`$, and for $`\delta \mathrm{\Delta }_{1,2}`$
$$\begin{array}{cc}\delta \varphi (1a_1(z))+b_1(z)\delta \mathrm{\Delta }_1\hfill & =0,\hfill \\ \delta \varphi (1a_2(z))+b_2(z)\delta \mathrm{\Delta }_2\hfill & =0,\hfill \\ \delta \varphi \delta \mathrm{\Delta }_1\delta \mathrm{\Delta }_2\hfill & =0,\hfill \end{array}$$
(3.14)
where
$$a_1=\frac{\mathrm{\Delta }_1^{^{\prime \prime }}}{f^{^{}}f^{^{\prime \prime }}}I_2,b_1=\frac{z}{f^{^{}}}I_2,I_2=I_2(z\mathrm{\Delta }_1^{^{}}/f^{^{}},z\mathrm{\Delta }_1^{^{\prime \prime }}/f^{\prime \prime }),$$
(3.15)
$$I_2(z^{},z^{\prime \prime })=\frac{N_2(\mathrm{d}\lambda )}{(\lambda z^{})(\lambda z^{\prime \prime })},$$
(3.16)
and $`a_2`$, $`b_2`$ can be obtained from $`a_1`$ and $`b_1`$ by replacing $`N_2`$ and $`\mathrm{\Delta }_1`$ by $`N_1`$ and $`\mathrm{\Delta }_2`$ in above formulas. For any $`y_0>0`$ consider the domain
$$E(y_0)=\{z:|\mathrm{Im}z|y_0,|\mathrm{Re}z||\mathrm{Im}z|\}.$$
(3.17)
If $`s(z)`$ is the Stieltjes transform (2.11) of a probability measure $`m,`$ then we have for $`zE(y_0)`$,
$$\left|\frac{\lambda m(\mathrm{d}\lambda )}{\lambda z}\right|=\left|\underset{|\lambda |M}{}+\underset{|\lambda |>M}{}\right|\frac{M}{y_0}+2\underset{|\lambda |>M}{}m(\mathrm{d}\lambda ),$$
i.e.
$$zs(z)=1+o(1),z\mathrm{},zE(y_0).$$
(3.18)
Analogously, by using this asymptotic relation and condition (2.19) we obtain that for $`z\mathrm{}`$, $`zE(y_0)`$
$$z^2I_{1,2}(z)=1+o(1),a_{1,2}(z)=o(1),b_{1,2}(z)=1+o(1).$$
Thus the determinant $`b_1b_2+b_1+b_2(a_2b_1+a_1b_2)`$ of system (3.14) is equal asymptotically to $`1`$. We conclude that if $`y_0`$ in (3.17) is big enough, then system (3.14) has only a trivial solution, i.e. system (2.18) is uniquely soluble. $`\mathrm{}`$
In what follows we use the notation
$$\underset{U(n)}{}\mathrm{}dU=\mathrm{}$$
(3.19)
Proof of Theorem 3.1. Because of unitary invariance of eigenvalues of a Hermitian matrices we can assume without loss of generality that the unitary matrix $`V`$ in (2.1) is set to unity, i.e. we can work with the random matrix (2.16). We will omit below the subindex $`n`$ in all cases when it will not lead to confusion. Write the resolvent identity (3.8) for the pair $`(H_1,H)`$ of (2.1):
$$G(z)=G_1(z)G_1(z)H_2G(z),$$
(3.20)
where
$$G(z)=(H_1+H_2z)^1,G_1(z)=(H_1z)^1.$$
Consider the matrix $`g_n(z)G(z),`$ where
$$g_n(z)=\frac{1}{n}\mathrm{Tr}G(z)=\frac{N_n(\mathrm{d}\lambda )}{\lambda z},\mathrm{Im}z0$$
(3.21)
is the Stieltjes transform of random measure (2.2). The resolvent identity (3.20) leads to the relation
$$g_n(z)G(z)=g_n(z)G_1(z)G_1(z)g_n(z)H_2G(z).$$
(3.22)
By using Proposition 3.2 with the matrix element $`\left((H_1+Mz)^1\right)_{ac}`$ as $`\mathrm{\Phi }(M)`$ we have in view of (3.9) and (3.11) - (3.12)
$$(G[X,H_2]G)_{ac}=0.$$
Choosing the Hermitian matrix $`X`$ with only $`(a,b)`$-th and $`(b,a)`$ non-zero entries, we obtain
$$G_{aa}(H_2G)_{bc}=(GH_2)_{aa}G_{bc}.$$
(3.23)
Applying to this relation the operation $`n^1\underset{a=1}{\overset{n}{}}`$ and taking into account the definition (3.21) of $`g_n(z)`$ we rewrite the last relation in the form
$$g_n(z)H_2G(z)=\delta _{2,n}(z)G(z),$$
where
$$\delta _{2,n}(z)=\frac{1}{n}\mathrm{Tr}H_2G(z).$$
(3.24)
Thus we can rewrite (3.22) as
$$g_n(z)G(z)=g_n(z)G_1(z)G_1(z)\delta _{2,n}(z)G(z).$$
(3.25)
Introduce now the centralized quantities
$$g_n^{}(z)=g_n(z)f_n(z),\delta _{2,n}^{}(z)=\delta _{2,n}(z)\mathrm{\Delta }_{2,n}(z),$$
(3.26)
where
$$f_n(z)=g_n(z),\mathrm{\Delta }_{2,n}(z)=\delta _{2,n}(z).$$
(3.27)
With these notations (3.25) becomes
$$f_n(z)G(z)=f_n(z)G_1(z)\mathrm{\Delta }_{2,n}(z)G_1(z)G(z)+R_{1,n}(z),$$
(3.28)
where
$$R_{1,n}(z)=g_n^{}(z)G(z)G_1(z)\delta _{2,n}^{}(z)G(z).$$
(3.29)
Besides, since
$$\begin{array}{c}n^1\mathrm{Tr}H^2=n^1\mathrm{Tr}(H_1+H_2)^22n^1\mathrm{Tr}H_1^2+2n^1\mathrm{Tr}H_2^2=\hfill \\ =2\lambda ^2N_{1,n}(\mathrm{d}\lambda )+2\lambda ^2N_{2,n}(\mathrm{d}\lambda )4m_24m_4^{1/2},\hfill \end{array}$$
(3.30)
we have
$$\mu _2\underset{n}{sup}(n^1\mathrm{Tr}H^2)=\underset{n}{sup}\lambda ^2N_n(\mathrm{d}\lambda )4m_24m_4^{1/2}<\mathrm{}.$$
(3.31)
Thus
$$g_n(z)=\frac{N_n(\mathrm{d}\lambda )}{\lambda z}=\frac{1}{z}+\widehat{g}_n(z),$$
where
$$\widehat{g}_n(z)=\frac{\lambda N_n(\mathrm{d}\lambda )}{(\lambda z)z}.$$
In view of (3.31)
$$|z\widehat{g}_n(z)||\mathrm{Im}z|^1|\lambda |N_n(\mathrm{d}\lambda )|\mathrm{Im}z|^1m_4^{1/4},$$
i.e. the asymptotic relation
$$g_n^1(z)=z\left(1+O\left(\frac{1}{|\mathrm{Im}z|}\right)\right),\mathrm{Im}z\mathrm{}$$
(3.32)
holds uniformly in $`n`$. We have also the simple bound
$$|g_n(z)||\mathrm{Im}z|^1$$
(3.33)
following from (3.4) and (3.7) and, in addition, according to Proposition 3.1 and (3.24), the bounds
$$|\delta _{2,n}(z)|m_4^{1/4}|\mathrm{Im}z|^1,$$
(3.34)
$$z\delta _{2,n}(z)=n^1\mathrm{Tr}H_2zG(z)=n^1\mathrm{Tr}H_2(1+HG(z)).$$
(3.35)
Hence, in view of (3.31)
$$\begin{array}{c}|z\delta _{2,n}(z)|(n^1\mathrm{Tr}H_2^2)^{1/2}+(n^1\mathrm{Tr}H_2^2)^{1/2}(n^1\mathrm{Tr}H^2G(z)G^{}(z))^{1/2}\\ m_4^{1/4}+2m_4^{1/2}/y_0,\end{array}$$
(3.36)
i.e. $`z\delta _{2,n}(z)`$ is uniformly bounded in $`n`$.
As a result of above bounds we have for $`|\mathrm{Im}z|y_0`$ uniformly in $`n`$
$$\mathrm{\Delta }_{2,n}(z)f_n^1(z)G_1(z)=O\left(\frac{1}{y_0}\right),y_0\mathrm{}$$
i.e. the matrix $`1\mathrm{\Delta }_{2,n}(z)f_n^1(z)G_1(z)`$ is invertible uniformly in $`n`$ and there is $`y_0`$ independent of $`n`$ and such that for $`|\mathrm{Im}z|y_0`$
$$(1+\mathrm{\Delta }_{2,n}(z)f_n^1(z)G_1(z))^12.$$
(3.37)
Thus (3.28) is equivalent to
$$\begin{array}{c}G(z)=(1+\mathrm{\Delta }_{2,n}(z)f_n^1(z)G_1(z))^1G_1(z)+\\ (1+\mathrm{\Delta }_{2,n}(z)f_n^1(z)G_1(z))^1f_n^1(z)R_{1,n}(z)\end{array}$$
or to
$$G(z)=G_1\left(z\mathrm{\Delta }_{2,n}(z)f_n^1(z)\right)+(1+\mathrm{\Delta }_{2,n}(z)f_n^1(z)G_1(z))^1f_n^1(z)R_{1,n}(z).$$
Applying to this relation the operation $`n^1\mathrm{Tr}`$ we obtain
$$f_n(z)=f_{1,n}(z\mathrm{\Delta }_{2,n}(z)f_n^1(z))+r_{1,n}(z),$$
(3.38)
where
$$f_{1,n}(z)=n^1\mathrm{Tr}G_1(z)=\frac{N_{1,n}(\mathrm{d}\lambda )}{\lambda z}$$
(3.39)
is the Stieltjes transform of the normalized counting measure of $`H_{1,n}`$ in (2.1) and
$$r_{1,n}(z)=n^1\mathrm{Tr}(1+\mathrm{\Delta }_{2,n}(z)f_n^1(z)G_1(z))^1f_n^1(z)R_{1,n}(z),$$
(3.40)
where $`R_{1,n}(z)`$ is defined in (3.29). We show in the next Theorem 3.2 that there exists a sufficiently big $`y_0>0`$ and $`C(y_0)>0`$, both independent of $`n`$ and such that if $`zE(y_0),`$ where $`E(y_0)`$ is defined in (3.17), then the variances
$$v_1(z)=|g_n^{}(z)|^2,v_2(z)=|\delta _{2,n}^{}(z)|^2$$
(3.41)
admit the bounds
$$v_1(z)\frac{C(y_0)}{n^2},v_2(z)\frac{C(y_0)}{n^2}.$$
(3.42)
These bounds, Proposition 3.1, (3.37), and Schwartz inequality for the expectation $`\mathrm{}`$ imply that uniformly in $`n`$ and in $`zE(y_0)`$
$$|r_{1,n}(z)|\frac{2C^{1/2}(y_0)}{n}(1+y_0^1)|f_n^2(z)n^1\mathrm{Tr}G(z)G^{}(z)|^2^{1/2}.$$
In view of (3.27), (3.32) and the identity $`zG(z)=1+HG(z)`$ we have
$$f_n^1(z)G(z)=z(1+O(y_0^1))G(z)=(1+O(y_0^1))(1HG(z)),$$
and since, by (3.3), (3.4) and (3.30)
$$\begin{array}{c}|n^1\mathrm{Tr}HG(z)|y_0^1n^1\mathrm{Tr}H^2\\ 2m_4^{1/4}y_0^1,|n^1\mathrm{Tr}H^2G(z)G^{}(z)|4m_4^{1/2}y_0^2,\end{array}$$
we obtain that for $`zE(y_0)`$
$$|r_{1,n}(z)|\frac{C_1(y_0)}{n},$$
(3.43)
where $`C_1(y_0)`$ is independent of $`n`$ and is bounded in $`y_0`$.
Furthemore, the bounds (3.33) and (3.34) imply that sequences $`\{f_n(z)\}`$ and $`\{\mathrm{\Delta }_{2,n}(z)\}`$ are analytic and uniformly in $`n`$ bounded for $`|\mathrm{Im}z|y_0>0`$. Thus the sequences are compact with respect to uniform convergence on compacts of the domain
$$D(y_0)=\{z:|\mathrm{Im}z|y_0>0\}.$$
(3.44)
In addition, according to the hypothesis of the theorem, the normalized counting measures $`N_{1,n}`$ of matrices $`H_{1,n}`$ converge weakly to a limiting probability measure $`N_1`$ Thus, their Stieltjes transforms (3.39) converge uniformly on compacts of (3.44) to the Stieltjes transform $`f_1`$ of $`N_1`$. Hence, if $`y_0>0`$ is large enough, there exist two analytic in (3.44) functions $`f`$ and $`\mathrm{\Delta }_2`$ verifying the relation
$$f(z)=f_1\left(z\frac{\mathrm{\Delta }_2(z)}{f(z)}\right),|\mathrm{Im}z|y_0.$$
This is the first equation of system (2.18). The second equation of the system follows from the argument above in which the roles $`H_1`$ and $`H_2`$ are interchanged, in particular the quantity $`n^1\mathrm{Tr}H_1G(z)`$ is denoted $`\mathrm{\Delta }_{1,n}(z)`$. As for the third equation, it is just the limiting form of the identity
$$n^1\mathrm{Tr}(H_{1,n}+H_{2,n}z)G(z)=1.$$
(3.45)
Thus, we have derived system (2.18). Its unique solubility in domain (3.17) where $`y_0`$ is large enough is proved in Proposition 3.3. Besides, all three functions $`f_n`$, $`\mathrm{\Delta }_{r,n},r=1,2`$ defined in (3.27) are a priory analytic for $`|\mathrm{Im}z|>0`$. Thus, their limits $`f,\mathrm{\Delta }_r,r=1,2`$ are also analytic for non-real $`z`$. In view of the weak compactness of probability measures and the continuity of the one-to-one correspondence between nonnegative measures and their Stieltjes transforms (see Proposition 2.1(v)) there exists a unique nonnegative measure $`N`$ such that $`f`$ admit the representation (2.4). The measure $`N`$ is a probability measure in view of (3.32) and.(2.14).
We conclude that the whole sequence $`\{f_n\}`$ of expectations (3.27) of the Stieltjes transforms $`g_n`$ (3.21) of measures (2.2) converges uniformly on compacts of $`D(y_0),`$ where $`D(y_0)`$ is defined in (3.44), to the limiting function $`f`$ verifying (2.18). This result, Theorem 3.2 and the Borel-Cantelli lemma imply that the sequence $`\{g_n(z)\}`$ converges with probability 1 to $`f(z)`$ for any fixed $`zD(y_0).`$ Since the convergence of a sequence of analytic functions on any countable set having an accumulation point in their common domain of definition implies the uniform convergence of the sequence on any compact of the domain, we obtain the convergence $`g_n`$ to $`f`$ with probability 1 on any compact of $`D(y_0)`$. Due to the continuity of the one-to-one correspondence between probability measures and their Stieltjes transforms (see Proposition 2.1(v)) the normalized eigenvalue counting measure (2.2) of the eigenvalues of random matrix (2.1) converge weakly with probability 1 to the nonrandom measure $`N`$ whose Stieltjes transform (2.4) satisfies (2.18).$`\mathrm{}`$
###### Theorem 3.2
Let $`H_n`$ be the random matrix of the form (2.1) satisfying the condition of Theorem 3.1. Denote
$$g_n(z)=n^1\mathrm{Tr}(H_nz)^1,\delta _{r,n}(z)=n^1\mathrm{Tr}H_{r,n}(H_nz)^1,r=1,2.$$
(3.46)
Then there exist $`y_0`$ and $`C(y_0)`$, both positive and independent of $`n`$ and such that the variances of random variables (3.46) admit the bounds for $`|\mathrm{Im}z|y_0`$
$$|g_n(z)g_n(z)|^2\frac{C(y_0)}{n^2}$$
(3.47)
$$|\delta _{r,n}(z)\delta _{r,n}(z)|^2\frac{C(y_0)}{n^2},r=1,2,.$$
(3.48)
if $`zE(y_0)`$, where $`E(y_0)`$ is defined in (3.17).
Proof. Because of the symmetry of the problem with respect to $`H_1`$ and $`H_2`$ in (2.1) it suffices to prove (3.48) for, say, $`\delta _{2,n}(z)`$. Besides, we will use bellow the notations $`g(z)`$ and $`\delta (z)`$ for $`g_n(z)`$ and $`\delta _{2,n}(z)`$ and the notations 1 and 2 for two values $`z_1`$ and $`z_2`$ of the complex spectral parameter $`z`$. We assume that $`|\mathrm{Im}z_{1,2}|y_0>0`$.
We will use the same approach as in the proof of Theorem 3.1, i.e. we will derive and study certain relations obtained by using Proposition 3.2 and the resolvent identity.
Consider the matrix
$$V_1=g^{}(1)G(2),$$
(3.49)
where $`g^{}(1)=g(1)g(1)`$. Its clear that $`n^1\mathrm{Tr}V_1`$ for $`z_1=z`$ and $`z_2=\overline{z}`$ is the variance (3.47), that we denoted $`v_1(z)`$ in (3.41):
$$|g^{}(z)|^2=n^1\mathrm{Tr}V_1|_{z_1=z,z_2=\overline{z}}=v_1(z).$$
(3.50)
In view of the resolvent identity (3.20) for the pair $`(H_1,H)`$ we have
$$V_1=G_1(2)W,$$
(3.51)
$$W=g^{}(1)H_2G(2).$$
(3.52)
Applying Proposition 3.2 to the function
$$\mathrm{\Phi }(M)=G_{aa}^{}(1)(MG(2))_{cd},$$
where $`G(z)=(H_1+Mz)^1`$, and
$`G^{}(z)`$ $`=`$ $`G(z)G(z)=`$
$`(H_1+Mz)^1{\displaystyle _{U(n)}}(H_1+U^{}BUz)^1𝑑U,`$
we obtain the relation
$$\begin{array}{c}(G(1)[X,H_2]G(1))_{aa}(H_2G(2))_{cd}+G_{aa}^{}(1)([X,H_2]G(2))_{cd}\\ G_{aa}^{}(1)(H_2G(2)[X,H_2]G(2))_{cd}=0,\end{array}$$
where the operation $`[\mathrm{},\mathrm{}]`$ is defined in (3.12). Choosing as $`X`$ the Hermitian matrix having only the $`(c,j)`$-th and $`(j,c)`$ non-zero entries, we obtain from the above relation the following one:
$$G_{ac}(1)(H_2G(1))_{ja}(H_2G(2))_{cd}+(G(1)H_2)_{ac}G_{ja}(1)(H_2G(2))_{cd}+$$
$$+G_{aa}^{}(1)\delta _{cc}(H_2G(2))_{jd}G_{aa}^{}(1)(H_2)_{cc}G_{jd}(2)$$
$$G_{aa}^{}(1)(H_2G(2))_{cc}(H_2G(2))_{jd}+G_{aa}^{}(1)(H_2G(2)H_2)_{cc}G_{jd}(2)=0$$
Applying to this relation the operation $`n^1\underset{ac}{}`$ and taking into account that
$$g^{}=n^1\underset{a}{}G_{aa}^{},$$
we have
$$\begin{array}{c}n^2[G^2(1),H_2]H_2G(2)+g^{}(1)H_2G(2)+\\ +g^{}(1)k(2)G(2)g^{}(1)\delta (2)H_2G(2)=0,\end{array}$$
(3.53)
where
$$k(z)=n^1\mathrm{Tr}K(z),K(z)=BG_U(z)BB,G_U(z)=UG(z)U^{}.$$
(3.54)
Introducing the centralized quantity (cf.(3.26))
$$k^{}=kk,$$
(3.55)
and using our notations (3.24) and (3.27), we can rewrite (3.53) as
$$(1\mathrm{\Delta }(2))W=k(2)V_1+R,$$
(3.56)
where
$$R=g^{}(1)\delta ^{}(2)H_2G(2)g^{}(1)k^{}(2)G(2)T_1,$$
(3.57)
and
$$T_1=n^2[G^2(1),H_2]H_2G(2).$$
In view of the uniform in $`n`$ bound (3.36)), the function $`1\mathrm{\Delta }(z)`$ is uniformly in $`n`$ bounded away from zero. Thus we have from (3.51), (3.52) and (3.56)
$$V_1=\left(1k(2)(1\mathrm{\Delta }(2))^1G_1(2)\right)^1(1\mathrm{\Delta }(2))^1G_1(2)R.$$
(3.58)
According to (3.54), (3.6) and (3.1), we have uniformly in $`n`$
$$|k(z)|y_0^1n^1\mathrm{Tr}B^2+|n^1\mathrm{Tr}B|y_0^1m_4^{1/2}+m_4^{1/4}<\mathrm{}.$$
(3.59)
This bound and universal bound (3.6) imply that the matrix $`(1k(z)(1\mathrm{\Delta }(z))^1G_1(z))`$ is uniformly in $`n`$ invertible if $`|\mathrm{Im}z|y_0`$ and $`y_0`$ is large enough, and hence the matrix
$$Q=\left(1k(z)(1\mathrm{\Delta }(z))^1G_1(z)\right)^1(1\mathrm{\Delta }(z))^1G_1(z)$$
admits the following bound for $`|\mathrm{Im}z|y_0`$ and sufficiently large $`y_0`$
$$Q\frac{C}{y_0},$$
(3.60)
where $`C`$ is an absolute constant.
Setting now in (3.58) $`z_1=z`$, $`z_2=\overline{z}`$ and applying to this relation the operation $`n^1\mathrm{Tr}`$ we obtain in the l.h.s. the variance $`v_1(z)`$ because of (3.50). As for the r.h.s., its terms can be estimated as follows in view of (3.57):
1. $$|g^{}(1)\delta ^{}(2)n^1\mathrm{Tr}QH_2G(2)|\alpha _{12}(y_0)v_1^{1/2}v_2^{1/2},$$
(3.61)
where $`v_2`$ is defined in (3.41) and because, according to (3.1), (3.3), (3.6) and (3.60),
$$\begin{array}{c}|n^1\mathrm{Tr}QH_2G(2)|(n^1\mathrm{Tr}Q^{}Q)^{1/2}(n^1\mathrm{Tr}H_2^2G(2)G^{}(2))^{1/2}\\ Cy_0^2m_4^{1/4}\alpha _{12}(y);\end{array}$$
(3.62)
2. $$|g^{}(1)k^{}(2)n^1\mathrm{Tr}QG(2)|\alpha _{13}(y_0)v_1^{1/2}v_3^{1/2},$$
(3.63)
where
$$v_3=|k^{}(z)|^2$$
(3.64)
because
$$\begin{array}{c}|n^1\mathrm{Tr}QG(2)|(n^1\mathrm{Tr}Q^{}Q)^{1/2}(n^1\mathrm{Tr}G(2)G^{}(2))^{1/2}\\ Cy_0^2\alpha _{13}(y_0);\end{array}$$
(3.65)
3. $$|n^3\mathrm{Tr}(Q[G^2(1),H_2]H_2G(2))|Cm_4^{1/2}y_0^4n^2\frac{\beta _1(y_0)}{n^2}.$$
Thus we obtain the inequality
$$v_1\alpha _{12}(y_0)v_1^{1/2}v_2^{1/2}+\alpha _{13}(y_0)v_1^{1/2}v_3^{1/2}+\frac{\beta _1(y_0)}{n^2},$$
(3.66)
where $`\alpha _{12}`$, $`\alpha _{13}`$ and $`\beta _1`$ are independent on $`n`$ and vanish as $`y_0\mathrm{}`$.
Now we are going to derive analogous inequalities for $`v_2`$ and $`v_3`$ defined in (3.41) and in (3.64) and to obtain the system
$$v_i\underset{j=1,ji}{\overset{3}{}}\alpha _{ij}v_i^{1/2}v_j^{1/2}+\frac{\beta _i(y_0)}{n^2},i=1,2,3.$$
(3.67)
To get the second inequality of the system we consider the matrix (cf. (3.49))
$$V_2=\delta ^{}(1)H_2G(2).$$
(3.68)
Applying to $`V_2`$ operation $`n^1\mathrm{Tr}`$ and setting $`z_1=z`$, $`z_2=\overline{z,}`$ we obtain the variance $`v_2`$ of (3.42). On the other hand, using Proposition 3.2 for the function
$$\mathrm{\Phi }(M)=(MG(1))_{aa}^{}(MG(2))_{cd},$$
we obtain, after performing in essence the same procedure as that used in the derivation of (3.53), in particular, choosing the Hermitian matrix $`X`$ with only the $`(c,j)`$-th and $`(j,c)`$ non-zero entries,
$$v_2=g(2)\delta ^{}(1)k(2)+\delta ^{}(1)\delta ^2(2)T_2,$$
(3.69)
where
$$T_2=n^3\mathrm{Tr}([G_U(1),K(1)]BG(2))$$
(3.70)
and $`K(z)`$, $`k(z)`$ are defined in (3.54). Using again centralized quantities (3.26) and (3.55), we can write
$$g(2)\delta ^{}(1)k(2)=g^{}(2)\delta ^{}(1)k(2)+g(2)\delta ^{}(1)k^{}(2)$$
and
$$\delta ^{}(1)\delta ^2(2)=\delta ^{}(1)\delta ^{}(2)\delta (2)+\delta ^{}(1)\delta ^{}(2)\delta (2).$$
Thus, in view of (3.33), (3.34), (3.59), and Schwarz inequality we have the bounds
$$g(2)\delta ^{}(1)k(2)v_1^{1/2}v_2^{1/2}m_4^{1/4}(1+m_4^{1/4}y_0^1)+v_2^{1/2}v_3^{1/2}y_0^1,$$
and
$$\delta ^{}(1)\delta ^2(2)2v_2m_4^{1/4}y_0^1.$$
These bounds and analogously obtained bound for $`T_2`$ in (3.70) lead for $`m_4^{1/4}y_0^11/4`$ to the second inequality (3.67), in which
$$\alpha _{21}(y_0)=4m_4^{1/4},\alpha _{23}(y_0)=2y_0^1,\beta _2=8m_4^{1/4}y_0^2.$$
(3.71)
To obtain the third inequality of (3.67) we may use the same scheme as above applied to the matrix $`V_3=k^{}(1)K(2)`$ (cf. (3.49) and (3.68)). However this requires rather tedious computations and the existence of the uniformly bounded in $`n`$ sixth moment $`m_6`$ of the measure $`N_{2,n}`$. For this reason we consider the quantity
$$n^1\mathrm{Tr}(BG_U(1)B)^{}G_U(2)B,$$
(3.72)
where $`G_U(z)`$ is defined in (3.54). As before we would like to obtain for this quantity a certain relation, basing on the invariance of the Haar measure with respect to the group shifts. To this end we will introduce the following function of the unitary matrix $`U`$:
$$(BUG(1)U^{}B)_{aa}^{}(UG(2)U^{}B)_{cd},$$
where $`G(z)=(H_1+U^{}BUz)^1`$ and we will use the analogue of (3.11) obtained from the left shift invariance of the Haar measure. This leads to the relation (cf. (3.53) and (3.69))
$$k^{}(1)g(2)K(2)+k^{}(1)\delta (2)G_U(2)Bk^{}(1)G_U(2)BT_3=0,$$
(3.73)
where
$$T_3=n^2G_U(1)BK(1)G_U(2)BK(1)BG_U(1)G_U(2)B.$$
We multiply (3.73) by $`B`$ from the left and introduce again the centralized quantities $`g^{}`$, $`\delta ^{}`$ and $`k^{}`$ defined in (3.26) and (3.55). We obtain
$$\begin{array}{c}(1\mathrm{\Delta }(2)f(2)B)k^{}(1)K(2)=k^{}(1)g^{}(2)BK(2)+\\ +k^{}(1)\delta ^{}(2)BG_U(2)B+BT_3.\end{array}$$
In view of (3.32) and (3.36) the imaginary part of the function $`1\mathrm{\Delta }(z)`$ is uniformly in $`n`$ bounded away from zero if $`|\mathrm{Im}z|`$ is large enough. Since $`B`$ is a Hermitian matrix, the matrix
$$S=(1\mathrm{\Delta }(2)f(2)B)^1$$
(3.74)
admits the bound
$$S=|f(2)|^1((1\mathrm{\Delta }(2))f^1(2)B)^1|f(2)|^1\left|\mathrm{Im}\frac{1\mathrm{\Delta }(2)}{f(2)}\right|^1.$$
By using (3.28) and (3.34) we find that for $`zE(y_0)`$, where $`E(y_0)`$ is defined in (3.17) with sufficiently big $`y_0`$, we have the uniform in $`n`$ inequality $`|f(2)\mathrm{Im}(1\mathrm{\Delta }(2))f^1(2)|1/2,`$ i.e.
$$S2.$$
(3.75)
This leads to the relation
$$\begin{array}{c}V_3k^{}(1)K(2)=k^{}(1)g^{}(2)SBK(2)+\\ +k^{}(1)\delta ^{}(2)SBG_U(2)B+SBT_3.\end{array}$$
(3.76)
We apply to this relation the operation $`n^1\mathrm{Tr}`$, set $`z_1=z`$, $`z_2=\overline{z}`$ and estimate the contribution of the two first terms of the r.h.s. as (3.76) as above, using in addition (3.75). We obtain
$$\begin{array}{c}|n^1\mathrm{Tr}SBK(2)|4m_4^{1/2}\alpha _{31}(y_0),\hfill \\ |n^1\mathrm{Tr}SBG_U(2)B|4m_4^{1/2}y_0^1\alpha _{32}(y_0).\hfill \end{array}$$
(3.77)
To estimate the third term of the r.h.s. of (3.76) we use the identity
$$SB=f^1(2)+(1\mathrm{\Delta }(2))f^1(2)S,$$
the asymptotic relations (3.32) and (3.34) and the bound (3.75). This yields the bound $`SB4y_0`$. By using this bound and the same reasoning as in obtaining other bounds above, we obtain
$$|n^1\mathrm{Tr}SBT_3|\frac{Cm_4}{y_0^2n^2}\frac{\beta _3}{n^2},$$
where $`C`$ is an absolute constant.
Let us introduce new variables
$$u_1=y_0v_1^{1/2},u_2=v_2^{1/2},u_3=v_3^{1/2}$$
(3.78)
Then we obtain from (3.67) and (3.62), (3.65), (3.71), and (3.77) the system
$$u_i^2\underset{j=1,ji}{\overset{3}{}}a_{ij}u_iu_j+\frac{\gamma _i}{n^2},$$
(3.79)
in which the coefficients $`\{a_{ij},ij\}`$ have the form $`a_{ij}=y_0^1b_{ij},`$ where $`b_{ij}`$ are bounded in $`y_0`$ and in $`n`$ as $`y_0\mathrm{}`$ and $`n\mathrm{}`$. By choosing $`y_0`$ sufficiently big (and then fixing it) we can guarantees that $`0a_{ij}1/4,ij`$. Thus summing the three relations (3.79) we can write the result in the form $`(\widehat{a}u,u)\gamma /n^2`$ where $`\gamma =\gamma _1+\gamma _2+\gamma _3`$ and ($`\widehat{a})_{ij}=\delta _{ij}+(1\delta _{ij})/4,i,j=1,2,3.`$ Since the minimum eigenvalue of the matrix $`\widehat{a}`$ is $`1/2,`$ we obtain from (3.78)bounds (3.47) and (3.48).$`\mathrm{}`$
## 4 Convergence in Probability
In this Section we prove Theorem 2.1. Since, according to Theorem 3.2 the randomness of $`U_n`$ in (2.1) (or (2.16)) provides already vanishing the variance of the Stieltjes transform of the NCM (2.2), we have only to prove that the additional randomness due to the matrices $`A_n`$ and $`B_n`$ in (2.1) does not destroy this property. We will prove this fact first for $`A_n`$ and $`B_n`$ whose norms are uniformly bounded in $`n`$ (see Lemma 4.1 below ), and then we will treat the general case of Theorem 2.1 by using a certain truncating procedure.
###### Proposition 4.1
Let $`\{m_n\}`$ be a sequence of random non-negative unit measures on the line and $`\{s_n\}`$ be the sequence of their Stieltjes transforms (2.11). Then the sequence $`\{m_n\}`$ converges weakly in probability to a nonrandom non-negative unit measure $`m`$ if and only if the sequence $`\{s_n\}`$ converges in probability for any fixed $`z`$ belonging to a compact $`K\{zC:\mathrm{Im}z>0\}`$ to the Stieltjes transform $`f`$ of the measure $`m`$.
Proof. Let us prove first the necessity. According to the hypothesis for any continuous and having a compact support function $`\phi (\lambda )`$ we have
$$\underset{n\mathrm{}}{lim}𝐏\left\{\left|\phi (\lambda )m(d\lambda )\phi (\lambda )m_n(d\lambda )\right|>\epsilon \right\}=0.$$
(4.1)
Let $`\chi (\lambda )`$ be a continuous function that is equal to 1 if $`|\lambda |<A`$ and is equal to $`0`$ if $`|\lambda |>A+1`$ for some $`A>0.`$ Then
$$|s(z)s_n(z)|\left|\frac{\chi (\lambda )m(d\lambda )}{\lambda z}\frac{\chi (\lambda )m_n(d\lambda )}{\lambda z}\right|+\frac{2}{\mathrm{min}\{\mathrm{dist}\{z,\pm A\}\}}.$$
According to (4.1) the first term in the r.h.s. of this inequality converges in probability to zero. Since $`A`$ is arbitrary, we obtain the required assertion.
To prove sufficiency we assume that for any $`zK`$
$$\underset{n\mathrm{}}{lim}𝐏\{|s(z)s_n(z)|>\epsilon \}=0.$$
(4.2)
This relation and the inequality (cf. (2.12))
$$|s_n(z)|\underset{zK}{\mathrm{max}}|\mathrm{Im}z|^1y_0^1<\mathrm{}$$
(4.3)
imply that
$$\underset{n\mathrm{}}{lim}𝐄\{|s(z)s_n(z)|\}=0,$$
(4.4)
i.e. the sequence $`\{s_n(z)\}`$ converges to zero in mean. We have also the inequality
$$|s_n^{^{}}(z)|y_0^2<\mathrm{}.$$
(4.5)
Inequalities (4.3) and (4.5) imply that the sequence $`\{s_n\}_{n=1}^{\mathrm{}}`$ of random analytic functions is uniformly bounded and equicontinuous. Thus, for any $`\eta >0`$ we can construct in $`K`$ a finite $`\eta `$-network, i.e. a set $`\{z_l\}_{l=1}^{p(\eta )}`$ such that for any $`zK`$ there exists $`z_l`$ satisfying the inequality $`|zz_l|\eta `$. Then we have for $`\varphi _n(z)s_n(z)s(z)`$, $`S_l=\{z:|zz_l|\eta \}`$, and $`\eta =y_0^2\epsilon /2,`$ where $`\epsilon `$ is arbitrary
$$\underset{K}{sup}|\varphi _n(z)|=\underset{l=1\mathrm{}p(\eta )}{\mathrm{max}}\underset{zKS_l}{sup}|\varphi _n(z)|\epsilon +\underset{l=1}{\overset{p(\eta )}{}}|\varphi _n(z_l)|,$$
and hence
$$𝐄\{\underset{K}{sup}|\varphi _n(z)|\}\epsilon +\underset{l=1}{\overset{p(\eta )}{}}𝐄\{|\varphi _n(z_l)|\}.$$
This inequality and (4.4) imply that
$$\underset{n\mathrm{}}{lim}𝐄\{\underset{zK}{sup}|s(z)s_n(z)|\}=0.$$
(4.6)
Assume now that the statement is false, i.e. the sequence $`\{m_n\}`$ does not converges weakly in probability to $`m`$. It means that there exists a continuous function $`\phi `$ of a compact support, a subsequence $`\{n_k\}`$ and some $`\epsilon >0`$ such that
$$\underset{n_k\mathrm{}}{lim}𝐏\left\{\left|\phi (\lambda )m(d\lambda )\phi (\lambda )m_{n_k}(d\lambda )\right|\epsilon \right\}=\xi >0.$$
(4.7)
On the other hand, we have from (4.6) and the Tchebyshev inequality that for any $`r`$ there exists an integer $`n(r)`$ such that for $`nn(r)`$
$$𝐏\left\{\underset{zK}{sup}|\varphi _n(z)|r^1\right\}1\xi /2.$$
(4.8)
Hence, one can select from the sequence $`\{n_k\}`$ a subsequence $`\{n_k^{}\}`$ such that inequalities (4.7) and (4.8) are both satisfied. Denote by $`𝒜`$ and by $``$ the events whose probabilities are written in the l.h.s. of (4.7) and (4.8). Then $`𝐏\{𝒜\}`$ $`𝐏\{𝒜\}+`$ $`𝐏\{\}1\xi /2`$. Hence, for any $`n_k^{}`$ there exists a realization $`\omega _{n_k^{}}`$ belonging to the both sets $`𝒜`$ and $``$, i.e. for which the both inequalities
$$\left|\phi (\lambda )m(d\lambda )\phi (\lambda )m_{n_k^{}}(d\lambda )\right|\epsilon ,\underset{zK}{sup}|\varphi _{n_k^{}}(z)|r^1$$
(4.9)
are valid. In view of the compactness of the family of the random analytic functions $`\{s_n\}`$ with respect to the uniform in $`K`$ convergence and the weak compactness of the family of random measure $`\{m_n\}`$ there exists a subsequence $`\{n_k^{^{\prime \prime }}\}`$ of $`\{n_k^{^{}}\}`$ and a subsequence of realizations $`\{\omega _{n_k^{\prime \prime }}\}`$ such that the subsequence $`\{m_{n_k^{\prime \prime }}\}`$ corresponding to these realizations converges weakly to a certain measure $`\stackrel{~}{m}`$ and we have in view of (4.7)
$$\left|\phi (\lambda )m(d\lambda )\phi (\lambda )\stackrel{~}{m}(d\lambda )\right|\epsilon >0.$$
(4.10)
On the other hand, in view of (4.9) and the continuity of the correspondence between measures and their Stieltjes transforms (see Proposition 2.1(v)), the subsequence $`\{s_{n_k^{\prime \prime }}\}`$ converges uniformly on $`K`$ to $`s(z)`$, the Stieltjes transform of the measure $`m`$. This is incompatible with (4.10), because of the one-to-one correspondence between measures and their Stieltjes transforms. $`\mathrm{}`$
###### Remark 1
Since the Stieltjes transforms of non-negative and normalized to unity measures are analytic and bounded for non-real $`z`$, we can replace the requirement of their convergence for any $`z`$ belonging to a certain compact of $`_\pm `$ by the convergence for any $`z`$ belonging to any interval of the imaginary axis, i.e. for $`z=iy,`$ $`y[y_1,y_2],`$ $`y_1>0.`$
###### Remark 2
The arguments, used in the proof of the proposition prove also that if $`\{m_n\}`$ is a sequence of random non-negative measures converging weakly in probability to a nonrandom non-negative measure $`m`$, then the Stieltjes transforms $`s_n`$ of $`m_n`$ and the Stieltjes transform $`s`$ of $`m`$ are related as follows
$$\underset{n\mathrm{}}{lim}𝐄\{\underset{zK}{sup}|s_n(z)s(z)|\}=0$$
(4.11)
for any compact $`K`$ of $`_\pm `$.
###### Lemma 4.1
Let $`H_n`$ be the random $`n\times n`$ matrix of the form (2.1) in which $`A_n`$ and $`B_n`$ are random Hermitian matrices, $`U_n`$ and $`V_n`$ are random unitary matrices distributed each according to the normalized to unity Haar measure on $`𝐔(n)`$ and $`A_n`$, $`B_n`$, $`U_n`$ and $`V_n`$ are mutually independent. Assume that the normalized counting measures $`N_{r,n},r=1,2`$ of matrices $`A_n`$ and $`B_n`$ converge in probability as $`n\mathrm{}`$ to non-random non-negative unit measures $`N_r,r=1,2`$ respectively and that
$$\underset{n}{sup}A_nT<\mathrm{},\underset{n}{sup}B_nT<\mathrm{}.$$
(4.12)
Then the normalized counting measure of $`H_n`$ converges in probability to a non-random unity measure $`N`$ whose Stieltjes transform $`f(z)`$ is a unique solution of system (2.18) in the class of functions $`f(z)`$, $`\mathrm{\Delta }_r(z),r=1,2`$ analytic for $`\mathrm{Im}z0`$ and satisfying conditions (2.12) - (2.14) and (2.19)
Proof. In view of Proposition 4.1 it suffices to show that $`lim_n\mathrm{}𝐄\{|g_n(z)f(z)|\}=0`$ for any $`z`$ belonging to a certain compact of $`_\pm `$. Moreover, according to Remark 1 after Proposition 4.1, we can restrict ourselves to a certain interval of the imaginary axis, i.e. to
$$z=iy,y[y_1,y_2],0<y_1<y_2<\mathrm{}.$$
(4.13)
Since the condition (4.12) of the lemma implies evidently the condition (3.1) of Theorem 3.1 and Theorem 3.2, all the results obtained in these theorems are valid in our case for any fixed realization of random matrices $`A_n`$ and $`B_n`$. In addition, all $`n`$-independent estimating quantities entering various bounds in the proofs of these theorems and depending on the forth moment $`m_4`$ in (3.1) and on $`y_0`$ will depend now on $`T`$ and on $`y_1`$ and $`y_2`$ in (4.13), but not on particular realizations of random matrices $`A_n`$ and $`B_n.`$ We will denote below all these quantities simply by the unique symbol $`C`$ that may have different value in different formulas.
In particular, denoting as above by $`\mathrm{}`$ the expectation with respect to the Haar measure and using (3.42), we can write that
$$𝐄\{|g_n(z)g_n(z)|\}𝐄\{|v_1^{1/2}(z)|\}\frac{C}{n}.$$
Thus, it suffices to show that
$$\underset{n\mathrm{}}{lim}𝐄\{|g_n(z)f(z)|\}=0,z=iy,y[y_1,y_2],$$
(4.14)
where $`y_1`$ is big enough. Introduce the quantities
$$\gamma _n(y)=iy(g_n(iy)f(iy)),\gamma _{r,n}(y)=\delta _{r,n}(iy)\mathrm{\Delta }_r(iy),r=1,2.$$
(4.15)
By using the second equation of system (2.18) we can write the identity
$$\gamma _n(y)=iy[f_2(iyt_{1,n}(y))f_2(iyt_1(y))]+\epsilon _{1,n}(y),$$
(4.16)
where
$$\epsilon _{1,n}(y)=iy[g_n(iy)f_2(iyt_{1,n}(y))],$$
(4.17)
$$t_{1,n}(y)=\frac{\delta _{1,n}(iy)}{g_n(iy)},t_1(y)=\frac{\mathrm{\Delta }_1(iy)}{f(iy)}.$$
(4.18)
We have
$$\begin{array}{c}𝐄\{|\epsilon _{1,n}(y)|\}y_2𝐄\{|g_n(iy)g_{2,n}(iyt_{1,n}(y))|\}+\\ 𝐄\{|g_{2,n}(iyt_{1,n}(y))f_2(iyt_{1,n}(y))|\}.\end{array}$$
(4.19)
The analogues of (3.38) - (3.39) in our case are:
$$g_n(z)=g_{2,n}(z\delta _{1,n}(z)g_n(z)^1)+\widehat{r}_{1,n}(z),$$
(4.20)
where
$$g_{2,n}(z)=n^1\mathrm{Tr}G_2(z)=\frac{N_{2,n}(\mathrm{d}\lambda )}{\lambda z},$$
is the Stieltjes transform of random NCM $`N_{2,n}`$ of $`H_{2,n}`$,
$$\begin{array}{c}\widehat{r}_{1,n}(z)=g_n^{}(z)n^1\mathrm{Tr}P^1g_n(z)^1G(z)\\ \delta _{1,n}^{}(z)n^1\mathrm{Tr}P^1g_n(z)^1G_2(z)G(z),\end{array}$$
the symbol $`\mathrm{}`$ denotes the expectation with respect the Haar measure on $`U(n),`$ $`P=1G_2(z)t_{1,n}(z),`$ and
$$g_n^{}(z)=g_n(z)g_n(z),\delta _{1,n}^{}(z)=\delta _{1,n}(z)\delta _{1,n}(z)$$
(4.21)
are the respective random variables centralized by the partial expectations with respect to the Haar measure. In addition, we have the analogue of (3.43)
$$\left|\widehat{r}_{1,n}(z)\right|\frac{C}{n}.$$
This leads to the following bound for the first term in the r.h.s. of (4.19):
$$𝐄\{|g_n(iy)g_{1,n}(iyt_{2,n}(y)|\}𝐄\{|\widehat{r}_{1,n}(iy)|\}\frac{C}{n}.$$
To show that the second term also vanishes as $`n\mathrm{}`$, we use the analogues of (3.32) and (3.36)
$$\left|g_{1,n}(iy)+\frac{1}{iy}\right|\frac{T}{y^2},|\delta _{2,n}(iy)|\frac{T}{y},$$
which imply that
$$|t_{1,n}(y)|2T,$$
(4.22)
if $`y_1`$ is big enough. Thus
$$𝐄\{|g_{2,n}(iyt_{1,n}(y))f_2(iyt_{1,n}(y))|\}\underset{|\zeta |T}{sup}𝐄\{|g_{2,n}(iy+\zeta )f_1(iy+\zeta )|\}.$$
The r.h.s of this inequality tends to zero as $`n\mathrm{}`$ in view of the hypothesis of Theorem 2.1 and Remark 2 after Proposition 4.1. Thus, there exist $`0<y_1<y_2<\mathrm{}`$ such that for all $`y[y_1,y_2],`$ $`lim_n\mathrm{}𝐄\{|\epsilon _{1,n}(y)|\}=0.`$ Analogous arguments show that $`lim_n\mathrm{}𝐄\{|\epsilon _{2,n}(y)|\}=0`$, where $`\epsilon _{2,n}(y)`$ is defined in (4.17) and in (4.18) where the indices 1 and 2 are interchanged. Thus we have
$$\underset{n\mathrm{}}{lim}𝐄\{|\epsilon _{r,n}(y)|\}=0,r=1,2.$$
(4.23)
Consider now the first term in the l.h.s. of (4.16). In view of (2.5) we can write this term in the form
$$[f_2(iyt_{1,n}(y))f_2(iyt_1(y))]=\frac{\delta _{1,n}}{fg_n}I_2\gamma _n+\frac{iy}{f}I_2\gamma _{1,n}=a_1\gamma _n+b_1\gamma _{1,n},$$
(4.24)
where $`I_{2,}a_1`$ and $`b_1`$ are defined by formulas (3.15) and (3.16), in which we have to replace $`\mathrm{\Delta }_1^{},`$ $`\mathrm{\Delta }_1^{^{\prime \prime }},`$ $`f^{}`$ and $`f^{\prime \prime }`$ by $`\mathrm{\Delta }_1,`$ $`\delta _{1,n},`$ $`f`$ and $`g_n`$ respectively. Denote by $`\mathrm{\Phi }=\{\mathrm{\Phi }_{ij}\}_{i,j=1}^3`$ the matrix defined by the l.h.s. of system (3.14) and by $`\mathrm{\Gamma }=\{\mathrm{\Gamma }_i\}_{i=1}^3`$ the vector with components $`\mathrm{\Gamma }_1=\gamma _n,`$ $`\mathrm{\Gamma }_2=\gamma _{1,n},`$ $`\mathrm{\Gamma }_3=\gamma _{2,n}.`$ Then we have from (4.16), (4.23) and (4.24)
$$𝐄\{|(\mathrm{\Phi }\mathrm{\Gamma })_1|\}𝐄\{|\epsilon _{1,n}|\}.$$
(4.25)
Interchanging in the above arguments indices $`1`$ and $`2`$ we obtain also that
$$𝐄\{|(\mathrm{\Phi }\mathrm{\Gamma })_2|\}𝐄\{|\epsilon _{2,n}|\}.$$
(4.26)
Besides, applying to the identity $`G(z)(H_1+H_2z)=1`$ the operation $`n^1Tr\mathrm{}`$ and subtracting from the result the third equation of system (2.18), we obtain the one more relation
$$𝐄\{|(\mathrm{\Phi }\mathrm{\Gamma })_3|\}=0.$$
(4.27)
It follows from the proof of Proposition 3.3 that the matrix $`\mathrm{\Phi }`$ is invertible if $`y_1`$ is big enough. Denote by $`\mathrm{}_1`$ the $`l^1`$-norm of $`^3`$ and by $`\mathrm{}`$ the induced matrix norm. Then we have
$$𝐄\{\mathrm{\Gamma }_1\}𝐄\{\mathrm{\Phi }^1\mathrm{\Phi }\mathrm{\Gamma }_1\}𝐄^{1/2}\{\mathrm{\Phi }^1^2\}𝐄^{1/2}\{\mathrm{\Phi }\mathrm{\Gamma }_1^2\}.$$
(4.28)
It follows from our arguments above that all entries of the matrices $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^1`$ and all components of the vector $`\mathrm{\Gamma }`$ are bounded uniformly in $`n`$ and in realizations of random matrices $`A_n`$,$`B_n,U_n`$ and $`V_n`$ in (2.1). Thus we have
$$\mathrm{\Phi }^1\underset{i,j=1}{\overset{3}{}}|(\mathrm{\Phi }^1)_{ij}|C,\mathrm{\Phi }\mathrm{\Gamma }_1\underset{i,j=1}{\overset{3}{}}|\mathrm{\Phi }_{ij}||\mathrm{\Gamma }|_jC.$$
These bounds and (4.25) - (4.28) imply that
$$𝐄\{\mathrm{\Gamma }_1\}C^{3/2}(𝐄\{|\epsilon _{2,n}|\}+𝐄\{|\epsilon _{2,n}|\})^{1/2}.$$
In view of (4.23) this inequality imply (4.14), i.e. the assertion of the lemma.$`\mathrm{}`$
Now we extend the result of Lemma 4.1 for the case of unbounded $`A_n`$ and $`B_n`$, having the limiting NCM’s with the finite first moments. We will apply the standard in probability truncation technique, whose random matrix version was used already in .
Proof of Theorem 2.1. Without loss of generality we can assume that
$$\underset{n}{sup}|\lambda |𝐄\{N_{1,n}(\mathrm{d}\lambda )\}m_1<\mathrm{}.$$
(4.29)
For any $`T>0`$ introduce the matrices $`A_n^T`$ and $`B_n^T`$ replacing eigenvalues $`A_n`$ and $`B_n`$ lying in $`]T,\mathrm{}[`$ by $`T`$ and eigenvalues lying in $`]\mathrm{},T]`$ by $`T`$. Denote by $`N_{r,n}^T,r=1,2`$ the NCM of $`A_n^T`$ and $`B_n^T`$. It is clear that for any $`T>0`$ and $`r=1,2`$, the sequence $`\{N_{r,n}^T\}_{n1}`$ converge weakly in probability to the measures $`N_r^T`$ as $`n\mathrm{}`$, where $`N_r^T`$ are analogously defined via $`N_r`$ and have their supports in $`[T,T]`$, and that for each $`r=1,2`$ the sequence $`\{N_r^T\}_{T1}`$ converge weakly to $`N_r`$ as $`T\mathrm{}`$. Denote by $`N_n^T,r=1,2`$ the NCM of $`H_n^T=H_{1,n}^T+H_{2,n}^T=V_n^{}A_n^TV_n+U_n^{}B_n^TU_n`$. According to linear algebra, if $`M_r,r=1,2`$ are two Hermitian $`n\times n`$ matrices, then
$$\mathrm{rank}(M_1+M_2)\mathrm{rank}M_1+\mathrm{rank}M_2,$$
(4.30)
and if $`\{\mu _{r,l}\}_{l=1}^n,r=1,2`$ are eigenvalues of $`M_r,r=1,2`$, then for any Borel set $`\mathrm{\Delta }`$
$$|\mathrm{\#}\{\mu _{1,l}\mathrm{\Delta }\}\mathrm{\#}\{\mu _{2,l}\mathrm{\Delta }\}|\mathrm{rank}(M_1M_2).$$
By using these facts we find that
$$\begin{array}{c}|N_n(\mathrm{\Delta })N_n^T(\mathrm{\Delta })|\frac{1}{n}\mathrm{rank}(H_nH_n^T)\frac{1}{n}\mathrm{rank}(A_nA_n^T)+\\ +\frac{1}{n}\mathrm{rank}(B_nB_n^T)N_{1,n}(\backslash ]T,T[)+N_{2,n}(\backslash ]T,T[),\end{array}$$
(4.31)
valid for any Borel set $`\mathrm{\Delta }`$. As a result, the Stieltjes transform $`g_n^T`$ of $`N_n^T`$ and the Stieltjes transform $`g_n`$ of $`N_n`$ are related as follows:
$$|g_n^T(z)g_n(z)|\frac{\pi }{|\mathrm{Im}z|}(N_{1,n}(\backslash ]T,T[)+N_{2,n}(\backslash ]T,T[)),$$
hence
$$𝐄\{|g_n^T(z)g_n(z)|\}\frac{\pi }{|\mathrm{Im}z|}(𝐄\{N_{1,n}(\backslash ]T,T[)\}+𝐄\{N_{2,n}(\backslash ]T,T[)\}).$$
(4.32)
and
$$\underset{n\mathrm{}}{lim}𝐄\{N_{r,n}(]T,T[)\}1N_r(]T,T[)=o(1),T\mathrm{}.$$
Since the norms of matrices $`H_1^T`$ and $`H_2^T`$ are bounded, the results of the Lemma 4.1 are applicable to the function $`g_n^T(z)`$, so that, in particular, for any non-real $`z`$ it converges in probability as $`n\mathrm{}`$ to a function $`f^T(z)`$ satisfying the system
$$\begin{array}{cc}\hfill f^T(z)& =f_1^T\left(z\frac{\mathrm{\Delta }_2^T(z)}{f^T(z)}\right)\hfill \\ \hfill f^T(z)& =f_2^T\left(z\frac{\mathrm{\Delta }_1^T(z)}{f^T(z)}\right)\hfill \\ \hfill f^T(z)& =\frac{1\mathrm{\Delta }_1^T(z)\mathrm{\Delta }_2^T(z)}{z}\hfill \end{array}.$$
In addition, since $`𝐄\{g_n^T(z)\}`$ and $`𝐄\{\delta _{1,n}^T(z)\}`$ are bounded uniformly in $`n`$ and $`T`$ for $`zE(y_0):`$
$$\begin{array}{c}|𝐄\{g_n^T(z)\}|\frac{1}{y_0},\hfill \\ |𝐄\{\delta _{1,n}^T(z)\}|\frac{1}{y_0}|\lambda |𝐄\{N_{1,n}^T(\mathrm{d}\lambda )\}\frac{1}{y_0}|\lambda |𝐄\{N_{1,n}(\mathrm{d}\lambda )\}\frac{m_1}{y_0},\hfill \end{array}$$
we have
$$|f^T(z)|\frac{1}{y_0},|\mathrm{\Delta }_1^T(z)|\frac{m_1}{y_0}$$
(4.33)
Thus, there exists a sequence $`T_k\mathrm{}`$ such that sequences of analytic functions $`\{f^{T_k}(z)\}`$ and $`\{\mathrm{\Delta }_1^{T_k}(z)\}`$ converge uniformly on any compact of the $`E(y_0)`$ of (4.32). In addition, the measures $`N_r^{T_k},r=1,2`$ converge weakly to the limiting measures $`N_r,r=1,2`$. Hence, there exist three analytic functions $`f(z)`$, $`\mathrm{\Delta }_1(z)`$ and $`\mathrm{\Delta }_2(z)=zf(z)+1\mathrm{\Delta }_1(z)`$ verifying (2.18). Besides, because of (4.33) and (3.1) for $`zE(y_0)`$ we have
$$|\mathrm{\Delta }_1(z)|\frac{m_1}{y_0},\mathrm{and}\mathrm{\Delta }_2(z)=o(1)\mathrm{as}y_0\mathrm{}.$$
As a result of relations above, $`f(z)`$ and $`\mathrm{\Delta }_r(z),r=1,2`$ satisfy the conditions of Proposition 3.3, hence they are defined uniquely.
Furthermore, we have
$$𝐄\{|g_n(z)f(z)|\}𝐄\{|g_n(z)g_n^{T_k}(z)|\}+𝐄\{|g_n^{T_k}(z)f^{T_k}(z)|\}+|f^{T_k}(z)f(z)|.$$
Hence in view of (4.32), arguments above on convergence of $`f^{T_k}`$ to $`f`$ , and Lemma 4.1 we conclude that for each $`zE(y_0)`$
$$\underset{n\mathrm{}}{lim}𝐄\{|g_n(z)f(z)|\}=0.$$
In view of Proposition 4.1 this implies that the NCM (2.2) of random matrices (2.1) converges weakly in probability as $`n\mathrm{}`$ to the non-random measure, whose Stieltjes transform is a unique solution of system (2.18).$`\mathrm{}`$
## 5 Properties of the Solution
Here we will consider several simple properties of the limiting eigenvalue counting measure described by Theorem 2.1, i.e. the measure, whose Stieltjes transform is a solution of (2.18) satisfying (2.12)–(2.14). We refer the reader to works and references therein for a rather complete collection of results on properties of the measure, resulting from the binary operation in the space of the probability measures, defined by a version of system (2.18). This binary operation is called free additive convoluton.
(i) Assume that the supports of the limiting eigenvalue measures of the matrices $`A_n`$ and $`B_n`$ are bounded, i.e. there exist $`\mathrm{}<a_{r,}b_r<\mathrm{},r=1,2,`$ such that
$$\mathrm{supp}N_r[a_r,b_r],r=1,2.$$
(5.1)
Then
$$\mathrm{supp}N[a_1+a_2,b_1+b_2].$$
(5.2)
Proof. Denote by $`\{\lambda _l\}_{l=1}^n`$ and by $`\{\lambda _{r,l}\}_{l=1}^n,r=1,2`$ eigenvalues of $`H_n`$ and $`H_{r,n}`$ in (2.1) respectively. Then, according to the linear algebra (cf.(4.31)),
$$\mathrm{\#}\{\lambda _l\backslash [a_1+a_2,b_1+b_2]\}\mathrm{\#}\{\lambda _{1,l}\backslash [a_1,b_1]\}+\mathrm{\#}\{\lambda _{2,l}\backslash [a_2,b_2]\}.$$
In view of Theorem 2.1 and (5.1) this leads to the relation $`N(\mathrm{}\sigma )=0`$, i.e. to (5.2).
(ii). Examples. 1. Consider the case when $`A_n=B_n`$, i.e. $`N_1=N_2`$. In this case system (2.18) will have the form
$$f(z)=f_1\left(\frac{z}{2}\frac{2}{f(z)}\right).$$
(5.3)
Take $`N_1=N=\alpha `$ $`\delta _0+(1\alpha )`$ $`\delta _a`$ where $`0\alpha 1`$, $`a>0`$ and $`\delta _\lambda `$ is the unit measure concentrated at $`\lambda `$. Then
$$f_1(z)=\frac{\alpha }{z}+\frac{1\alpha }{az}$$
and (2.18) reduces to the quadratic equation
$$z(z2a)f^2+2a(12\alpha )f1=0,$$
whose solution satisfying (2.12) - (2.14) is
$$f(z)=\frac{a(12\alpha )\sqrt{(z\lambda _+)(z\lambda _{})}}{z(z2a)},\lambda _\pm =a(1\pm 2\sqrt{\alpha (1\alpha )}).$$
By using (2.15) we find that the limiting measure in this case has the form
$$N=(2\alpha 1)_+\delta _0+(12\alpha )_+\delta _{2a}+N^{},$$
(5.4)
where $`x_+=\mathrm{max}(0,x),`$ and
$$N^{}(\mathrm{d}\lambda )=\frac{1}{\pi }\frac{\sqrt{(\lambda _+\lambda )(\lambda \lambda _{})}}{\lambda (\lambda 2a)}\chi _{[0,2a]}(\lambda )\mathrm{d}\lambda $$
(5.5)
is the absolute continuous measure of the mass $`12\alpha .`$ Here $`\chi _\mathrm{\Delta }(\lambda )`$ is the indicator of the set $`\mathrm{\Delta }`$. In the cases $`\alpha =0,1`$ (5.4) is $`\delta _{2a}`$ and $`\delta _0`$ respectively, and in the case $`\alpha =1/2`$ (5.4) has no atoms, but only the square root singularities
$$N^{}(\mathrm{d}\lambda )=\frac{1}{\pi \sqrt{\lambda (2a\lambda )}}\chi _{[0,2a]}(\lambda )\mathrm{d}\lambda $$
(5.6)
Formulas (5.3)–(5.6) shows that:
* the result (5.2) is optimal with respect to the endpoints of the measures $`N_r,r=1,2`$ and $`N`$;
* in the case when $`N_1=N_2`$ have atoms of the mass $`\mu >1/2`$ at the same point then the measure $`N`$ has also an atom of the mass $`(2\mu 1)`$ (for general results of this type see ).
However, in general the support of $`N`$ is strictly included in the sum of supports of measures $`N_r,r=1,2`$, i.e. the inclusion in the r.h.s part of (5.3) is strict. This can be illustrated by the following two examples.
2. Take again $`N_1=N_2`$, where now
$$N_1(\mathrm{d}\lambda )=\frac{1}{\pi \sqrt{a^2\lambda ^2)}}\chi _{[a,a]}(\lambda )\mathrm{d}\lambda .$$
is the arcsin law. This measure corresponds to the matrix ensemble (2.37) with
$$V(\lambda )=\{\begin{array}{ccc}0,\hfill & |\lambda |<1,\hfill & \\ \mathrm{},\hfill & |\lambda |>1.\hfill & \end{array}$$
(5.7)
In this case equation (5.3) is again quadratic and leads to
$$N(\mathrm{d}\lambda )=\frac{\sqrt{3a^2\lambda ^2)}}{\pi (4a^2\lambda ^2)}\chi _{[\sqrt{3}a,\sqrt{3}a]}(\lambda )\mathrm{d}\lambda .$$
3. In the next example we take
$$N_r(\mathrm{d}\lambda )=\frac{1}{4\pi a_r^2}\sqrt{8a_r^2\lambda ^2}\chi _{[2\sqrt{2}a_r,2,\sqrt{2}a_r]}(\lambda )\mathrm{d}\lambda ,r=1,2,$$
i.e. the both measures are the semicircle laws (2.31). Then it is easy to find that $`N`$ is also the semicircle measure with the parameter $`a^2=`$ $`a_1^2+a_2^2.`$ This case was indicated in . It can be easily deduced from the law of addition of the R-transforms of Voiculescu , because in this case $`R_r(f)=2a_r^2f`$. For further properties of the measure $`N`$ in the case when one of $`N_r,r=1,2`$ is the semicircle law see .
(iii). Suppose that one of the measures $`N_r(\mathrm{d}\lambda ),r=1,2`$ is absolute continuous with respect to the Lebesgue measure, i.e., say, $`N_1(\mathrm{d}\lambda )=\rho _1(\lambda )\mathrm{d}\lambda ,`$ and
$$\overline{\rho }_1=\mathrm{ess}\underset{\lambda }{sup}|\rho _\mathit{1}(\lambda )|<\mathrm{},.$$
Then $`N`$ is also absolute continuous with respect to the Lebesgue measure, i.e.$`N(\mathrm{d}\lambda )=\rho (\lambda )\mathrm{d}\lambda ,`$ and
$$\mathrm{ess}\underset{\lambda }{sup}|\rho _1(\lambda )|=\overline{\rho }_1<\mathrm{}.$$
(5.8)
Proof. Indeed, since the function $`z_1^{}=z\mathrm{\Delta }_{2,1}/f(z)`$ is analytic for non-real $`z`$, the number of its zeros in any compact of $`\backslash `$ is finite. Thus, for any $`\lambda `$ there exists a sequence $`\{z_n\}`$ of non-real numbers such that $`z_n\lambda `$ as $`n\mathrm{}`$ and $`\mathrm{Im}`$ $`z_n^{}0.`$ Hence, we have from the first equation of system (2.18) for $`z_n^{}=`$ $`\lambda _n^{}+i\epsilon _n^{}`$
$$\frac{1}{\pi }\mathrm{Im}f(z)=\frac{1}{\pi }\frac{\epsilon _r^{}\rho _r(\mu )\mathrm{d}\mu }{(\mu \lambda _r^{})^2+(\epsilon _r^{})^2}\overline{\rho }_1\frac{1}{\pi }\frac{\epsilon _r^{}\mathrm{d}\mu }{(\mu \lambda _r^{})^2+(\epsilon _r^{})^2}=\overline{\rho }_1.$$
This relation and the inversion formula (2.15) yield (5.8). For more general results in this direction see the recent paper .
## 6 Discussion
In this Section we comment on several topics related to those studied above.
1. In this paper we deal with Hermitian and unitary matrices, i.e. we assume that the matrices $`A_n`$ and $`B_n`$ in (2.1) are Hermitian and $`U_n`$ and $`V_n`$ are unitary. It is natural also to consider the case of real symmetric $`A_n`$ and $`B_n`$ and orthogonal $`U_n`$ and $`V_n`$. This case can be handled by using the analogue of formula (3.11) of the orthogonal group $`O(n)`$. Indeed, it is easy to see that this analogue has the form
$$\underset{O(n)}{}\mathrm{\Phi }^{^{}}(O^TMO)[X,O^TMO]dO=0,$$
where $`O^T`$ is the transposed to $`O`$ and $`X`$ is a real symmetric matrix. By using this formula we obtain instead of (3.23)
$$G_{aa}(H_2G)_{bc}+G_{ab}(H_2G)_{ac}=(GH_2)_{aa}G_{bc}+(GH_2)_{ab}G_{bc}.$$
The second terms in both sides of this formula give two additional terms
$$n^1G^TH_2G+n^1H_2G^TG.$$
in the analogue of (3.40). These terms, however, produce the asymptotically vanishing contribution because, in view of (3.3), (3.6) and (3.37), we have
$$\left|n^2\mathrm{Tr}(1+\mathrm{\Delta }_{2,n}f_n^1G_1)^1G_1(G^TH_2G+H_2G^TG)\right|\frac{2}{ny_0^3}m_4^{1/4}.$$
Similar and also negligible as $`n\mathrm{}`$ terms appear in analogues of formulas (3.53), (3.69) and (3.73) of the proof of Theorem 3.2. As the result, we obtain in this case the same system (2.18), defining the Stieltjes transform of the limiting eigenvalue counting measure of the analogue of (2.1) with the real symmetric $`A_n`$ and $`B_n`$ and orthogonal Haar-distributed $`U_n`$ and $`V_n`$.
2. As was mentioned in the Introduction, our main result, Theorem 2.1, can be viewed as an extension of the result by Speicher , obtained by the moment method and valid for uniformly in $`n`$ bounded matrices $`A_n`$ and $`B_n`$ in (2.1). Both results are analogues for randomly rotated matrices of old results of (see (2.24) and (2.33)) on the form of the limiting eigenvalue counting measure of the sum of an arbitrary matrix and certain random matrices (see (2.20) and (2.26)), in particular, Gaussian random matrices (2.28). In this case, however, there exists another model, proposed by Wegner that combines properties of random matrices, having all entries roughly of the same order, and of random operators, whose entries decay sufficiently fast in the distance from the principal diagonal (see e.g. ). A simple, but rather non-trivial version of the Wegner model corresponds to the selfadjoint operator $`H`$ acting in $`l^2(^d)\times 𝐂^n`$ and defined by the matrix
$$H(x,j;y,k)=v(xy)\delta _{jk}+\delta (xy)f_{jk}(x)$$
(6.1)
where $`x,y^d`$, $`j,k=1,\mathrm{},n`$, $`\delta (x)`$ is the $`d`$-dimensional Kronecker symbol,
$$v(x)=\overline{v}(x),\underset{x^d}{}|v(x)|<\mathrm{},$$
(6.2)
and $`f(x)=\{f_{jk}(x)\}_{j,k=1}^n`$, $`x^d`$ are independent for different $`x`$ and identically distributed $`n\times n`$ Hermitian matrices, whose distribution for any $`x`$ is given by (2.28). According to (see also ) asymptotic for $`n\mathrm{}`$ properties of operator (6.1) resemble, in many aspects, asymptotic properties of matrices (2.28). The ”free” analogue of the Wegner model was proposed in . In this case i.i.d. matrices $`f(x)`$ have the form
$$f(x)=U_n^{}(x)B_nU_n(x)$$
(6.3)
where $`B_n`$ is as in (2.1) and $`U_n(x)`$, $`x^d`$ are i.i.d. unitary $`n\times n`$ matrices whose distribution is given by the Haar measure on $`U(n)`$. By using a version of the moment method, similar to that of paper , or, rather, its formal scheme, the authors derived the limiting form of
$$𝐄\left\{n^1\underset{j=1}{\overset{n}{}}G(x,j;y,j)\right\},$$
where $`G(x,j;y,k)`$ is the matrix (the Green function) of the resolvent $`(Hz)^1`$ of (6.1) - (6.3). The authors also found a certain second moment of the Green function. This moment is necessary to compute the a.c. conductivity via the Kubo formula. Because of the moment method results of are valid for uniformly bounded in $`n`$ matrices $`B_n`$ in (6.3), similar to results for matrices (2.1) obtained in . By using a natural extension of the differentiation formula (3.11) and the technique developed in to analyze the Wegner model, the results of paper can be extended to the case of arbitrary matrices $`B_n`$ in (6.3), because in this case the role of condition (2.17) of Theorem 2.1) plays condition (6.2).
3. As was mentioned before asymptotic properties of random matrices are of considerable interest in the certain branches of the operator algebra theory and related branch of the non-commutative probability theory, known as free probability (see and references therein). Here large random matrices is an important example of the asymptotically free non-commutative random variables, providing a sufficiently reach analytic model of the abstract notion of freeness of elements of an operator algebra. The most widely used examples of asymptotically free families of non-commutative random variables are Gaussian random matrices and unitary Haar-distributed random matrices. The proof of asymptotic freeness of unitary matrices given in reduces to that for complex Gaussian matrices basing on the observation that the unitary part of the polar decomposition of complex Gaussian matrix with independent entries is the Haar-distributed unitary matrix. This method requires certain technicalities because of the singularity of the polar decomposition at zero. On the other hand, the differentiation formula (3.11) allows one to prove directly similar statements. Here is an example of results of this type (related results are proved in ).
###### Theorem 6.1
Let $`k`$ be a positive integer, $`\{T_{r,n}\}_{r=1}^k`$ be a set of $`n\times n`$matrices, such that
$$\underset{rk;k,l,n}{sup}n^1\mathrm{Tr}(T_{r,n}^{}T_{r,n})^l<\mathrm{},$$
(6.4)
and let $`U_n`$ be the unitary and Haar-distributed random matrix. If for any $`k`$
$$\underset{n\mathrm{}}{lim}n^1\mathrm{Tr}T_{r,n}=0,r=1,\mathrm{},k,$$
(6.5)
then for any set of non-zero integers such that $`\{m_r\}_{r=1}^k`$, $`_{r=1}^km_r=0`$
$$\underset{n\mathrm{}}{lim}n^1\mathrm{Tr}U_n^{m_1}T_{1,n}\mathrm{}U_n^{m_k}T_{k,n}=0,$$
(6.6)
where $``$ denotes the integration with respect to the Haar measure over $`U(n)`$.
###### Remark 1
The theorem is trivially true in the case when $`_{r=1}^km_r0`$.
In the two subsequent lemmas we omit the subindex $`n`$.
###### Lemma 6.1
Let $`\{T_i\}_{i=1}^k`$ be a set of $`n\times n`$ matrices and $`U`$ is the Haar-distributed unitary matrix. Then for any set of non-zero integers $`\{m_i\}_{i=1}^k`$, $`_{i=1}^km_i=0`$ the following identity holds:
$$n^1\mathrm{Tr}U^{m_1}T_1\mathrm{}U^{m_k}T_k=\underset{l_1=2}{\overset{m_1}{}}n^1\mathrm{Tr}U^{l_11}n^1\mathrm{Tr}(U^{m_1l_1+1}T_1\mathrm{}U^{m_k}T_k)$$
$$\underset{r\{2,\mathrm{},k\},m_r>0}{}\underset{l_r=1}{\overset{m_r}{}}n^1\mathrm{Tr}(U^{m_1}T_1\mathrm{}T_{r1}U^{l_r1})n^1\mathrm{Tr}(U^{m_rl_r+1}T_r\mathrm{}U^{m_k}T_k)+$$
(6.7)
$$\underset{r\{2,\mathrm{},k\},m_r<0}{}\underset{l_r=1}{\overset{m_r}{}}n^1\mathrm{Tr}(U^{m_1}T_1\mathrm{}T_{r1}U^{l_r})n^1\mathrm{Tr}(U^{m_r+l_r}T_r\mathrm{}U^{m_k}T_k).$$
Proof. Without loss of generality assume that $`m_1>0`$. Then, using the analogue of formula (3.11) for the average $`[U^{m_1}T_1\mathrm{}U^{m_k}T_k]_{ab}`$, we obtain for any Hermitian $`X`$
$$\underset{r\{1,\mathrm{},k\},m_r>0}{}\underset{l_r=1}{\overset{m_r}{}}[U^{m_1}T_1\mathrm{}T_{r1}U^{l_11}XU^{m_rl_r+1}T_r\mathrm{}U^{m_k}T_k]_{ab}+$$
$$\underset{r\{2,\mathrm{},k\},m_r<0}{}\underset{l_r=1}{\overset{m_r}{}}[U^{m_1}T_1\mathrm{}T_{r1}U^{l_r}XU^{m_r+l_r}T_r\mathrm{}U^{m_k}T_k]_{ab}=0$$
(6.8)
Choosing as $`X`$ the Hermitian matrix having only $`(c,d)`$-th and $`(d,c)`$-th non-zero entries, setting then $`a=c`$ and $`b=d`$ and applying to the result the operation $`n^2_{a,b}`$, we obtain (6.7).$`\mathrm{}`$
###### Lemma 6.2
Under the conditions (6.4) and (6.5) the variance $`D=|\xi ^{}|^2`$ of the random variable
$$\xi =n^\mathit{1}\mathrm{Tr}L,L=U^{m_\mathit{1}}T_\mathit{1}\mathrm{}U^{m_k}T_k$$
(6.9)
is of the order $`n^2`$ as $`n\mathrm{}`$.
Proof. Using the same technique as that in Lemma 6.1 for $`\overline{L_{ab}}L_{cd}`$ we obtain the relation
$$D=\underset{l_1=2}{\overset{m_1}{}}\overline{\xi }^{}n^1\mathrm{Tr}U^{l_11}n^1\mathrm{Tr}(U^{m_1l_1+1}T_1\mathrm{}U^{m_k}T_k)$$
$$\underset{r\{2,\mathrm{},k\},m_r>0}{}\underset{l_r=1}{\overset{m_r}{}}\overline{\xi }^{}n^1\mathrm{Tr}(U^{m_1}T_1\mathrm{}T_{r1}U^{l_r1})n^1\mathrm{Tr}(U^{m_rl_r+1}T_r\mathrm{}U^{m_k}T_k)+$$
(6.10)
$$\begin{array}{c}\underset{r\{2,\mathrm{},k\},m_r<0}{}\underset{l_r=1}{\overset{m_r}{}}\overline{\xi }^{}n^1\mathrm{Tr}(U^{m_1}T_1\mathrm{}T_{r1}U^{l_r})n^1\mathrm{Tr}(U^{m_r+l_r}T_r\mathrm{}U^{m_k}T_k)\\ +n^2\mathrm{\Phi },\end{array}$$
where
$$\mathrm{\Phi }=\underset{r\{1,\mathrm{},k\},m_r>0}{}\underset{l_r=1}{\overset{m_r}{}}n^1\mathrm{Tr}(U^{m_rl_r+1}T_r\mathrm{}T_kU^{m_1}T_1\mathrm{}T_{r1}U^{l_r1})^{}L+$$
$$+\underset{r\{2,\mathrm{},k\},m_r<0}{}\underset{l_r=1}{\overset{m_r}{}}n^1\mathrm{Tr}(U^{m_r+l_r}T_r\mathrm{}T_kU^{m_1}T_1\mathrm{}T_{r1}U^{l_r})^{}L.$$
We have obviously for $`k=m=1`$
$$n^1\overline{\mathrm{Tr}(UT)^{}}n^1\mathrm{Tr}(UT)\frac{1}{n^2}n^1\mathrm{Tr}(TT^{}).$$
We proceed further by induction. In view of condition (6.4) and Proposition 3.1 we have the bound
$$|n^1\mathrm{Tr}(U^{m_1}T_{r_1}\mathrm{}U^{m_p}T_{r_p})|C_,^2$$
(6.11)
where $`C`$ may depend only on $`p`$. Now, since $`n^1\mathrm{Tr}U^l=0,l0,`$ the summands of the first term in r.h.s. of (6.10) can be estimated as follows
$$\left|\overline{\xi }^{}n^1\mathrm{Tr}(U^{l_1})n^1\mathrm{Tr}(U^{m_1l_1+1}T_1\mathrm{}U^{m_k}T_k)\right|C\sqrt{D}\sqrt{|n^1\mathrm{Tr}(U^{l_1})^{}|^2}.$$
(6.12)
Likewise, by using the cyclic property of the trace, the identity $`a^{}bc=a^{}b^{}c+a^{}c^{}b`$, Schwarz inequality, and (6.11), we obtain for the second term in the right-hand side of (6.10) the following estimates for $`r2`$
$$\begin{array}{c}\left|\overline{\xi }^{}n^1\mathrm{Tr}(U^{m_1}T_1\mathrm{}T_{r1}U^{l_r1})n^1\mathrm{Tr}(U^{m_rl_r+1}T_r\mathrm{}U^{m_k}T_k)\right|\\ C\sqrt{D}\{\sqrt{|n^1\mathrm{Tr}(U^{m_1+l_r1}T_1\mathrm{}U^{m_{r1}}T_{r1})^{}|^2}+\\ +\sqrt{|n^1\mathrm{Tr}(U^{m_rl_r+1}T_r\mathrm{}U^{m_k}T_k)^{}|^2}\}\end{array}$$
(6.13)
The third term in the right-hand side of (6.10) can be estimated analogously. The forth term is of the order $`1/n^2`$ in view of (6.9). By the induction hypothesis the expectations under square roots in the r.h.s. of (6.13) and (6.12) are of the order $`n^2`$. This leads to the inequality
$$D\frac{C_1}{n}\sqrt{D}+\frac{C_2}{n^2},$$
where $`C_1`$ and $`C_2`$ are independent of $`n`$. This implies the bound $`D=O(n^2)`$.$`\mathrm{}`$
Proof of Theorem 6.1. We use Lemma 6.1 and again the induction. We have first
$$n^1\mathrm{Tr}U^mT_1U^mT_2=n^1\mathrm{Tr}T_1n^1\mathrm{Tr}T_2=0.$$
In general case we use Lemma 6.2 to factorize asymptotically the moments in the r.h.s. of (6.7). In the resulting relation the expressions $`n^1\mathrm{Tr}U^{m_{r_1}}T_{r_1}\mathrm{}U^{m_{r_s}}T_{r_s}`$ are zero for any collection $`(T_{r_1},\mathrm{},T_{r_s})`$ and any $`n`$, if $`_{i=1}^sm_{r_i}0,`$ and tend to zero as $`n\mathrm{}`$ if $`_{i=1}^sm_{r_i}=0`$ in view of the induction hypothesis and condition (6.5). This leads to (6.6).$`\mathrm{}`$
###### Remark 2
A simple version of the above arguments allows us to prove that the normalized counting measure of the Haar distributed unitary matrices converges with probability one to the uniform distribution on the unit circle. Indeed, consider again the Stieltjes transform $`g_n`$ of this measure, supported now on the unit circle. By the spectral theorem for unitary matrices we have
$$g_n(z)=n^1\mathrm{Tr}G(z),G(z)=(Uz)^1,|z|1.$$
(6.14)
We can then obtain the following identities
$$\mathrm{Tr}G^2(z)U=0,g_n(z)n^1\mathrm{Tr}G(u)U=0,$$
(6.15)
$$g_n(z_1)n^1\mathrm{Tr}G(z_1)Ug(z_2)+n^3\mathrm{Tr}G(z_1)G(z_2)UG(z_2)=0.$$
(6.16)
By using the obvious relations
$$G^{}(z)=G^2(z),G(0)=U^1,G(\mathrm{})=0,$$
we obtain from the first of identities (6.15)
$$f_n(z)g_n(z)=\{\begin{array}{ccc}0,\hfill & |z|<1,\hfill & \\ z^1,\hfill & |z|>1.\hfill & \end{array}.$$
This relation shows that the expectation of the normalized counting measure of $`U`$ is the uniform distribution on the unit circle, the fact that follows easily from the shift invariance of the Haar measure. Now the second identity (6.15) and (6.16) lead to the bound
$$|g_n(z)|^2\frac{C(r_0)}{n^2},|z|r_{0,}$$
where $`C(r_0)`$ is independent of $`n`$ and finite if $`r_0`$ is small enough. This bound and arguments analogous to those used in the proof of Theorem 3.1 imply that the normalized eigenvalue counting measure of unitary Haar distributed random matrices converges with probability one to the uniform distribution on the unit circle. This fact as well as the analogous fact for the orthogonal group can be deduced from the works by Dyson (see e.g. ), where the joint probability distribution of all $`n`$ eigenvalues of the Haar distributed unitary or orthogonal matrices was found and studied. This technique is more powerful but also more complex than that used above and based on rather elementary means.
Acknowledgements. V. Vasilchuk is thankful to Laboratoire de Physique Mathématique et Géométrie de l’Université Paris-7 for hospitality and to Ministère des Affaires Etrangères de France for the financial support.
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# Damping of condensate collective modes due to equilibration with the non-condensate
## I Introduction
The collective oscillations of a condensate at zero temperature $`T=0`$ are well described by the solutions of the linearized Gross-Pitaevskii (GP) time-dependent equation of motion for the condensate wavefunction $`\mathrm{\Phi }(𝐫,t)`$. At finite temperatures, the condensate dynamics is modified by interactions with the non-condensate atoms in the thermal cloud, which has the effect of renormalizing and damping the condensate oscillations. Recently the coupled dynamics of the condensate and the thermal cloud has been the subject of several theoretical studies . Such calculations lead to a generalized GP equation for $`\mathrm{\Phi }(𝐫,t)`$ and some appropriate Boltzmann-like kinetic equation describing the dynamics of the non-condensate atoms. In the present work, we make use of the recent formulation of Zaremba, Nikuni, and Griffin (ZNG) to discuss a new kind of damping of condensate oscillations that arises from collisions between the condensate and non-condensate components. In contrast to Ref. , which discussed the collision-dominated hydrodynamic regime, here we discuss the collisionless regime. Within the well-known Thomas-Fermi (TF) approximation, we derive a generalized Stringari wave equation describing the condensate normal modes that is valid at finite $`T`$ and includes damping due to the fact that the condensate is not in equilibrium with the thermal cloud. This new source of damping is in addition to the usual Landau and Beliaev damping considered in the collisionless region at finite $`T`$ .
Our theory can be used to generalize any discussion based on the usual GP equation at $`T=0`$. This simplicity is due to our neglect of any dynamics of the thermal cloud. Available studies of collective modes at finite $`T`$ suggest that, for any given mode symmetry, one mode mainly involves motion of the condensate (with a small out-of-phase motion of the non-condensate). This mode is a natural extension of the $`T=0`$ oscillation of a pure condensate and should be described by our theory. The other mode, of the same symmetry, mainly involves the motion of the thermal cloud (with a small in-phase motion of the condensate) and can be viewed as the natural extension of the oscillations above the critical temperature $`T_{\mathrm{BEC}}`$ . Our present calculations do not apply to such “normal fluid” oscillations, which include the Kohn mode at the trap frequency.
## II Derivation of Model
Our starting point is the finite $`T`$ generalized GP equation derived by ZNG (see also Refs. and )
$$i\mathrm{}\frac{\mathrm{\Phi }}{t}=\left[\frac{\mathrm{}^2}{2m}^2+U_{\mathrm{ext}}+gn_c+2g\stackrel{~}{n}i\mathrm{}R\right]\mathrm{\Phi },$$
(1)
where the interaction parameter $`g=4\pi \mathrm{}^2a/m`$, $`a`$ is the s-wave scattering length, $`n_c(𝐫,t)=|\mathrm{\Phi }(𝐫,t)|^2`$, and $`\stackrel{~}{n}(𝐫,t)`$ is the non-condensate local density. The damping term in (1) is given by $`R(𝐫,t)\mathrm{\Gamma }_{12}(𝐫,t)/2n_c(𝐫,t)`$, with
$$\mathrm{\Gamma }_{12}(𝐫,t)=\frac{d𝐩}{(2\pi \mathrm{})^3}C_{12}[f(𝐩,𝐫,t),\mathrm{\Phi }(𝐫,t)].$$
(2)
This involves the collision integral $`C_{12}[f,\mathrm{\Phi }]`$ describing collisions of condensate atoms with the thermal atoms, which also enters the approximate semi-classical kinetic equation for the single-particle distribution function (valid for $`k_\mathrm{B}Tgn_{c0}`$ and $`k_\mathrm{B}T\mathrm{}\omega _0`$)
$$\frac{f}{t}+\frac{𝐩}{m}\mathbf{}f\mathbf{}U\mathbf{}_𝐩f=C_{12}[f,\mathrm{\Phi }]+C_{22}[f,\mathrm{\Phi }].$$
(3)
Here the collision integral $`C_{22}[f]`$ describes binary collisions between non-condensate atoms. It does not change the number of condensate atoms and hence does *not* appear explicitly in (1). These coupled equations (1)-(3), along with equations (23a) and (23b) of Ref. defining the collision integrals $`C_{12}`$ and $`C_{22}`$, were derived in the semi-classical approximation. However, they are expected to contain all the essential physics in trapped Bose-condensed gases at finite $`T`$, in both the collisionless and hydrodynamic domains. They assume that the atoms in the thermal cloud are well-described by the single-particle Hartree-Fock spectrum $`\stackrel{~}{\epsilon }_p(𝐫,t)=p^2/2m+U(𝐫,t)`$, where $`U(𝐫,t)=U_{\mathrm{ext}}(𝐫)+2g[n_c(𝐫,t)+\stackrel{~}{n}(𝐫,t)]`$. We expect this semi-classical description to break down only for very low temperatures where the Bogoliubov excitation spectrum is more appropriate . Our entire discussion is within what is called the Popov approximation in that we have ignored all effects associated with the anomalous pair correlations $`\stackrel{~}{m}(𝐫,t)=\stackrel{~}{\psi }(𝐫,t)\stackrel{~}{\psi }(𝐫,t)`$.
The coupled equations (1)-(3) have been used to derive the generalized two-fluid hydrodynamic equations in the collision-dominated region described by a local-equilibrium Bose distribution . They also have been recently used to give a detailed analysis of condensate growth by quenching the thermal cloud distribution . In these papers, (1)-(3) are solved with both the condensate and non-condensate being treated dynamically and allowed to be out of equilibrium. The key limitation of the present paper is that we only consider the dynamics of the condensate, with the thermal cloud being in static equilibrium. This assumption allows a simple theoretical development and should be adequate for out-of-phase modes. The key role of the condensate and non-condensate being out of diffusive equilibrium was first stressed in a series of papers by Gardiner and coworkers . These were based on a kinetic master equation formalism quite different from what we use, and no application was made to the damping of condensate collective modes.
It is important to understand what is meant by the *collisionless regime* and to clarify how this terminology relates to the present work. Above $`T_{\mathrm{BEC}}`$ (where $`C_{12}=0`$), equation (3) reduces to the Boltzmann equation describing a normal gas . In this case, the collisionless region is well defined and corresponds to having $`\omega _i\tau _{\mathrm{cl}}1`$, where $`\omega _i`$ is the collective mode frequency of the gas, on the order of the trap frequency, and $`\tau _{\mathrm{cl}}`$ can be approximated by the mean time between collisions described by the classical Boltzmann collision integral $`C_{22}`$. In static equilibrium, this collision rate for a uniform gas is given by
$$\frac{1}{\tau _{\mathrm{cl}}}=\sqrt{2}n\sigma \overline{v},$$
(4)
where $`n`$ is the density of atoms, $`\sigma =8\pi a^2`$ is the quantum collision cross section, and $`\overline{v}`$ is the average speed of an atom in the gas. Both above and below $`T_{\mathrm{BEC}}`$, the analogous collision time corresponding to collision processes described by $`C_{22}`$ in (3) will give an estimate of the lifetime of a *single-particle* excitation in the thermal cloud. This is distinct from the physics given by $`C_{12}`$, which describes the collisions of condensate atoms with atoms from the thermal cloud. In particular, one finds that if the thermal cloud is described by the equilibrium Bose distribution $`(f=f^0)`$, then $`C_{22}[f^0,\mathrm{\Phi }]=0`$ but $`C_{12}[f^0,\mathrm{\Phi }]0`$. Thus $`C_{12}`$ will give rise to damping of condensate oscillations even when the thermal cloud is treated statically. Of course, in the collisionless region, there is another source of damping arising from the dynamical mean-field coupling between the condensate and thermal cloud that is also included in (1) and (3); this is Landau damping , which will be discussed below.
### A Static Popov approximation
In the present work, we use these equations to calculate the damped normal modes of the condensate given by the solutions of (1) assuming that the non-condensate atoms always remain in static thermal equilibrium. For our model, this means we take
$$f(𝐩,𝐫,t)f^0(𝐩,𝐫)=\frac{1}{e^{\beta [p^2/2m+U_0(𝐫)\stackrel{~}{\mu }_0]}1},$$
(5)
where $`\stackrel{~}{\mu }_0`$ is the equilibrium chemical potential of the non-condensate and $`U_0(𝐫)=U_{\mathrm{ext}}(𝐫)+2g[n_{c0}(𝐫)+\stackrel{~}{n}_0(𝐫)]`$. The detailed analysis given by ZNG shows that the Bose-Einstein distribution in (5) is a stationary solution to (3) when the condensate and non-condensate are in diffusive equilibrium, which requires $`\stackrel{~}{\mu }_0=\mu _{c0}`$, where $`\mu _{c0}`$ is the equilibrium chemical potential of the condensate as described by (1).
Using our finite $`T`$ “static Popov” approximation, equation (1) can be simplified to
$$i\mathrm{}\frac{\mathrm{\Phi }}{t}=\left[\frac{\mathrm{}^2}{2m}^2+U_{\mathrm{ext}}+gn_c+2g\stackrel{~}{n}_0i\mathrm{}R_0\right]\mathrm{\Phi },$$
(6)
which describes the condensate motion within the static thermal cloud. Here $`\stackrel{~}{n}_0`$ is the equilibrium density of the non-condensate and the damping term $`R_0`$ is calculated using
$$\mathrm{\Gamma }_{12}^0(𝐫,t)\frac{d𝐩}{(2\pi \mathrm{})^3}C_{12}[f^0(𝐩,𝐫),\mathrm{\Phi }(𝐫,t)].$$
(7)
Notice that $`\mathrm{\Gamma }_{12}^0(𝐫,t)`$ depends on time now only through $`\mathrm{\Phi }(𝐫,t)`$. Using the explicit general expression for $`C_{12}`$ given in Eq. (23b) of ZGN, one finds (see also Ref. )
$$\mathrm{\Gamma }_{12}^0(𝐫,t)=\frac{n_c(𝐫,t)}{\tau _{12}(𝐫,t)}\left[e^{\beta (\stackrel{~}{\mu }_0\epsilon _c(𝐫,t))}1\right],$$
(8)
where we have defined the $`C_{12}`$ collision time
$`{\displaystyle \frac{1}{\tau _{12}(𝐫,t)}}`$ $``$ $`{\displaystyle \frac{2g^2}{(2\pi )^5\mathrm{}^7}}{\displaystyle 𝑑𝐩_1𝑑𝐩_2𝑑𝐩_3\delta (𝐩_c+𝐩_1𝐩_2𝐩_3)}`$ (9)
$`\times `$ $`\delta (\epsilon _c+\stackrel{~}{\epsilon }_{p_1}\stackrel{~}{\epsilon }_{p_2}\stackrel{~}{\epsilon }_{p_3})(1+f_1^0)f_2^0f_3^0.`$ (10)
Here the condensate atom local energy is $`\epsilon _c(𝐫,t)=\mu _c(𝐫,t)+\frac{1}{2}mv_c^2(𝐫,t)`$ with the non-equilibrium condensate chemical potential
$$\mu _c(𝐫,t)=\frac{\mathrm{}^2\mathbf{}^2\sqrt{n_c}}{2m\sqrt{n_c}}+U_{\mathrm{ext}}+gn_c+2g\stackrel{~}{n}_0.$$
(11)
The condensate atom momentum is $`𝐩_c=m𝐯_c`$, and $`f_i^0=f^0(𝐫,𝐩_i)`$. We have introduced the usual condensate velocity defined in terms of the phase $`\theta `$ of the condensate $`\mathrm{\Phi }(𝐫,t)=\sqrt{n_c(𝐫,t)}e^{i\theta (𝐫,t)}`$ as $`𝐯_c=\mathrm{}\mathbf{}\theta (𝐫,t)/m`$. A closed set of equations for $`\mathrm{\Phi }(𝐫,t)`$ is given by (6) and its complex conjugate combined with (8) \- (11).
We note that in terms of $`n_c`$ and $`𝐯_c`$, (6) is completely equivalent to the coupled equations
$`{\displaystyle \frac{n_c}{t}}+\mathbf{}(n_c𝐯_c)`$ $`=`$ $`\mathrm{\Gamma }_{12}^0[f^0,\mathrm{\Phi }]`$ (12)
$`m\left({\displaystyle \frac{}{t}}+𝐯_c\mathbf{}\right)𝐯_c`$ $`=`$ $`\mathbf{}\mu _c.`$ (13)
It is easy to see from (8) that when the condensate is in equilibrium with the thermal cloud according to $`\mu _c\mu _{c0}=\stackrel{~}{\mu }_0`$, $`\mathrm{\Gamma }_{12}^0(𝐫,t)`$ then vanishes. It is clear that the description of the system given by equations (5) and (6), or equivalently (13), is valid only if the condensate is slightly perturbed from equilibrium and the condensate motion is essentially uncoupled from that of the thermal cloud, which we can then treat statically. In order to describe the condensate oscillations about equilibrium, we use the quantum hydrodynamic variables $`n_c(𝐫,t)=n_{c0}(𝐫)+\delta n_c(𝐫,t)`$ and $`𝐯_c(𝐫,t)=\delta 𝐯_c(𝐫,t)`$, where $`n_{c0}(𝐫)`$ is the equilibrium density of the condensate with the associated equilibrium chemical potential $`\mu _{c0}`$. Alternatively, one may work with the fluctuations of $`\mathrm{\Phi }(𝐫,t)`$ and derive coupled Bogoliubov equations generalized to include the effect of the $`R_0`$ damping term. This generalization will be discussed elsewhere .
### B Finite-$`T`$ Stringari wave equation
From (13), we can obtain linearized equations of motion for the condensate fluctuations $`\delta n_c`$ and $`\delta 𝐯_c`$. We use the fact that, to lowest order in the fluctuations from static equilibrium, (8) reduces to
$$\delta \mathrm{\Gamma }_{12}^0=\frac{\beta n_{c0}(𝐫)}{\tau _{12}^0(𝐫)}\delta \mu _c(𝐫,t),$$
(14)
where the “equilibrium” $`C_{12}`$ collision rate is defined by
$`{\displaystyle \frac{1}{\tau _{12}^0(𝐫)}}`$ $``$ $`{\displaystyle \frac{2g^2}{(2\pi )^5\mathrm{}^7}}{\displaystyle 𝑑𝐩_1𝑑𝐩_2𝑑𝐩_3\delta (𝐩_1𝐩_2𝐩_3)}`$ (15)
$`\times `$ $`\delta \left({\displaystyle \frac{p_1^2p_2^2p_3^2}{2m}}gn_{c0}\right)(1+f_1^0)f_2^0f_3^0.`$ (16)
In the present paper, we restrict ourselves to the Thomas-Fermi limit, valid for large $`N_c`$,
$`{\displaystyle \frac{\delta n_c}{t}}+\mathbf{}(n_{c0}\delta 𝐯_c)`$ $`=`$ $`{\displaystyle \frac{1}{\tau ^{}}}\delta n_c`$ (17)
$`m{\displaystyle \frac{\delta 𝐯_c}{t}}`$ $`=`$ $`g\mathbf{}\delta n_c.`$ (18)
The collision time $`\tau ^{}(𝐫)`$ describes collisions between the condensate and non-condensate atoms when the condensate is perturbed away from equilibrium,
$$\frac{1}{\tau ^{}(𝐫)}=\frac{gn_{c0}(𝐫)}{k_\mathrm{B}T}\frac{1}{\tau _{12}^0(𝐫)}.$$
(19)
In the TF approximation the equilibrium distribution reduces to $`f_i^0=[e^{\beta (p_i^2/2m+gn_{c0})}1]^1`$. The new term on the right-hand side of (17) causes damping of the condensate fluctuations due to the lack of collisional detailed balance between the condensate and the static thermal cloud. We note that this collision time is only a function of $`𝐫`$ through its dependence on the static condensate density $`n_{c0}(𝐫)`$. Plots of $`1/\tau ^{}(𝐫)`$ will be discussed below.
We can easily combine (17) and (18) to obtain what we shall refer to as the finite $`T`$ Stringari wave equation
$$\frac{^2\delta n_c}{t^2}\frac{g}{m}\mathbf{}(n_{c0}\mathbf{}\delta n_c)=\frac{1}{\tau ^{}}\frac{\delta n_c}{t}.$$
(20)
Equation (20) is the main result of this paper. If we neglect the right-hand side, we obtain the undamped finite $`T`$ Stringari normal modes $`\delta n_c(𝐫,t)=\delta n_i(𝐫)e^{i\omega _it}`$ given by the solution of
$$\frac{g}{m}\mathbf{}\left[n_{c0}(𝐫)\mathbf{}\delta n_i(𝐫)\right]=\omega _i^2\delta n_i(𝐫).$$
(21)
As has been noted by several authors in recent papers , $`n_{c0}(𝐫)`$ at finite $`T`$ can be well approximated by the TF condensate profile at $`T=0`$ but with the number of atoms in the condensate $`N_c(T)`$ now being a function of temperature, since the static mean field of the non-condensate plays such a minor role. With this approximation for $`n_{c0}(𝐫)`$, the solutions of the finite $`T`$ Stringari equation (21) will be identical to those at $`T=0`$, since the $`T=0`$ Stringari frequencies do not depend on the magnitude of $`N_c`$. Of course, as shown by calculations solving the coupled Bogoliubov equations , the TF approximation breaks down when $`N_c10^4`$. Thus the condensate collective mode frequencies will always become temperature dependent close to $`T_{\mathrm{BEC}}`$, where the TF approximation is no longer valid. We generalize our present discussion to deal with this region in Ref. .
We can use the undamped Stringari modes as a basis set to solve (20) and find the damping of these modes. Writing $`\delta n_c(𝐫)=_ic_i\delta n_i(𝐫)`$, and using the orthogonality condition $`𝑑𝐫\delta n_i(𝐫)\delta n_j(𝐫)=\delta _{ij}`$, one obtains the following algebraic equations for the coefficients $`c_i`$
$$\omega ^2c_i=\omega _i^2c_ii\omega \underset{j}{}\gamma _{ij}c_j,$$
(22)
where
$$\gamma _{ij}=𝑑𝐫\delta n_i(𝐫)\delta n_j(𝐫)/\tau ^{}(𝐫).$$
(23)
Assuming the damping is small (we are in the collisionless region), (22) is easily solved using perturbation theory by setting $`\gamma _{ij}=0`$ for $`ij`$. This gives the damped Stringari frequency (to lowest order) $`\mathrm{\Omega }_i=\omega _ii\mathrm{\Gamma }_i`$, with
$$\mathrm{\Gamma }_i\frac{\gamma _{ii}}{2}=\frac{1}{2}𝑑𝐫\frac{\delta n_i(𝐫)^2}{\tau ^{}(𝐫)}.$$
(24)
This result for $`\mathrm{\Gamma }_i`$ is reasonable, namely it involves an average over $`1/\tau ^{}(𝐫)`$ weighted with respect to the undamped density fluctuations of the Stringari wave equation (21). We find that coupling to other modes ($`\gamma _{ij}0`$) is extremely small.
## III Results
### A Homogeneous gas
Before treating the trapped gas, it is useful to first apply our theory to a homogeneous gas, which was considered previously in Ref. in connection with the collision-dominated hydrodynamic region. For a homogeneous gas, $`\tau ^{}`$ is independent of position and then (24) reduces to $`\mathrm{\Gamma }_i=1/2\tau ^{}`$. Although our model in the present paper applies only to the collisionless region, it is useful to compare the inter-component collision time in both the collisionless and hydrodynamic regimes. In ZNG, it was shown that the inter-component collision time $`\tau _\mu `$ in the hydrodynamic region is given by $`\tau _\mu =\sigma \tau ^{}`$, where the temperature-dependent factor $`\sigma `$ (not to be confused with the collision cross section) depends on various thermodynamic functions. In Fig. 1 we compare $`1/\tau ^{}`$ and $`1/\tau _\mu `$ as functions of $`T`$. We see that $`\sigma `$ dramatically alters the inter-component relaxation rate $`1/\tau _\mu `$ appropriate to the hydrodynamic regime, as compared to $`1/\tau ^{}`$ involved in the collisionless regime. For completeness, in Fig. 1 we also plot the often-used classical collision time given by (4) as well as $`\tau _{12}^0`$ defined in (16).
In a uniform Bose gas at finite temperatures, the Landau damping ($`\omega =cqi\mathrm{\Gamma }_L`$) of condensate modes has been evaluated in several recent papers . Working within the full second-order Beliaev approximation, one finds
$$\mathrm{\Gamma }_L=\left(\frac{3\pi }{8}\right)\frac{ak_\mathrm{B}Tq}{\mathrm{}}.$$
(25)
This is clearly quite different from our inter-component damping $`\mathrm{\Gamma }=1/2\tau ^{}`$, as plotted in Fig. 1. Landau damping originates from the interaction of a condensate collective mode with the excitations of the thermal cloud but is not associated with $`C_{12}`$ collisions which give rise to $`\tau ^{}`$.
In the context of our generalized GP equation in (1), Landau damping comes from the fluctuations in the thermal cloud induced by the condensate mean field,
$$\delta \stackrel{~}{n}=\stackrel{~}{\chi }_0(2g\delta n_c).$$
(26)
In the finite temperature region of interest, $`\stackrel{~}{\chi }_0`$ can be approximated as the density response function of a non-interacting gas of atoms with a spectrum $`\stackrel{~}{\epsilon }_p`$ and chemical potential $`\mu _{c0}`$. For a uniform gas, one sees that using (26) in (1), with $`R=0`$, gives condensate modes satisfying $`\omega ^2=c^2q^2(1+4g\stackrel{~}{\chi }_0)`$ and thus $`\omega =cqi\mathrm{\Gamma }_L`$, where
$$\mathrm{\Gamma }_L=2gcq\mathrm{Im}\stackrel{~}{\chi }_0(q,\omega =cq).$$
(27)
Evaluating $`\mathrm{Im}\stackrel{~}{\chi }_0`$ in the limit of small $`q`$ , one finds $`\mathrm{\Gamma }_L=\frac{4}{3}ak_\mathrm{B}Tq/\mathrm{}`$. Apart from the slightly larger numerical coefficient, this agrees with the exact result given in equation (25) .
### B Trapped gas
We now turn to explicit calculations of the inter-component damping rate using our model for a trapped gas. In order to calculate $`\mathrm{\Gamma }_i`$ for a given mode, the equilibrium chemical potential $`\mu _{c0}=\stackrel{~}{\mu }_0`$ must be calculated self-consistently for a given total number $`N`$ of atoms at a given temperature $`T`$. In the TF approximation, the procedure is straight forward (see, for example, Ref. ). In the following, we consider a harmonic trap with axial symmetry $`U_{\mathrm{ext}}(𝐫)=\frac{1}{2}m\omega _\rho ^2(\rho ^2+\lambda ^2z^2)`$, where $`\lambda =\omega _z/\omega _\rho `$ is the anisotropy parameter. In the TF approximation, the condensate density takes the explicit form $`n_{c0}=[\mu _{c0}m\omega _\rho ^2(\rho ^2+\lambda ^2z^2)/2]/g`$ within the TF radius, and the condensate chemical potential is $`\mu _{c0}=\frac{1}{2}\mathrm{}\omega _\rho [15\lambda N_ca/\rho _0]^{2/5}`$, where $`\rho _0=\sqrt{\mathrm{}/m\omega _\rho }`$. The form of the Stringari normal modes $`\delta n_i(𝐫)`$ is given explicitly in the literature . We mainly consider the breathing mode ($`n=1,l=0`$) for which $`\omega _{10}=\sqrt{5}\omega _\rho `$, for $`\lambda =1`$.
We choose experimentally accessible parameters in the following calculations for the collisionless region. However, we do not compare our results to the two available experiments on damping of normal modes at finite $`T`$, since the TF approximation is not valid for most of the JILA data , and the MIT experiment is approaching the collision-dominated hydrodynamic limit where the dynamics of the condensate and non-condensate become more strongly coupled. For <sup>87</sup>Rb the scattering length is $`a5.7`$ nm . We first consider a spherically symmetric trap $`\lambda =1`$, with trap frequency $`\nu _r=10`$ Hz, and we take $`N=2\times 10^6`$. In the collisionless limit, we require $`\omega _i\tau _{\mathrm{cl}}1`$, taking $`\tau _{\mathrm{cl}}`$ as defined in (4). For a trapped gas, we obtain an upper limit on $`1/\tau _{\mathrm{cl}}`$ by taking the density in the center of the trap $`n(r=0)`$, which gives $`1/\tau _{\mathrm{cl}}=8a^2N\omega _\rho ^3m/(\pi k_\mathrm{B}T)`$. For the parameters we use, $`\omega _{10}\tau _{\mathrm{cl}}19`$ (compared to $`\omega _{02}\tau _{\mathrm{cl}}20`$ for the JILA data , and $`\omega _{02}\tau _{\mathrm{cl}}2`$ for the MIT data ).
In Fig. 2a we plot $`1/\tau ^{}(𝐫)`$ vs. position for $`T=0.9T_{\mathrm{BEC}}`$ and $`T=0.5T_{\mathrm{BEC}}`$. We see that the collision rate increases steadily up to the condensate boundary, but as $`T`$ increases, $`1/\tau ^{}(𝐫)`$ becomes relatively constant. The behavior of $`1/\tau ^{}(𝐫)`$ just seems to be mimicking the behavior of the non-condensate density $`\stackrel{~}{n}(𝐫)`$, which we plot in the inset of Fig. 2b along with the condensate density. The condensate mean-field pushes the non-condensate density out of the center of the trap, a well known result . We also show the breathing mode density fluctuation in Fig. 2b. The sharp cusp of $`\stackrel{~}{n}(𝐫)`$ and the sudden drop of $`\delta n_i(𝐫)`$ and $`1/\tau ^{}(𝐫)`$ at the condensate boundary are all unphysical artifacts of the TF approximation. Inclusion of the kinetic energy pressure in a more accurate calculation would have the effect of smoothing out this behavior at the boundary. To illustrate the effect of the kinetic energy pressure, we also show in Fig. 2b the breathing mode obtained by solving the $`T=0`$ coupled Bogoliubov equations. We estimate that an improved treatment which includes the kinetic energy pressure will modify our estimate of $`\mathrm{\Gamma }_{10}`$ by about $`1020\%`$ .
In Fig. 3a we plot the damping rate $`\mathrm{\Gamma }_{10}`$ for the breathing mode ($`n=1`$, $`l=0`$) shown in Fig. 2b as a function of temperature up to $`T=0.95T_{\mathrm{BEC}}`$, where $`N_c7\times 10^4`$. At higher temperatures, the Thomas-Fermi approximation will start to break down and the mode frequencies become temperature dependent .
Landau damping of condensate modes in trapped gases has also been discussed in some detail in recent papers . These papers give results that are in qualitative agreement with the expression for a uniform gas in (25), with $`q=\omega _\rho /c`$ and evaluating the Bogoliubov sound velocity $`c`$ for the density at the center of the trap . This simple estimate for Landau damping is plotted in Fig. 3 for comparison. We note that it is larger but comparable to the inter-component collisional damping that we consider. Clearly, a fully satisfactory theory of finite $`T`$ damping of normal modes must include *both* Landau damping as well as the damping we consider in this paper due to the condensate being out of diffusive equilibrium with the non-condensate.
Our theory is easily applied to anisotropic traps. In Fig. 3b we show the damping of $`m=0,2`$ quadrupole modes for an axially symmetric trap with $`\lambda =\sqrt{8}`$. Here we choose a slightly tighter trap $`\nu _r=23`$ Hz, and we take $`N=1\times 10^6`$ (in this case, $`N_c3\times 10^4`$ at $`T=0.95T_{\mathrm{BEC}}`$). For these parameters, we find $`\omega _{20}\tau _{\mathrm{cl}}6`$. In Fig. 3b we see that Landau damping is about twice as large as our inter-component collisional damping.
It is instructive to also consider the dependence of the mode damping $`\mathrm{\Gamma }_i`$ on the total population $`N`$. In Fig. 4, we show a shaded surface plot of $`\mathrm{\Gamma }_{10}`$ for the breathing mode in an isotropic trap as a function of $`T`$ and $`N`$. The white line at $`N=2\times 10^6`$ corresponds to the solid line plotted in Fig. 3a. As one might expect, the inter-component damping rate increases with increasing total population (since the density is increasing). It is also important to realize that in current experiments, the data taken is for a broad range of $`N`$ due to evaporative cooling losses . The idealized fixed-$`N`$ line can never be achieved in practice and one is instead dealing with a curve like the longer line on the surface in Fig. 4.
## IV Conclusion
In summary, we have calculated a new damping mechanism of condensate collective modes due to collisions with the thermal cloud, based on the finite-$`T`$ equations derived in Ref. . The essential mechanism involves the lack of diffusive equilibrium between the condensate and the thermal cloud , which also plays a key role in the theory of condensate growth . Here we have carried-out the first explicit calculation of this damping mechanism for a trapped gas in the collisionless regime (this inter-component damping has recently also been evaluated in the collision-dominated hydrodynamic regime ). In recent discussions of the damping of condensate collective modes in the collisionless region, the mechanism we consider is omitted. One instead focuses on the dynamical mean-field coupling between the condensate and thermal cloud, which gives rise to Landau and Beliaev damping . While we have not considered it in detail, we have indicated how we could include Landau damping by considering the non-condensate fluctuations in (1) induced by the condensate mean field . Comparing Landau damping to the additional mechanism we have calculated, we find that the two are comparable in size. Further experimental studies of the collective modes at finite temperatures are needed to clarify the relative importance of these different sources of damping.
In this paper, we have argued that a good first estimate of the inter-component damping of condensate collective modes can be obtained by coupling it to a static thermal cloud. A more systematic theory is clearly desirable in which the collisionless dynamics of the thermal cloud is allowed for. However, as noted in the introduction, we do not believe that this will lead to significant corrections to the inter-component damping of out-of-phase condensate modes in which the motion of the thermal cloud is not significant. In future work, we hope to discuss the damping due to $`C_{12}`$ collisions of collective modes that mainly involve the motion of the thermal cloud, with the condensate being treated statically.
This work grew out of discussions with E. Zaremba about the role of $`C_{12}`$ collisions in the collisionless region. We also thank T. Nikuni for useful comments. This research was supported by a grant from NSERC.
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# 1 Introduction
## 1 Introduction
In the development of vertex operator algebra theory, one of the most important problems is to construct new solvable vertex operator (super)algebras in the sense that irreducible modules and fusion rules can be completely determined and that intertwining operators can be explicitly constructed. To a certain extent, such algebras give rise to solvable physical models. One of many ways to get such vertex operator (super)algebras is to consider certain extensions of some well known algebras. For example (\[MS\], \[Li5\]), the vertex operator algebra $`L(k,0)`$ associated to the affine Lie algebra $`\widehat{s}l_2`$ with a positive even integral level $`k`$ can be extended to a vertex operator (super)algebra $`L(k,0)+L(k,k)`$. When $`k`$ is odd, $`L(k,0)+L(k,k)`$ does not have an extended vertex operator (super)algebra structure because of the failure of the locality. (For a certain class of vertex operator algebras, e.g., vertex operator algebras associated to positive-definite even lattices, as proved in \[DL\], the sum of a copy of each irreducible module does have a nice structure, called an abelian intertwining algebra. See \[DL\], \[FFR\], \[M\] and \[Hua2\] for notions of various generalized structures.)
In \[FM\], Feigin and Miwa constructed a family of extended vertex operator (super)algebras $`A_k`$ from the vertex operator algebras $`L(k,0)`$ associated to the affine Lie algebra $`\widehat{s}l_2`$ with an arbitrary positive integral level $`k`$, and they classified all irreducible modules and determined all fusion rules. In addition they obtained very interesting results on the monomial basis for irreducible modules. This paper was mainly motivated by \[FM\] and \[DLM2\]. As the main results of this paper we generalize their results except the monomial basis result to affine Lie algebras of other types by using a different approach.
The algebras $`A_k`$ were defined in \[FM\] by a set of mutually local vertex operators (or fields). On the other hand, in terms of vertex operator algebra language, $`A_k`$ are extensions (by an infinite sum of irreducible modules) of vertex operator algebra $`L(k,0)M(1,0)`$, where $`M(1,0)`$ is the vertex operator algebra associated to an infinite-dimensional Heisenberg Lie algebra of rank one, or a single free bosonic field. The essential building block of $`A_k`$ is the irreducible $`L(k,0)M(1,0)`$-module $`L(k,k)M(1,\alpha )`$ for some $`\alpha C`$.
The $`L(k,0)`$-module $`L(k,k)`$ has been known to be a simple current (\[FG\], \[GW\], \[SY\]) in the sense that the left multiplication of the equivalence class $`[L(k,k)]`$ in the Verlinde algebra gives rise to a permutation on the standard basis. It was known to physicists (cf. \[MS\], \[SY\]) that a simple current with integer weights can be included to generate an extended vertex operator algebra. In \[Li5\], as an exercise by using an explicit construction of simple currents given in \[Li4\] we studied the extension of a certain vertex operator algebra by a self-dual (or order 2) simple current where $`L(k,0)+L(k,k)`$ is a special case. For such extended algebras, all their irreducible modules were classified and the complete reducibility of every module was proved. A little bit latter, the results of \[Li5\] were greatly extended in \[DLM2\].
The construction of simple currents given in \[Li4\] was based on a result of \[Li2\]. Let $`V`$ be a vertex operator algebra and let $`h`$ be a weight one primary vector in $`V`$ such that the component operators $`h(m)`$ of the vertex operator $`Y(h,z)`$ satisfy the Heisenberg algebra relation and such that $`h(0)`$ is semisimple on $`V`$ with only rational eigenvalues. Clearly, $`\sigma _h:=e^{2\pi ih(0)}`$ is an automorphism of $`V`$. Define
$$\mathrm{\Delta }(h,z)=z^{h(0)}\mathrm{exp}\left(\underset{n=1}{\overset{\mathrm{}}{}}\frac{h(n)}{n}(z)^n\right),$$
an element of $`(\mathrm{End}V)\{z\}`$. It was proved in \[Li2\] that for any $`V`$-module $`W`$,
$$(W^{(h)},Y_h(,z)):=(W,Y(\mathrm{\Delta }(h,z),z))$$
is a $`\sigma _h`$-twisted $`V`$-module. In particular, this gives an (untwisted) $`V`$-module if $`h(0)`$ acting on $`V`$ has only integral eigenvalues. It was furthermore proved in \[Li4\] that if $`I(,z)`$ is an intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$ (in the sense of \[FHL\]), then $`I(\mathrm{\Delta }(h,z),z)`$ is an intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W_3^{(h)}}{W_1W_2^{(h)}}\right)`$. By using this result and the invertiability of $`\mathrm{\Delta }(h,z)`$, it was proved that $`V^{(h)}`$ is a simple current. As a matter of fact, for certain well known vertex operator algebras almost all simple currents can be constructed in this way. For example, when $`V`$ is the vertex operator algebra $`V_L`$ associated with a positive-definite even lattice $`L`$, it was proved that all irreducible $`V_L`$-modules can be constructed this way. In this case, this construction is intimately related to the construction of twisted modules by using shifted vertex operators in \[Le\]. When $`V=L(\mathrm{},0)`$ associated to an affine Lie algebra $`\widehat{g}`$ with $`gE_8`$ and with a positive integral level $`\mathrm{}`$, all the simple currents can also be constructed this way. The merit of this construction of a simple current is on the canonicalness of the vector space and the vertex operator map (in terms of the algebra $`V`$ and the element $`h`$). With this construction, certain intertwining operators can be constructed canonically. The canonical construction of intertwining operators of certain types is the basis of \[Li5\] and \[DLM2\].
The essential results of \[DLM2\] can be described as follows: Let $`V`$ be a vertex operator algebra and $`H`$ be a subspace of $`V_{(1)}`$ such that the components $`h(n)`$ of vertex operators $`Y(h,z)`$ for $`hH,nZ`$ satisfy the Heisenberg algebra relation. Let $`L`$ be a subgroup of $`H`$ such that for every $`\alpha L`$, $`\alpha (0)`$ acts semisimply on $`V`$ with only integral eigenvalues. Consider the space $`V[L]=C[L]V`$, where for $`\alpha L`$, $`e^\alpha V`$ is identified with $`V^{(\alpha )}`$ equipped with the $`V`$-module structure $`Y_\alpha `$. Extend the $`V`$-module structure on $`V[L]`$ in a certain canonical way to a vertex operator map $`Y`$ on $`V[L]`$. Then it was proved that $`V[L]`$ equipped with the defined vertex operator map $`Y`$ is a generalized vertex algebra in the sense of \[DL\]. It was proved that $`V[L]`$ is a vertex operator (super)algebra when $`L`$ satisfies certain conditions. When $`V=M_𝐡(1,0)`$ with $`𝐡=C_ZL`$, where $`L`$ is an integral lattice, we have $`V[L]=V_L`$ (\[FLM\], \[DL\]). Then we may view $`V[L]`$ as a generalization of $`V_L`$. In \[DLM2\], we were mainly interested in the case $`V=L(k,0)`$ associated to an affine algebra $`\widehat{g}`$ with a positive integral level $`k`$. Because $`L(k,0)`$ has only finitely many irreducible modules up to equivalence, each irreducible $`V`$-module $`V^{(\alpha )}`$ in $`V[L]`$ is not multiplicity-free. Having noticed this we proved that $`V[L]`$ has a quotient algebra $`\overline{V[L]}`$ such that every irreducible $`V`$-module in $`\overline{V[L]}`$ is multiplicity-free and that $`V[L]`$ and $`\overline{V[L]}`$ contain the same number of non-isomorphic irreducible $`V`$-modules. An irreducible $`V`$-module $`W`$ was also extended to $`W[L]`$ on which $`V[L]`$ acts and it was proved that $`W[L]`$ is in general a twisted $`V[L]`$-module with respect to a certain automorphism of $`V[L]`$. With this result, under the assumption of complete reducibility of $`V`$-modules, all irreducible $`V[L]`$-modules were classified and a complete reducibility theorem for $`V[L]`$-modules was proved.
In this paper, to generalize Feigin and Miwa’s construction we apply the results of \[DLM2\] by taking $`V=L(k,0)M_𝐡^{}(1,0)`$ and by choosing $`𝐡^{}`$ and $`L`$ appropriately, depending on the type of $`g`$. In this case, each irreducible $`V`$-module in $`V[L]`$ is multiplicity-free and $`V[L]`$ is a simple algebra. On the other hand, since the category of $`V`$-modules is not semisimple, the results of \[DLM2\] for the complete reducibility of $`V[L]`$-modules do not apply to this case directly. The complete reducibility of $`V[L]`$-modules is proved here. We also naturally extend intertwining operators for modules in the category of $`V`$-modules to intertwining operators for modules in the category of $`V[L]`$-modules. Using this result we are able to derive a formula of fusion rules for $`V[L]`$-modules in terms of fusion rules for $`V`$-modules.
This paper is organized as follows: In Section 2, we recall the classification of irreducible modules for certain vertex operator algebras and recall a construction of simple currents. Most part of this section is preliminary and the only new result is about the fusion rules for simple currents for $`L(k,0)`$. In Section 3 we recall and refine some of the results of \[DLM2\] on the extended vertex (super)algebra $`V[L]`$. We furthermore study the multiplicity-free case. In Section 4, we apply the results of Section 3 to construct extended vertex operator (super)algebras associated to an affine algebra $`\widehat{g}`$.
## 2 Vertex operator algebras and simple currents
The extended vertex operator (super) algebras we shall construct are based on vertex operator algebras $`L_g(\mathrm{},0)`$, $`M_𝐡(1,0)`$, $`V_L`$ and their representations. For this reason, in this section we shall recall the relevant information about these algebras and we also recall from \[Li4\] a construction of simple currents for a certain type of vertex operator algebras, including $`L_g(\mathrm{},0)`$, $`M_𝐡(1,0)`$, $`V_L`$. For the vertex operator algebra $`L_g(\mathrm{},0)`$ associated to an affine algebra $`\widehat{g}`$ (not type $`E_8^{(1)}`$) with a positive integral level $`\mathrm{}`$, we prove that the equivalence classes of the simple currents form an abelian group isomorphic to $`P^{}/Q^{}`$, where $`P^{}`$ and $`Q^{}`$ are the co-weight and co-root lattices of $`g`$.
### 2.1 Vertex operator algebras $`L_g(\mathrm{},0)`$, $`M_𝐡(1,0)`$ and $`V_L`$
We shall use standard definitions and notations as given in \[FHL\] and \[FLM\] and we also use \[K\] and \[H\] as our references for (Kac-Moody) Lie algebras. Following \[DL\] we use the term “vertex (super)algebra” for an object that satisfies all the axioms defining the notion of vertex operator (super)algebra except the two grading restrictions.
Let $`g`$ be a finite-dimensional simple Lie algebra, $`𝐡`$ a Cartan subalgebra, and $`,`$ the normalized killing form such that the square length of a long root is 2. Let
$`\widehat{g}=gC[t,t^1]Cc`$ (2.1)
be the affine Lie algebra.
Let $`\mathrm{}`$ be a complex number such that $`\mathrm{}h^{}`$, where $`h^{}`$ is the dual Coxeter number of $`g`$. Let $`C_{\mathrm{}}`$ be the one-dimensional $`(gC[t]+Cc)`$-module on which $`c`$ acts as scalar $`\mathrm{}`$ and $`gC[t]`$ acts as zero. Form the generalized Verma $`\widehat{g}`$-module
$`M_g(\mathrm{},0)=U(\widehat{g})_{U(gC[t]+Cc)}C_{\mathrm{}}.`$ (2.2)
It was well known (cf. \[FF\], \[FZ\], \[Li1\], \[MP\]) that $`M_g(\mathrm{},0)`$ has a natural vertex operator algebra structure. Furthermore, the category of weak $`M_g(\mathrm{},0)`$-modules in the sense that all the axioms defining the notion of module except those involving grading hold is canonically equivalent to the category of restricted (cf. \[K\]) $`\widehat{g}`$-modules of level $`\mathrm{}`$ in the sense that for every vector $`w`$ of the module, $`(gt^nC[t])w=0`$ for $`n`$ sufficiently large.
###### Remark 2.1
More generally, let $`W`$ be an arbitrary vector space. An element $`a(z)`$ of $`(\mathrm{End}W)[[z,z^1]]`$ is called a vertex operator if $`a(z)wW((z))`$ for every $`wW`$. Two vertex operators $`a(z)`$ and $`b(z)`$ are said to be mutually local if there exists a nonnegative integer $`N`$ such that
$`(z_1z_2)^N[a(z_1),b(z_2)]=0`$ (2.3)
(cf. \[DL\], (1.4)). It was proved in \[Li1\] (Corollary 3.2.11) that any set of mutually local vertex operators on $`W`$ automatically generates a vertex algebra, which is a canonical vector subspace of $`(\mathrm{End}W)[[z,z^1]]`$, and that $`W`$ is a natural module for this vertex algebra.
For $`\lambda 𝐡^{}`$, denote by $`L_g(\mathrm{},\lambda )`$ the irreducible highest weight $`\widehat{g}`$-module of level $`\mathrm{}`$ with highest weight $`\lambda `$. Each $`L_g(\mathrm{},\lambda )`$ is an irreducible $`M_g(\mathrm{},0)`$-module possibly with infinite-dimensional homogeneous subspaces.
Denote by $`\theta `$ the highest long root of $`g`$. For a positive integer $`\mathrm{}`$, set
$`P_{\mathrm{}}=\{\lambda P_+|\lambda ,\theta \mathrm{}\},`$ (2.4)
where $`P_+`$ is the set of dominant integral weights of $`g`$. Then $`L(\mathrm{},\lambda )`$ is an integrable $`\widehat{g}`$-module if and only if $`\lambda P_{\mathrm{}}`$ \[K\]. The following result was known (cf. \[DL, Proposition 13.17\], \[FZ, Theorem 3.1.3\], \[Li1, Propositions 5.2.4 and 5.2.5\], \[MP, Theorems 5.9, 5.14 and 5.15\]):
###### Proposition 2.2
Let $`\mathrm{}`$ be a positive integer. Then (1) The set of irreducible $`L_g(\mathrm{},0)`$-modules is exactly the set of irreducible highest weight integrable (or standard) $`\widehat{g}`$-modules of level $`\mathrm{}`$. (2) Every $`L_g(\mathrm{},0)`$-module is completely reducible.
The following stronger result was obtained in \[DLM1\] (Theorem 3.7):
###### Proposition 2.3
Let $`\mathrm{}`$ be a positive integer. Then every weak $`L_g(\mathrm{},0)`$-module is a direct sum of irreducible highest weight integrable (or standard) $`\widehat{g}`$-modules of level $`\mathrm{}`$. In particular, every irreducible weak $`L_g(\mathrm{},0)`$-module is an (ordinary) $`L_g(\mathrm{},0)`$-module.
A vertex operator algebra with the property that every weak module is a direct sum of irreducible (ordinary) modules is said to be regular \[DLM1\].
Let $`𝐡`$ be a finite-dimensional vector space equipped with a nondegenerate symmetric bilinear form $`,`$, i..e, a finite-dimensional abelian Lie algebra equipped with a nondegenerate symmetric invariant bilinear form. Let $`\widehat{𝐡}=𝐡C[t,t^1]+Cc`$ be the affine Lie algebra. We have
$`\widehat{𝐡}=\widehat{𝐡}_Z𝐡,`$ (2.5)
where $`\widehat{𝐡}_Z=_{n0}𝐡t^n+Cc`$ is a Heisenberg algebra and $`𝐡`$ is central. Similar to the construction of $`M_g(\mathrm{},0)`$ we construct a space $`M_𝐡(1,0)`$, and just like $`M_g(\mathrm{},0)`$, $`M_𝐡(1,0)`$ is a vertex operator algebra. For any $`\alpha 𝐡^{}`$ $`(=𝐡)`$, let $`Ce^\alpha `$ be a one-dimensional $`(𝐡C[t]+Cc)`$-module on which $`𝐡tC[t]`$ acts as zero, $`c`$ acts as $`1`$ and $`h=h(0)`$ acts as scalar $`\alpha ,h`$ for $`h𝐡`$. Form the induced module
$`M_𝐡(1,\alpha )=U(\widehat{𝐡})_{U(𝐡C[t]+Cc)}Ce^\alpha M_𝐡(1,0)Ce^\alpha \text{(linearly)}.`$ (2.6)
As in the case of $`M_g(\mathrm{},0)`$, $`M_𝐡(1,\alpha )`$ is an irreducible $`M_𝐡(1,0)`$-module. Furthermore, from \[FLM\] the lowest $`L(0)`$-weight of $`M_𝐡(1,\alpha )`$ is
$`\mathrm{\Delta }_\alpha ={\displaystyle \frac{1}{2}}\alpha ,\alpha .`$ (2.7)
On the other hand, clearly every irreducible $`M_𝐡(1,0)`$-module is isomorphic to $`M_𝐡(1,\alpha )`$ for some $`\alpha `$. It follows from the complete reducibility of certain $`\widehat{𝐡}_Z`$-modules (\[LW\], \[K\]) that an $`M_𝐡(1,0)`$-module on which $`𝐡`$ semisimply acts is completely reducible. In general, an $`M_𝐡(1,0)`$-module may not be completely reducible.
Let $`P`$ be a rational lattice of finite rank with the $`Z`$-bilinear form $`,`$. Set
$`𝐡=C_ZP,`$ (2.8)
and extend $`,`$ to a $`C`$-bilinear form on $`𝐡`$.
Denote by $`C[P]`$ the group algebra. Set
$`V_P=C[P]M_𝐡(1,0),`$ (2.9)
equipped with the standard $`M_𝐡(1,0)`$-module (or $`\widehat{𝐡}_Z`$-module) structure. That is, $`V_P`$ is a direct sum of irreducible $`M_𝐡(1,0)`$-modules $`M_𝐡(1,\alpha )Ce^\alpha M_𝐡(1,0)`$ for $`\alpha P`$.
It was proved in \[FLM\] (cf. \[B\]) that when $`P=L`$ is even and positive-definite, $`V_L`$ has a natural simple vertex operator algebra structure which extends the $`M_𝐡(1,0)`$-module structure. (It follows from Proposition 3.22 that such a vertex operator algebra structure is unique up to equivalence.) Furthermore, let $`L^o`$ be the dual lattice of $`L.`$ Then $`V_{L^o}`$ is a natural $`V_L`$-module and $`V_{\beta +L}`$ is an irreducible $`V_L`$-module for $`\beta L^o.`$
###### Proposition 2.4
(1) Let $`\{\beta _1,\mathrm{},\beta _r\}`$ be a complete set of representatives of cosets of $`L`$ in $`L^o`$. Then $`\{V_{\beta _1+L},\mathrm{},V_{\beta _r+L}\}`$ is a complete set of representatives of equivalent classes of irreducible $`V_L`$-modules. (2) Every $`V_L`$-module is completely reducible. (3) $`V_L`$ is regular, i.e., every weak $`V_L`$-module is a direct sum of irreducible (ordinary) $`V_L`$-modules.
The assertion (1) was proved in \[FLM\] and \[D1\], (2) was proved in \[Guo\] and (3) was proved in \[DLM1\].
###### Remark 2.5
For a general rational lattice $`P`$, $`V_P`$ is not a vertex (operator) algebra because of the involvement of non-local vertex operator. A notion of generalized vertex algebra was introduced in \[DL\] with a generalized Jacobi identity as one of the main axioms and it was proved in \[DL, Theorem 5.1 and Remark 9.11\] that $`V_P`$ is a generalized vertex algebra.
### 2.2 Simple currents and a construction
We first recall from \[FHL\] the definition of an intertwining operator.
###### Definition 2.6
Let $`V`$ be a vertex operator algebra and let $`W_1,W_2`$ and $`W_3`$ be $`V`$-modules. An intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_1}\right)`$ is a linear map $`I`$ from $`W_1`$ to $`(\mathrm{Hom}(W_2,W_3))\{z\}`$, (where for a vector space $`U`$, $`U\{z\}`$ is defined to be the vector space of $`U`$-valued formal series in $`z`$ with arbitrary complex powers of $`z`$), satisfying the following properties: for $`w_1W_1,w_2W_2`$,
$`I(w_1,z)w_2z^{\gamma _1}W_3[[z]]+\mathrm{}+z^{\gamma _n}W_3[[z]]`$ (2.10)
for some (finitely many) complex numbers $`\gamma _1,\mathrm{},\gamma _n`$, and for $`vV,w_1W_1`$,
$`[L(1),I(w_1,z)]={\displaystyle \frac{d}{dz}}I(w_1,z),`$ (2.11)
$`z_0^1\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)Y(v,z_1)I(w_1,z_2)z_0^1\delta \left({\displaystyle \frac{z_2z_1}{z_0}}\right)I(w_1,z_2)Y(v,z_1)`$ (2.12)
$`=`$ $`z_2^1\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)I(Y(v,z_0)w_1,z_2).`$
All intertwining operators of type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$ form a vector space denoted by $`I_{W_1W_2}^{W_3}`$. The dimension of $`I_{W_1W_2}^{W_3}`$ is called the fusion rule, denoted by $`N_{W_1,W_2}^{W_3}`$. Clearly, the fusion rule only depends on the equivalence class of each $`W_i`$.
The following are among the immediate consequences of the Jacobi identity (2.12):
$`[v_n,I(w_1,z_2)]={\displaystyle \underset{i0}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)z_2^nI(v_iw_1,z_2)`$ (2.13)
(the commutator formula) and
$`I(v_nw_1,z_2)`$ (2.14)
$`=`$ $`\mathrm{Res}_{z_1}\left((z_1z_2)^nY(v,z_1)I(w_1,z_2)(z_2+z_1)^nI(w_1,z_2)Y(v,z_1)\right)`$
(the iterate formula) for $`nZ`$. Conversely, the commutator and iterate formulas imply the Jacobi identity.
###### Remark 2.7
If $`V=L_g(\mathrm{},0)`$ and $`W_1=L_g(\mathrm{},\lambda )`$, for $`w_1L(\lambda )`$ (the lowest $`L(0)`$-weight subspace of $`L_g(\mathrm{},\lambda )`$), the commutator formula (2.13) gives
$`[a_n,I(w_1,z_2)]=z_2^nI(aw_1,z_2)`$ (2.15)
for $`agL_g(\mathrm{},0)`$ and for $`nZ`$. In many literatures such as \[TK\] (and \[FM\]), an intertwining operator was defined on $`L(\lambda )`$ with the properties (2.11) and (2.15) as the defining axioms. However, it can be proved (cf. \[TK\], \[Li3, Li6\]) that the two definitions give rise to the same fusion rules.
Let $`V`$ be a vertex operator algebra. We denote by $`\mathrm{Irr}(V)`$ the set of equivalence classes of irreducible modules and for an irreducible $`V`$-module $`W`$, denote by $`[W]`$ the equivalence class. $`V`$ is said to be quasi-rational \[MS\] if all fusion rules associated with irreducible modules are finite and if for any $`[W_1],[W_2]\mathrm{Irr}(V)`$, $`N_{[W_1],[W_2]}^{[W_3]}=0`$ for all but finitely many $`[W_3]\mathrm{Irr}(V)`$.
For a quasi-rational vertex operator algebra $`V`$, the Verlinde algebra $`𝒜(V)`$ is defined to be an algebra (over $`C`$) with $`\mathrm{Irr}(V)`$ as a basis and with the fusion rules as the structural constants, i.e.,
$`[W_i][W_j]={\displaystyle \underset{[W_k]\mathrm{Irr}(V)}{}}N_{[W_i],[W_j]}^{[W_k]}[W_k].`$ (2.16)
When $`V`$ is simple, it is easy to show that $`[V]`$ is the unit. It follows immediately from \[FHL, HL2\] that $`𝒜(V)`$ is commutative. Under certain conditions, Huang \[Hua1\] established the associativity, but for a general $`V`$ the associativity is still an unsolved problem.
Let $`V=L_{sl(2)}(k,0)`$ with a positive integer $`k`$. It is well known (\[GW\], \[TK\], \[FZ\]) that the Verlinde algebra has the following relations:
$`[L(k,i)][L(k,j)]={\displaystyle \underset{r=\text{max}(ij,ji)}{\overset{\text{min}(i+j,2kij)}{}}}[L(k,r)].`$ (2.17)
The following definition is due to Schellekens and Yankielowicz \[SY\].
###### Definition 2.8
An irreducible $`V`$-module $`W`$ is called a simple current if the associated matrix of the left multiplication of $`[W]`$ with respect to the standard basis of the Verlinde algebra is a permutation. The order of the associated matrix as a group element is called the order of $`W`$.
We now recall a construction of simple currents from \[Li4\]. In the following, one may think of $`V`$ as one of, or more general, any tensor product of the following vertex operator algebras:
$$L_g(\mathrm{},0),M_𝐡(1,0),V_L,L_g(\mathrm{},0)M_𝐡(1,0),L_g(\mathrm{},0)V_L.$$
Let $`\alpha V_{(1)}`$ satisfying the following conditions:
$`L(n)\alpha =\delta _{n,0}\alpha ,\alpha (n)\alpha =\delta _{n,1}\gamma \mathrm{𝟏}\text{for }nZ_+,`$ (2.18)
where $`Y(\alpha ,z)=_{nZ}\alpha (n)z^{n1}`$, i.e., $`\alpha (n)=\alpha _n`$, and $`\gamma `$ is a fixed complex number. Notice that condition (2.18) implies that $`\alpha `$ is a primary vector and that operators $`\alpha (n)`$ satisfy the Heisenberg algebra relation
$`[\alpha (m),\alpha (n)]=m\gamma \delta _{m+n,0}\text{ for }m,nZ.`$ (2.19)
Furthermore, assume that $`\alpha (0)`$ acts semisimply on $`V`$. It is clear that $`e^{2\pi i\alpha (0)}`$ is an automorphism of $`V`$ and that $`e^{2\pi i\alpha (0)}=1`$ if and only if $`\alpha (0)`$ has only integral eigenvalues on $`V`$. If each $`\alpha (0)`$ has only rational eigenvalues on $`V`$ and if $`V`$ is finitely generated, then $`e^{2\pi i\alpha (0)}`$ is of finite order. Define
$`\mathrm{\Delta }(\alpha ,z)=z^{\alpha (0)}\mathrm{exp}\left({\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\alpha (k)}{k}}(z)^k\right).`$ (2.20)
This is a well defined element of $`(\mathrm{End}W)\{z\}`$ for any weak $`V`$-module $`W`$ on which $`\alpha (0)`$ semisimply acts. The following result is a special case of Proposition 5.4 of \[Li2\].
###### Proposition 2.9
Let $`\alpha ,\mathrm{\Delta }(\alpha ,z)`$ be given as before. Assume that $`\alpha (0)`$ has only integral eigenvalues on $`V`$. Let $`W`$ be any (irreducible) weak $`V`$-module. Set
$`(W^{(\alpha )},Y_\alpha (,z))=(W,Y(\mathrm{\Delta }(\alpha ,z),z)).`$ (2.21)
Then $`(W^{(\alpha )},Y_\alpha )`$ carries the structure of an (irreducible) weak $`V`$-module.
As a convention, by $`V`$-module $`W^{(\alpha )}`$ we mean the $`V`$-module $`(W^{(\alpha )},Y_\alpha )`$.
Recall the following result from \[Li4\] (Proposition 2.5):
###### Proposition 2.10
Let $`\alpha ,\mathrm{\Delta }(\alpha ,z)`$ be as in Proposition 2.9. Let $`W_1`$ and $`W_2`$ be weak $`V`$-modules and $`f`$ be a $`V`$-homomorphism from $`W_1`$ to $`W_2`$. Then $`f`$ is also a $`V`$-homomorphism from $`W_1^{(\alpha )}`$ to $`W_2^{(\alpha )}`$.
###### Remark 2.11
In view of Propositions 2.9 and 2.10, we obtain a canonical functor $`F_\alpha `$ from the category of weak $`V`$-modules to itself in the obvious way. Since
$`\mathrm{\Delta }(\alpha ,z)\mathrm{\Delta }(\alpha ,z)=\mathrm{\Delta }(\alpha ,z)\mathrm{\Delta }(\alpha ,z)=1,`$ (2.22)
we easily see that $`F_\alpha `$ is the inverse functor of $`F_\alpha `$. Therefore, $`F_\alpha `$ is an isomorphism.
The following result \[Li4, Proposition 2.4\] is a generalization of Proposition 2.9:
###### Proposition 2.12
Let $`\alpha ,\mathrm{\Delta }(\alpha ,z)`$ be given as in Proposition 2.9. Let $`I`$ be an intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$. Define
$`I_\alpha (w_1,z)w_2=I(\mathrm{\Delta }(\alpha ,z)w_1,z)w_2`$ (2.23)
for $`w_1W_1,w_2W_2`$. Then $`I_\alpha `$ is an intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W_3^{(\alpha )}}{W_1W_2^{(\alpha )}}\right)`$.
The following results \[Li4, Corollary 2.12, Theorem 2.15, Proposition 3.2\] give a construction of simple currents:
###### Theorem 2.13
Let $`V`$ be a simple vertex operator algebra and let $`\alpha V_{(1)}`$ be such that (2.18) holds and such that $`\alpha (0)`$ has only integral eigenvalues on $`V`$. If for each irreducible (ordinary) $`V`$-module $`W`$, the weak module $`W^{(\alpha )}`$ is an (ordinary) $`V`$-module, then $`V^{(\alpha )}`$ is a simple current. Furthermore, for any irreducible $`V`$-module $`W`$,
$`[W][V^{(\alpha )}]=[W^{(\alpha )}].`$ (2.24)
The following lemma gives more information about $`V^{(\alpha )}`$.
###### Lemma 2.14
Let $`V,\alpha `$ be as in Theorem 2.13. Then $`L(0)`$ acts semisimply on $`V^{(\alpha )}`$ with eigenvalues in $`\frac{1}{2}\gamma +Z`$, where $`\alpha (1)\alpha =\gamma \mathrm{𝟏}`$.
Proof. From \[Li5, (3.18)\], we have
$`\mathrm{\Delta }(\alpha ,z)\omega =\omega +\alpha z^1+{\displaystyle \frac{1}{2}}\alpha (1)\alpha z^2=\omega +\alpha z^1+{\displaystyle \frac{1}{2}}\gamma \mathrm{𝟏}z^2,`$ (2.25)
where $`\omega `$ is the Virasoro element. Then
$`Y_\alpha (\omega ,z)=Y(\mathrm{\Delta }(\alpha ,z)\omega ,z)=Y(\omega ,z)+z^1Y(\alpha ,z)+{\displaystyle \frac{1}{2}}\gamma z^2.`$ (2.26)
In terms of components we have
$`L_\alpha (m)=L(m)+\alpha (m)+{\displaystyle \frac{1}{2}}\gamma \delta _{m,0}`$ (2.27)
for $`mZ`$, where $`Y_\alpha (\omega ,z)=_{mZ}L_\alpha (m)z^{m2}`$. Since $`L(0)`$ and $`\alpha (0)`$ act semisimply on $`V`$ with integral eigenvalues, $`L_\alpha (0)`$ acts semisimply on $`V`$ with eigenvalues in $`\frac{1}{2}\gamma +Z`$. That is, $`L(0)`$ acts semisimply on $`V^{(\alpha )}`$ with eigenvalues in $`\frac{1}{2}\gamma +Z`$. $`\mathrm{}`$
Since any irreducible weak module is an ordinary module for a regular vertex operator algebra, from Theorem 2.13 we immediately have:
###### Corollary 2.15
Let $`V`$ be a regular vertex operator algebra and let $`\alpha V_{(1)}`$ be given as in Theorem 2.13. Then $`V^{(\alpha )}`$ is a simple current. In particular, this is true when $`V`$ is a tensor product from the following algebras:
$$L_g(\mathrm{},0),V_L,L_g(\mathrm{},0)V_L,$$
where $`\mathrm{}`$ is a positive integer and $`L`$ is a positive-definite even lattice.$`\mathrm{}`$
The following result was obtained in \[Li4\]:
###### Proposition 2.16
Let $`L`$ be a positive definite even lattice. (1) For $`\beta L^o`$, as a $`V_L`$-module, $`V_L^{(\beta )}`$ is isomorphic to $`V_{L+\beta }`$. (2) Every irreducible $`V_L`$-module is a simple current. (3) The Verlinde algebra is canonically isomorphic to the group algebra $`C[L^o/L]`$.
Previously, intertwining operators for $`V_L`$ were explicitly constructed, fusion rules were calculated and (3) was proved in \[DL\]. (Of course, (3) implies (2).)
Though vertex operator algebra $`M_𝐡(1,0)`$ is not regular, the same proof of Proposition 2.16 gives the following result:
###### Proposition 2.17
(1) For $`h𝐡`$, as an $`M_𝐡(1,0)`$-module, $`M_𝐡(1,0)^{(h)}`$ is isomorphic to $`M_𝐡(1,h)`$. (2) Every irreducible $`M_𝐡(1,0)`$-module is a simple current. (3) The Verlinde algebra is canonically isomorphic to the group algebra $`C[𝐡]`$. $`\mathrm{}`$
###### Remark 2.18
Suppose that $`V=V^1V^2`$ is a tensor product vertex operator algebra. Then $`V_{(1)}^1`$ and $`V_{(1)}^2`$ are canonical subspaces of $`V_{(1)}`$. Suppose $`\alpha =\alpha ^{}+\alpha ^{\prime \prime }`$ where $`\alpha ^1V_{(1)}^1,a^2V_{(1)}^2`$. Then $`\alpha `$ satisfies (2.18) if and only if both $`\alpha ^1`$ and $`\alpha ^2`$ satisfy (2.18). Furthermore, $`\alpha (0)`$ acting on $`V`$ has only integral eigenvalues if and only if $`\alpha ^i(0)`$ acting on $`V^i`$ has only integral eigenvalues for $`i=1,2`$. Let $`W=W_1W_2`$, where $`W_1,W_2`$ are $`V_1`$ and $`V_2`$-modules, respectively. Since $`[\alpha ^1(m),\alpha ^2(n)]=0`$ for $`m,nZ`$, we have
$`\mathrm{\Delta }(\alpha ,z)=\mathrm{\Delta }(\alpha ^1,z)\mathrm{\Delta }(\alpha ^2,z).`$ (2.28)
Then we have
$`W^{(\alpha )}=W_1^{(\alpha ^1)}W_2^{(\alpha ^2)}.`$ (2.29)
Let $`g,𝐡,,`$ be as in Section 2.1. Let $`\{e_i,f_i|i=1,\mathrm{},n\}`$ be the Chevalley generators with simple roots $`\alpha _1,\mathrm{},\alpha _n`$ and simple coroots $`\alpha _1^{},\mathrm{},\alpha _n^{}`$. Let $`\lambda _i`$ ($`i=1,\mathrm{},n`$) be the fundamental weights for $`g`$. Let $`Q=_{i=1}^nZ\alpha _i`$ be the root lattice and let $`P=_{i=1}^nZ\lambda _i`$ be the weight lattice.
Let $`h^{(1)},\mathrm{},h^{(n)}𝐡`$ be the fundamental co-weights, i.e.,
$`\alpha _i(h^{(j)})=\delta _{i,j}\text{ for }i,j=1,\mathrm{},n.`$ (2.30)
Set
$`Q^{}`$ $`=`$ $`Z\alpha _1^{}+\mathrm{}+Z\alpha _n^{},`$ (2.31)
$`P^{}`$ $`=`$ $`Zh^{(1)}+\mathrm{}+Zh^{(n)},`$ (2.32)
the co-root lattice and the co-weight lattice.
Let
$`\theta ={\displaystyle \underset{i=1}{\overset{n}{}}}a_i\alpha _i`$ (2.33)
be the highest long root. Then
$`\theta (h^{(i)})=a_i\text{ for }i=1,\mathrm{},n.`$ (2.34)
We shall need to know which $`a_i`$ equal $`1`$. The following is a list for such $`a_i`$ (cf. \[K\]):
$`A_n:`$ $`a_1,\mathrm{},a_n`$
$`B_n:`$ $`a_1`$
$`C_n:`$ $`a_n`$
$`D_n:`$ $`a_1,a_{n1},a_n`$
$`E_6:`$ $`a_1,a_5`$
$`E_7:`$ $`a_6.`$ (2.35)
###### Remark 2.19
Simple roots for type $`E`$ ($`E_6`$, $`E_7`$, $`E_8`$) were numbered differently in \[H\] and \[K\]. Here, we use the numbering system of \[K\].
Let $`\mathrm{\Lambda }_0,\mathrm{},\mathrm{\Lambda }_n`$ be the fundamental weights of $`\widehat{g}`$ \[K\]. Then each $`\lambda _i`$ for $`1in`$ is naturally extended to $`\mathrm{\Lambda }_i`$. From the Dynkin diagram (\[K\], TABLE Aff 1) we find that $`a_i=1`$ if and only if the vertices $`0`$ and $`i`$ are in the same orbit under the automorphism group of the affine Dynkin diagram. We point out that if $`a_i=1`$, then $`\mathrm{\Lambda }_i`$ is of level one.
The following proposition was proved in \[Li4\] (Proposition 3.5, Remark 3.8):
###### Proposition 2.20
Let $`\mathrm{}`$ be a complex number with $`\mathrm{}h^{}`$, where $`h^{}`$ is the dual Coxeter number of $`g`$. If the coefficient $`a_i`$ of $`\alpha _i`$ in $`\theta `$ is $`1`$, then as an $`L(\mathrm{},0)`$-module
$`L(\mathrm{},0)^{(h^{(i)})}L(\mathrm{},\mathrm{}\lambda _i).`$ (2.36)
Furthermore, if $`\mathrm{}`$ is a positive integer, $`L(\mathrm{},\mathrm{}\lambda _i)`$ is a simple current for $`L(\mathrm{},0)`$.
###### Remark 2.21
It was known (\[FG\], \[F\]) that $`L(\mathrm{},\mathrm{}\lambda _i)`$ with $`a_i=1`$ are all the simple currents except the level $`2`$ simple current $`L(\mathrm{\Lambda }_7)`$ for $`g`$ of type $`E_8`$.
###### Remark 2.22
The element $`\mathrm{\Delta }(h^{(i)},z)`$ gives rise to an automorphism $`\psi _i`$ of affine Lie algebra $`\widehat{g}`$ via
$`\psi _i(Y(a,z))=Y_{h^{(i)}}(a,z)=Y(\mathrm{\Delta }(h^{(i)},z)a,z)`$ (2.37)
for $`ag`$, where $`Y(a,z)=_{nZ}a(n)z^{n1}`$ is the generating series of $`a`$. That is,
$`\psi _i(\alpha _i^{}(n))=\alpha _i^{}(n)+\delta _{n,0}\mathrm{},\psi _i(e_i(n))=e_i(n+1),\psi _i(f_i(n))=f_i(n1);`$ (2.38)
$`\psi _i(\alpha _j^{}(n))=\alpha _j^{}(n),\psi _i(e_j(n))=e_j(n),\psi _i(f_j(n))=f_j(n)\text{for }ji,nZ,`$ (2.39)
and
$`\psi _i(f_\theta (n))=f_\theta (n1)\text{for }nZ.`$ (2.40)
The vector $`\mathrm{𝟏}`$ in $`L(\mathrm{},0)^{(h^{(i)})}`$ is a highest weight vector of weight $`\mathrm{}\mathrm{\Lambda }_i`$. More general automorphisms of this type was recently studied in \[FS\]. Note that the composition of the representation on $`L(\mathrm{},0)`$ with the corresponding Dynkin diagram automorphism of $`\widehat{g}`$ also gives $`L(\mathrm{},\mathrm{}\lambda _i)`$. But $`\psi _i`$ are not Dynkin diagram automorphisms. Dynkin diagram automorphisms played important roles in \[FG\], \[F\], \[SY\] and \[FM\].
###### Remark 2.23
For $`\alpha 𝐡`$, from (2.25) we have
$`L_\alpha (m)=L(m)+\alpha (m)+{\displaystyle \frac{1}{2}}\mathrm{}\alpha ,\alpha \delta _{m,0}`$ (2.41)
for $`mZ`$, recall that $`Y_\alpha (\omega ,z)=_{mZ}L_\alpha (m)z^{m2}`$ with $`\omega `$ being the Virasoro element. Then the $`L_{h^{(i)}}(0)`$-weight of $`\mathrm{𝟏}`$ is $`\frac{1}{2}\mathrm{}h^{(i)},h^{(i)}`$. Since under the new action $`Y(\mathrm{\Delta }(h^{(i)},z),z)`$ on $`L(\mathrm{},0)`$, $`\mathrm{𝟏}`$ is still a highest weight vector for $`\widehat{g}`$, the lowest $`L(0)`$-weight of $`L(\mathrm{},\mathrm{}\lambda _i)`$ is $`\frac{1}{2}\mathrm{}h^{(i)},h^{(i)}`$. Of course, there is a formula for the lowest weight of any irreducible module (cf. \[DL\]).
From now on we assume that $`k`$ is a positive integer. Let $`h𝐡`$. Note that for any root vector $`e_\alpha `$ of $`g`$ and for $`mZ`$,
$`[h(0),e_\alpha (m)]=\alpha (h)e_\alpha (m).`$ (2.42)
Since $`h(0)\mathrm{𝟏}=0`$ and $`L(k,0)`$ is generated by $`\widehat{g}`$ from $`\mathrm{𝟏}`$, using (2.42) we see that $`h(0)`$ has only integral eigenvalues on $`L(k,0)`$ if and only if $`hP^{}`$.
Define a map $`\pi `$ from $`P^{}`$ to the Verlinde algebra $`𝒱(L(k,0))`$ of $`L(k,0)`$ by
$`\pi (\alpha )=[L(k,0)^{(\alpha )}]\text{ for }\alpha P^{}.`$ (2.43)
The map $`\pi `$ naturally extends to a linear map from the group algebra $`C[P^{}]`$ to $`𝒱(L(k,0))`$. We abuse the notation $`\pi `$ for this extension also.
###### Proposition 2.24
The linear map $`\pi `$ is an algebra homomorphism.
Proof. For $`\alpha ,\beta P^{}`$, since $`[\alpha (r),\beta (s)]=0`$ for $`r,s0`$, we have
$`\mathrm{\Delta }(\alpha +\beta ,z)=\mathrm{\Delta }(\alpha ,z)\mathrm{\Delta }(\beta ,z),`$ (2.44)
recall (2.20). Then we get
$`L(k,0)^{(\alpha +\beta )}(L(k,0)^{(\alpha )})^{(\beta )}`$ (2.45)
as $`L(k,0)`$-modules. In view of Theorem 2.13 we have
$`[L(k,0)^{(\alpha )}][L(k,0)^{(\beta )}]=[(L(k,0)^{(\alpha )})^{(\beta )}]=[L(k,0)^{(\alpha +\beta )}].`$ (2.46)
Thus $`\pi `$ is an algebra homomorphism. $`\mathrm{}`$
It follows immediately that $`\pi (P^{})`$ is an abelian group. The following result gives important kernel elements of $`\pi `$ as a group homomorphism on $`P^{}`$.
###### Proposition 2.25
We have
$`[L(k,0)^{(\alpha )}]=[L(k,0)]\text{ for }\alpha Q^{}.`$ (2.47)
Furthermore, for any irreducible $`L(k,0)`$-module $`W`$, we have
$`[W^{(\alpha )}]=[W]\text{ for }\alpha Q^{}.`$ (2.48)
Proof. It suffices to prove (2.47) for $`\alpha =\alpha _i^{}`$. For $`1in`$, set $`r=\frac{2}{\alpha _i,\alpha _i}`$, a positive integer. From \[H\] or \[K\] we have
$`\alpha _i^{},\alpha _i^{}={\displaystyle \frac{4}{\alpha _i,\alpha _i}}=2r.`$ (2.49)
Since
$`[e_i(1),f_i(1)]=\alpha _i^{}(0)+e_i,f_ic=\alpha _i^{}(0)+{\displaystyle \frac{1}{2}}\alpha _i^{},\alpha _i^{}c=\alpha _i^{}(0)+rc,`$ (2.50)
$`_i:=Ce_i(1)+Cf_i(1)+C(\alpha _i^{}+rc)`$ is a subalgebra of $`\widehat{g}`$ isomorphic to $`sl(2)`$. Furthermore, $`L(k,0)`$, being an integrable $`\widehat{g}`$-module, is an integrable $`_i`$-module. Clearly, $`\mathrm{𝟏}`$ is a highest weight vector of weight $`rk`$ for $`_i`$. Then $`f_i(1)^{rk}\mathrm{𝟏}0`$. From (2.25) we have
$`L_{\alpha _i^{}}(0)=L(0)+\alpha _i^{}(0)+{\displaystyle \frac{1}{2}}k\alpha _i^{},\alpha _i^{}=L(0)+\alpha _i^{}(0)+kr,`$ (2.51)
where $`Y_{\alpha _i^{}}(\omega ,z)=_{nZ}L_{\alpha _i^{}}(n)z^{n2}`$. Then
$`L_{\alpha _i^{}}(0)f_i(1)^{rk}\mathrm{𝟏}=(L(0)+\alpha _i^{}(0)+kr)f_i(1)^{rk}\mathrm{𝟏}=(kr2kr+kr)f_i(1)^{rk}\mathrm{𝟏}=0.`$ (2.52)
That is, $`f_i(1)^{rk}\mathrm{𝟏}`$ is a non-zero element of $`L(k,0)^{(\alpha _i^{})}`$ of weight zero.
On the other hand, with $`L(k,0)`$ being regular, every irreducible module, in particular, $`L(k,0)^{(\alpha _i^{})}`$, is an integrable $`\widehat{g}`$-module of level $`k`$, which is unitary. Thus the lowest weights of $`L(k,0)^{(\alpha _i^{})}`$ are nonnegative. Consequently, the lowest weight of $`L(k,0)^{(\alpha _i^{})}`$ is $`0`$. Because $`L(k,0)`$ is the only irreducible $`L(k,0)`$-module with $`0`$ being the lowest weight, we must have $`L(k,0)^{(\alpha _i^{})}L(k,0)`$ as $`L(k,0)`$-modules.
For an irreducible $`L(k,0)`$-module $`W`$, using the first part and Theorem 2.13 we get
$`[W]=[W][V]=[W][V^{(\alpha )}]=[W^{(\alpha )}]\text{ for }\alpha Q^{}.`$ (2.53)
This completes the proof.$`\mathrm{}`$
The following result generalizes Proposition 2.20 with $`L=Q`$:
###### Theorem 2.26
The algebra homomorphism $`\pi `$ gives rise to an algebra isomorphism from the group algebra $`C[P^{}/Q^{}]`$ onto $`\pi (C[P^{}])`$. Furthermore,
$`\pi (P^{})=\{[L(k,k\lambda _i)]|a_i=1\}.`$ (2.54)
Proof. In view of Proposition 2.25, $`\pi `$ gives rise to an algebra homomorphism $`\overline{\pi }`$ from $`C[P^{}/Q^{}]`$ onto $`\pi (C[P^{}])`$. From \[H\] (Section 13.1), we have
$`|P/Q|=n+1,\mathrm{\hspace{0.33em}2},\mathrm{\hspace{0.33em}2},\mathrm{\hspace{0.33em}4},\mathrm{\hspace{0.33em}3},\mathrm{\hspace{0.33em}2},\mathrm{\hspace{0.33em}1},\mathrm{\hspace{0.33em}1},\mathrm{\hspace{0.33em}1}`$ (2.55)
for $`g`$ of type $`A_{n+1},B_n,C_n,D_n,E_6,E_7,E_8,F_4,G_2`$, respectively. Note that $`B_n`$ and $`C_n`$ are dual to each other and the others are self-dual. Then $`|P^{}/Q^{}|=|P/Q|`$ for all types. On the other hand, from Proposition 2.20, all $`[L(k,k\lambda _i)]`$ for $`i`$ with $`a_i=1`$, which are distinct basis elements $`𝒱(L(k,0))`$, are images of $`\pi `$. Then it follows immediately that (2.54) holds and $`\overline{\pi }`$ is an algebra isomorphism. $`\mathrm{}`$
The group structure of $`P/Q`$ for the nontrivial cases was given in \[H\] (Exercise 4 on page 71). Then we have:
For $`A_{n+1}`$, $`P^{}/Q^{}Z/(n+1)Z`$ with $`h^{(1)}+Q^{}`$ as a generator such that
$`mh^{(1)}+Q^{}=h^{(\overline{m})}+Q^{},`$ (2.56)
where $`\overline{m}`$ is the least nonnegative residue of $`m`$ modulo $`n+1`$.
For $`D_n`$ with $`n`$ being odd, $`P^{}/Q^{}Z/4Z`$ with $`h^{(n)}+Q^{}`$ as a generator such that
$`2h^{(n)}+Q^{}=h^{(1)}+Q^{},\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}3}h^{(n)}+Q^{}=h^{(n1)}+Q^{}.`$ (2.57)
For $`D_n`$ with $`n`$ being even, $`P^{}/Q^{}Z/2Z\times Z/2Z`$ with $`h^{(n1)}+Q^{}`$ and $`h^{(n)}+Q^{}`$ as generators such that
$`h^{(1)}+Q^{}=h^{(n1)}+h^{(n)}+Q^{}.`$ (2.58)
Then with Theorem 2.26 we immediately have the following results which have been known to physicists:
###### Corollary 2.27
The following fusion algebra relations hold:
For $`A_{n+1}`$, for $`0i,jn`$,
$`[L(k,k\lambda _i)][L(k,k\lambda _j)]=[L(k,k\lambda _{\overline{i+j}})].`$ (2.59)
For $`B_n`$,
$`[L(k,k\lambda _1)][L(k,k\lambda _1)]=[L(k,0)].`$ (2.60)
For $`C_n`$,
$`[L(k,k\lambda _n)][L(k,k\lambda _n)]=[L(k,0)].`$ (2.61)
For $`D_n`$ with odd $`n`$,
$`[L(k,k\lambda _n)]^2=[L(k,k\lambda _1)],[L(k,k\lambda _n)]^3=[L(k,k\lambda _{n1})],[L(k,k\lambda _n)]^4=[L(k,0)].`$ (2.62)
For $`D_n`$ with even $`n`$,
$`[L(k,k\lambda _1)]^2=[L(k,k\lambda _{n1})]^2=[L(k,k\lambda _n)]^2=[L(k,0)],`$ (2.63)
$`[L(k,k\lambda _{n1})][L(k,k\lambda _n)]=[L(k,k\lambda _1)].`$ (2.64)
For $`E_6`$,
$`[L(k,k\lambda _1)]^2=[L(k,k\lambda _5)],[L(k,k\lambda _1)]^3=[L(k,0)].`$ (2.65)
For $`E_7`$,
$`[L(k,k\lambda _6)]^2=[L(k,0)].\mathrm{}`$ (2.66)
We shall need the number $`h^{(i)},h^{(i)}`$ and the explicit expression of $`h^{(i)}`$ in terms of $`\alpha _j^{}`$. The expression of each $`\lambda _i`$ in terms of simple roots $`\alpha _j`$ was known (\[H\], page 69, Table 1). Suppose that for $`1in`$,
$`\lambda _i=a_{i1}\alpha _1+\mathrm{}+a_{in}\alpha _n,`$ (2.67)
$`h^{(i)}=b_{i1}\alpha _1^{}+\mathrm{}+b_{in}\alpha _n^{}.`$ (2.68)
Then
$`a_{ij}=\lambda _i(h^{(j)})=b_{ji}.`$ (2.69)
Furthermore, from \[H\] we get
$`h^{(i)},\alpha _j^{}=h^{(i)},t_{\alpha _j}{\displaystyle \frac{2}{\alpha _j,\alpha _j}}=\alpha _j(h^{(i)}){\displaystyle \frac{2}{\alpha _j,\alpha _j}}=\delta _{i,j}{\displaystyle \frac{2}{\alpha _j,\alpha _j}}.`$ (2.70)
Then
$`h^{(i)},h^{(i)}=b_{ii}{\displaystyle \frac{2}{\alpha _i,\alpha _i}}=a_{ii}{\displaystyle \frac{2}{\alpha _i,\alpha _i}}.`$ (2.71)
With a glance of Table Aff in \[K\] we see that if $`a_i=1`$, $`\alpha _i`$ is a long root, hence $`\alpha _i,\alpha _i=2`$. Therefore, we have obtained:
###### Lemma 2.28
If $`\lambda _i=a_{i1}\alpha _1+\mathrm{}+a_{in}\alpha _n`$ for $`1in`$, then
$`h^{(i)}=a_{1i}\alpha _i^{}+\mathrm{}+a_{ni}\alpha _i^{},`$ (2.72)
$`h^{(i)},h^{(i)}=a_{ii}{\displaystyle \frac{2}{\alpha _j,\alpha _j}}.`$ (2.73)
In particular, if $`a_i=1`$, we have $`h^{(i)},h^{(i)}=a_{ii}`$. $`\mathrm{}`$
## 3 Extension of vertex operator algebras by simple currents
We shall first recall some of the results from \[DLM2\] on $`V[L]`$, an extension of a certain vertex operator algebra $`V`$ by simple currents $`V^{(\alpha )}`$ parametrized by $`\alpha L`$, where $`L`$ is a lattice carrying an intrinsic structure of $`V`$. We then extend and refine those results. In Section 3.3, we concentrate a special class of $`V[L]`$. We classify all irreducible $`V[L]`$-modules, prove a complete reducibility theorem and we give a formula of fusion rules in the category of $`V[L]`$-modules in terms of fusion rules in the category of $`V`$-modules.
### 3.1 Extension of algebras
We shall first establish some basic assumptions which will remain in force throughout this section.
Let $`V`$ be a simple vertex operator algebra. As in Section 2, one may think of $`V`$ as a tensor product from the following vertex operator algebras:
$$L_g(\mathrm{},0),M_𝐡(1,0),V_L.$$
Let $`H`$ be a (finite-dimensional) subspace of $`V_{(1)}`$ satisfying the following conditions:
$`L(n)h=\delta _{n,0}h,h(n)h^{}=B(h,h^{})\delta _{n,1}\mathrm{𝟏}\text{for }nZ_+,h,h^{}H,`$ (3.1)
where $`Y(h,z)=h(n)z^{n1}`$ for $`hH`$, and $`B(,)`$ is assumed to be a nondegenerate symmetric bilinear form on $`H`$. We then identify $`H`$ with its dual $`H^{}`$. We also assume that for any $`hH`$, $`h(0)`$ acts semisimply on $`V`$. Then
$`V=_{\alpha H}V^{(0,\alpha )},\text{ where }V^{(0,\alpha )}=\{vV|h(0)v=B(\alpha ,h)v\text{for }hH\}.`$ (3.2)
Set
$`P=\{\alpha H|V^{(0,\alpha )}0\}.`$ (3.3)
As $`V`$ is simple, $`P`$ is a subgroup of $`H`$ (cf. \[LX\]). Then $`P`$ equipped with the bilinear form $`B`$ is a lattice.
Let $`L`$ be a subgroup of $`H`$ such that for each $`\alpha L`$, $`\alpha (0)`$ acting on $`V`$ has only integral eigenvalues. This amounts to $`LP^o`$, where
$`P^o=\{hH|B(h,\alpha )Z\text{ for }\alpha P\}`$ (3.4)
is the dual lattice of $`P`$. Note that the rank of $`P`$ may be less than $`dimH`$.
Let $`W`$ be a $`V`$-module. By Proposition 2.9, for $`\alpha L`$, we have a (weak) $`V`$-module
$`(W^{(\alpha )},Y_\alpha (,z)):=(W,Y(\mathrm{\Delta }(\alpha ,z),z)).`$ (3.5)
For convenience, we reformulate the construction of the $`V`$-module $`W^{(a)}`$ as follows: Set
$`W^{(\alpha )}=Ce^\alpha WW(\text{linearly}),`$ (3.6)
where $`e^\alpha `$ is a symbol for now and $`Ce^\alpha `$ is a one-dimensional vector space with $`e^\alpha `$ as a pre-chosen basis element. Then define
$`Y_\alpha (v,z)(e^\alpha w)=e^\alpha Y(\mathrm{\Delta }(\alpha ,z)v,z)w\text{ for }vV,wW.`$ (3.7)
Set
$`W[L]=_{\alpha L}W^{(\alpha )}=C[L]W,`$ (3.8)
equipped with the direct sum $`V`$-module structure. Now, the symbol $`e^\alpha `$ in the definition of $`W^{(\alpha )}`$ is considered as an element of the group algebra $`C[L]`$.
For $`\alpha L`$, we define a linear endomorphism $`\psi _\alpha `$ <sup>2</sup><sup>2</sup>2The map $`\psi _\alpha `$ here is the inverse of the map $`\psi _\alpha `$ in \[DLM2\] of $`W[L]`$ by
$`\psi _\alpha (e^\beta w)=e^{\alpha +\beta }w\text{ for }\beta L,wW.`$ (3.9)
Then
$`\psi _0=1,\psi _{\alpha +\beta }=\psi _\alpha \psi _\beta \text{ for }\alpha ,\beta L.`$ (3.10)
That is, $`\psi `$ gives rise to a representation of $`L`$ on $`W[L].`$
For $`\alpha L`$, we set (\[LW\], \[FLM\])
$`E^\pm (\alpha ,z)=\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{h(\pm n)}{\pm n}}z^n\right).`$ (3.11)
Then
$`\mathrm{\Delta }(\alpha ,z)=z^{\alpha (0)}E^+(\alpha ,z).`$ (3.12)
Next, we extend the domain of $`Y_\alpha `$ from $`V`$ to $`V[L]`$.
###### Definition 3.1
For $`uV^{(\alpha )},vV^{(\beta )}`$ with $`\alpha ,\beta L`$, we define
$$Y_\alpha (u,z)v=\psi _{\alpha +\beta }E^{}(\alpha ,z)Y(\psi _\alpha \mathrm{\Delta }(\beta ,z)u,z)\mathrm{\Delta }(\alpha ,z)\psi _\beta (v)V^{(\alpha +\beta )}\{z\}.$$
(3.13)
We then define a linear map $`\stackrel{~}{Y}(,z)`$ from $`V[L]`$ to $`(\mathrm{End}V[L])\{z\}`$ via $`\stackrel{~}{Y}(u,z)=Y_\alpha (u,z)`$ for $`uV^{(\alpha )}`$.
Note that the $`E^{}(\alpha ,z)`$ defined in \[DLM2\] is the $`E^{}(\alpha ,z)`$ defined in (3.11) (\[LW\], \[FLM\]). Then the definition of $`Y_\alpha `$ is is exactly the same as the one defined in \[DLM2\].
For $`\alpha P,hH`$, set
$`V^{(\alpha ,h)}=\psi _\alpha (V^{(0,h)}).`$ (3.14)
Then
$`V^{(\alpha )}=_{hP}V^{(\alpha ,h)}.`$ (3.15)
###### Definition 3.2
We define $`C`$-valued functions $`\eta `$ and $`C`$ on $`(L\times H)\times (L\times H)`$ by
$$\eta ((\alpha _1,h_1),(\alpha _2,h_2))=B(\alpha _1,\alpha _2)B(\alpha _1,h_2)B(\alpha _2,h_1)C,$$
(3.16)
$$C((\alpha _1,h_1),(\alpha _2,h_2))=e^{(B(\alpha _1,h_2)B(\alpha _2,h_1))\pi i}C^\times $$
(3.17)
for $`(\alpha _i,h_i)L\times H,i=1,2`$.
Then we have (\[DLM2\], Theorem 3.5):
###### Theorem 3.3
Let $`uV^{(\alpha ,h_1)},vV^{(\beta ,h_2)},wV^{(\gamma ,h_3)}`$ with $`\alpha ,\beta ,\gamma L,h_1,h_2,h_3P`$. Then
$`z_0^1\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)\left({\displaystyle \frac{z_1z_2}{z_0}}\right)^{\eta ((\alpha ,h_1),(\beta ,h_2))}\stackrel{~}{Y}(u,z_1)\stackrel{~}{Y}(v,z_2)w`$ (3.18)
$``$ $`C((\alpha ,h_1),(\beta ,h_2))z_0^1\delta \left({\displaystyle \frac{z_2z_1}{z_0}}\right)\left({\displaystyle \frac{z_2z_1}{z_0}}\right)^{\eta ((\alpha ,h_1),(\beta ,h_2))}\stackrel{~}{Y}(v,z_2)\stackrel{~}{Y}(u,z_1)w`$
$`=`$ $`z_2^1\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)\left({\displaystyle \frac{z_2+z_0}{z_1}}\right)^{\eta ((\alpha ,h_1),(\gamma ,h_3))}\stackrel{~}{Y}(\stackrel{~}{Y}(u,z_0)v,z_2)w.`$
Furthermore, for all $`vV[L]`$,
$`[L(1),\stackrel{~}{Y}(v,z)]={\displaystyle \frac{d}{dz}}\stackrel{~}{Y}(v,z).`$ (3.19)
Now we express $`\stackrel{~}{Y}`$ more explicitly.
###### Lemma 3.4
For $`\alpha ,\beta L,u,vV`$, we have
$`\stackrel{~}{Y}(e^\alpha u,z)(e^\beta v)=e^{\alpha +\beta }z^{B(\alpha ,\beta )}E^{}(\alpha ,z)Y(\mathrm{\Delta }(\beta ,z)u,z)E^+(\alpha ,z)(z)^{\alpha (0)}v.`$ (3.20)
In particular,
$`\stackrel{~}{Y}(e^\alpha \mathrm{𝟏},z)(e^\beta v)=e^{\alpha +\beta }z^{B(\alpha ,\beta )}E^{}(\alpha ,z)E^+(\alpha ,z)(z)^{\alpha (0)}v.`$ (3.21)
Proof. From Lemma 3.2 of \[DLM2\] we have
$`\psi _\alpha \mathrm{\Delta }(\beta ,z)=z^{B(\alpha ,\beta )}\mathrm{\Delta }(\beta ,z)\psi _a.`$ (3.22)
Note that the map $`\psi _\alpha `$ defined here is the map $`\psi _\alpha `$ defined in \[DLM2\]. Then (3.20) follows immediately. $`\mathrm{}`$
###### Remark 3.5
Let $`G=L\times P`$ be the product group. Suppose that there exists a positive integer $`T`$ such that $`\eta `$ restricted on $`G`$ is $`\frac{1}{T}Z/2Z`$-valued. The original theorem states that $`(V[L],\mathrm{𝟏},\omega ,\stackrel{~}{Y},T,G,\eta (,),C(,))`$ is a generalized vertex algebra in the sense of \[DL\]. This result is similar to Theorem 5.1 of \[DL\], which states that if $`L`$ is a rational lattice, then $`V_L`$ has a canonical generalized vertex algebra structure. In fact, by taking $`V=M_𝐡(1,0)`$ with $`𝐡=C_ZL`$, we have $`P=0`$ and $`V[L]=V_L`$.
In order to have vertex (super)algebras $`V[L]`$, we shall restrict ourselves to special $`L`$. We have already assumed that $`LP^o`$, or what is equivalent to, $`\alpha (0)`$ acting on $`V`$ has only integral eigenvalues for every $`\alpha L`$. Now we furthermore assume that $`(L,B)`$ is an integral lattice. Then
$`B(\alpha ,\alpha ),B(\alpha ,\beta )Z\text{ for }\alpha L,\beta P.`$ (3.23)
Recall from Lemma 2.14 that for $`\alpha L`$, the weights of $`V^{(\alpha )}`$ are contained in $`\frac{1}{2}B(\alpha ,\alpha )+Z`$. Then $`V[L]`$ is $`\frac{1}{2}Z`$-graded by $`L(0)`$-weights. Furthermore, the function $`\eta `$ restricted to $`(L\times P)\times (L\times P)`$ is $`Z`$-valued and
$`C((\alpha _1,h_1),(\alpha _2,h_2))=(1)^{B(\alpha _1,h_2)B(\alpha _2,h_1)},`$ (3.24)
(recall (3.16) and (3.17)). Then we have (\[DLM2\], (3.59)):
###### Corollary 3.6
Assume that $`LP^0`$ and $`(L,B)`$ is an integral lattice. For $`aV^{(\alpha )},bV^{(\beta )}`$ with $`\alpha ,\beta L`$, we have
$`z_0^1\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)\stackrel{~}{Y}(a,z_1)\stackrel{~}{Y}(b,z_2)(1)^{B(\alpha ,\beta )}z_0^1\delta \left({\displaystyle \frac{z_2z_1}{z_0}}\right)\stackrel{~}{Y}(b,z_2)\stackrel{~}{Y}(a,z_1)`$
$`=z_2^1\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)\stackrel{~}{Y}(\stackrel{~}{Y}(a,z_0)b,z_2).`$ (3.25)
###### Remark 3.7
Set
$`L^e=\{\alpha L|B(\alpha ,\alpha )2Z\}.`$ (3.26)
Clearly, $`L^e`$ is a subgroup of $`L`$ of index 2. If $`L^e2L^o`$, i.e.,
$`B(\alpha ,\beta )2Z\text{ for }\alpha L^e,\beta L,`$ (3.27)
we easily see that $`V[L]=_{\beta L}V^{(\beta )}`$ is a vertex superalgebra with
$`V[L]^0=V[L^e]=_{\alpha L^e}V^{(\alpha )},V[L]^1=_{\alpha LL^e}V^{(\alpha )}.`$ (3.28)
In particular, if $`L`$ is of rank one, clearly (3.27) holds, hence $`V[L]`$ is a vertex (super)algebra. However, without assuming (3.27) $`V[L]`$ equipped with the vertex operator map $`\stackrel{~}{Y}`$ may not be a vertex superalgebra. Even if $`L`$ is even, $`V[L]`$ may not be a vertex algebra unless $`B`$ is $`2Z`$-valued. Corollary 3.13 of \[DLM2\], which states that $`V[L]`$ equipped with $`\stackrel{~}{Y}`$ is a vertex superalgebra if $`L`$ is integral (without the condition (3.27), is incorrect.
As in \[FLM\] and \[DL\] for $`V_L`$, we shall need a $`2`$-cocycle $`ϵ`$ on $`L`$. Suppose that $`(L,B)`$ is an integral lattice of finite rank. Let $`\{\alpha _1,\mathrm{},\alpha _d\}`$ be a $`Z`$-basis for $`L`$. Define $`ϵ`$ to be the (uniquely determined) $`\{\pm 1\}`$-valued multiplicative function on $`L\times L`$ such that
$`ϵ(\alpha _i,\alpha _j)=(1)^{B(\alpha _i,\alpha _j)+B(\alpha _i,\alpha _i)B(\alpha _j,\alpha _j)}\text{ if }ij,\text{ and }1\text{ otherwise}`$ (3.29)
(cf. \[DL\], \[FLM\]). Note that $`ϵ(\alpha _i,\alpha _i)=1`$ because $`B(\alpha _i,\alpha _i)+B(\alpha _i,\alpha _i)B(\alpha _i,\alpha _i)`$ is even. Then
$`ϵ(\alpha ,\beta )ϵ(\beta ,\alpha )^1=(1)^{B(\alpha ,\beta )+B(\alpha ,\alpha )B(\beta ,\beta )}\text{ for }\alpha ,\beta L.`$ (3.30)
Define a vertex operator map $`Y`$ on $`V[L]`$ by
$`Y(e^\alpha u,z)(e^\beta v)=ϵ(\alpha ,\beta )\stackrel{~}{Y}(e^\alpha u,z)(e^\beta v)`$ (3.31)
for $`\alpha ,\beta L,u,vV`$. By Lemma 3.4,
$`Y(e^\alpha u,z)(e^\beta v)`$ (3.32)
$`=`$ $`ϵ(\alpha ,\beta )e^{\alpha +\beta }z^{B(\alpha ,\beta )}E^{}(\alpha ,z)Y(\mathrm{\Delta }(\beta ,z)u,z)E^+(\alpha ,z)(z)^{\alpha (0)}v.`$
From Corollary 3.6 we immediately obtain:
###### Proposition 3.8
Suppose that $`LP^o`$ and that $`(L,B)`$ is an integral lattice of finite rank. Let $`ϵ`$ be the $`\{\pm 1\}`$-valued multiplicative function on $`L\times L`$ defined in (3.29). Then $`V[L]`$ equipped with the vertex operator map $`Y`$ defined in (3.32) is a vertex (super)algebra with the even and odd parts being defined in (3.28). $`\mathrm{}`$
###### Remark 3.9
It was proved in \[DL\] (Theorem 6.7 and Remarks 6.17 and 12.38) that if $`L`$ is an integral lattice, $`V_L`$ is a vertex superalgebra.
Since we in this paper are mainly interested in vertex operator (super)algebras, for the rest of Section 3 we shall assume that $`LP^o`$ and $`(L,B)`$ is integral, and we fix the multiplicative function $`ϵ`$.
### 3.2 Extensions of modules and intertwining operators
We continue Section 3.1 to study the extension of an irreducible $`V`$-module, following \[DLM2\]. The extension $`W[L]`$ of an irreducible $`V`$-module $`W`$ is in general a twisted $`V[L]`$-module with respect to an automorphism of $`V[L]`$. We shall also study the extension of an intertwining operator.
Let $`W`$ be an irreducible $`V`$-module. By definition, homogeneous subspaces of $`W`$ are finite-dimensional. Since $`[L(0),h(0)]=0`$ for $`hH`$, $`H`$ preserves each homogeneous subspace of $`W`$ so that there exist $`0wW,hH=H^{}`$ such that $`h(0)w=B(h,h^{})w`$ for $`h^{}H`$. Since $`H`$ acts semisimply on $`V`$ (by assumption) and $`w`$ generates $`W`$ by $`V`$ (from the irreducibility of $`W`$), $`H`$ also acts semisimply on $`W`$. For any $`hH`$, we define
$`W^{(0,h)}=\{wW|h(0)w=B(h,h^{})w\text{for }h^{}H\}.`$ (3.33)
Set
$`P(W)=\{hH|W^{(0,h)}0\}.`$ (3.34)
Since $`W`$ is irreducible, $`P(W)`$ is an irreducible $`P(V)`$-set. Then $`P(W)=h+P(V)`$ for any $`hP(W)`$.
###### Definition 3.10
Let $`W`$ be an irreducible $`V`$-module and $`hP(W)`$. Define a linear endomorphism $`\sigma _W`$ of $`V[L]`$ by
$`\sigma _W(a)=e^{2\pi iB(\alpha ,h)}a\text{ for }aV^{(\alpha )}\text{ with }\alpha L.`$ (3.35)
Because $`LP(V)^o`$ and $`P(W)=h+P(V)`$, $`\sigma _W`$ is well defined, i.e., it does not depend on the choice of $`h`$.
###### Lemma 3.11
The defined linear endomorphism $`\sigma _W`$ of $`V[L]`$ is an automorphism of the vertex (super)algebra and $`\sigma _W=1`$ if and only if $`\alpha (0)`$ has only integral eigenvalues on $`W`$ for every $`\alpha L`$, i.e., $`P(W)L^o`$. Furthermore, if $`V`$ is finitely generated, $`\sigma _W`$ is of finite order if and only if $`\alpha (0)`$ has rational eigenvalues on $`W`$ for every $`\alpha L`$, or equivalently, $`B(\alpha ,h)Q`$ for $`\alpha L`$.
Proof. In view of (3.13), clearly, $`\sigma _W`$ is an automorphism of the vertex (super)algebra and $`\sigma _W=1`$ if and only if $`\alpha (0)`$ has only integral eigenvalues on $`W`$ for every $`\alpha L`$. Furthermore, when $`V`$ is finitely generated, $`\sigma _W`$ is of finite order if and only if $`\alpha (0)`$ has rational eigenvalues on $`W`$ for every $`\alpha L`$, or equivalently, $`B(\alpha ,h)Q`$ for $`\alpha L`$. $`\mathrm{}`$
Recall that $`W[L]=C[L]W=_{\alpha L}W^{(\alpha )}`$.
###### Definition 3.12
Let $`W`$ be an irreducible $`V`$-module. For $`aV^{(\alpha )},wW^{(\beta )},\alpha ,\beta L`$, we define $`Y_{W[L]}(a,z)wW^{(\alpha +\beta )}\{z\}`$ by
$`Y_{W[L]}(a,z)w=ϵ(\alpha ,\beta )\psi _{\alpha +\beta }E^{}(\alpha ,z)Y(\psi _\alpha \mathrm{\Delta }(\beta ,z)a,z)\mathrm{\Delta }(\alpha ,z)\psi _\beta (w)`$ (3.36)
(cf. (3.13)).
In terms of the notion of twisted module as defined in \[Le\], \[D2\] and \[FFR\] we have (\[DLM2\], Theorem 3.6, Corollary 3.14):
###### Proposition 3.13
Let $`W`$ be an irreducible $`V`$-module such that $`\alpha (0)`$ has rational eigenvalues on $`W`$ for every $`\alpha L`$. Then the following twisted Jacobi identity holds for $`uV^{(\alpha )},vV^{(\beta )},wW^{(\gamma ,h)}`$,
$`z_0^1\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)Y(u,z_1)Y(v,z_2)w(1)^{B(\alpha ,\alpha )B(\beta ,\beta )}z_0^1\delta \left({\displaystyle \frac{z_2z_1}{z_0}}\right)Y(v,z_2)Y(u,z_1)w`$ (3.37)
$`=`$ $`z_2^1\left({\displaystyle \frac{z_1z_0}{z_2}}\right)^{B(\alpha ,h)}\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)Y(Y(u,z_0)v,z_2)w.`$
Moreover, $`W[L]`$ is a $`\sigma _W`$-twisted $`V[L]`$-module. In particular, if $`P(W)L^o`$, i.e., $`\alpha (0)`$ acting on $`W`$ has only integral eigenvalues for $`\alpha L`$, $`W[L]`$ is an (untwisted) $`V[L]`$-module.
Next, we prove a functorial property.
###### Proposition 3.14
Let $`\sigma `$ be an automorphism of $`V[L]`$ of finite order such that
$`\sigma (V^{(\alpha )})=V^{(\alpha )}\text{ for }\alpha L,`$ (3.38)
$`\sigma (v)=v\text{ for }vV=V^{(0)}.`$ (3.39)
Let $`M`$ be a $`\sigma `$-twisted weak $`V[L]`$-module which is a direct sum of irreducible $`V`$-modules and let $`W`$ be an irreducible $`V`$-submodule of $`M`$. Then $`\sigma =\sigma _W`$ and there is a canonical $`V[L]`$-homomorphism $`\varphi [L]`$ from $`W[L]`$ to $`M`$ extending the embedding $`\varphi `$ of $`W`$ into $`M`$.
Proof. The following arguments are essentially the ones of \[DLM2\], Corollary 3.15 and Lemma 4.3.
Note that for $`\alpha L`$, $`Y_{W[L]}`$ restricted to $`V^{(\alpha )}\times W`$ is a nonzero intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W^{(\alpha )}}{V^{(\alpha )}W}\right)`$. Since $`Y_M(,z)\varphi `$ restricted to $`V^{(\alpha )}\times W`$ is an intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{M}{V^{(\alpha )}W}\right)`$ and $`[V^{(\alpha )}][W]=[W^{(\alpha )}]`$, it follows from Schur lemma (cf. \[FHL\]) that there exists a unique $`V`$-homomorphism $`\varphi _\alpha `$ from $`W^{(\alpha )}`$ to $`M`$ such that
$`Y_M(a,z)w=\varphi _\alpha (Y_{W[L]}(a,z)w)`$ (3.40)
for $`aV^{(\alpha )},wW`$. From the definition of a twisted module we have
$`z^{r_\alpha }Y_M(a,z)wM((z)),\varphi _\alpha (z^{B(\alpha ,h)}Y_{W[L]}(a,z)w)M((z)).`$ (3.41)
Then it follows from (3.40) that $`r_\alpha B(\alpha ,h)Z`$, hence $`\sigma (a)=\sigma _W(a)`$. Thus $`\sigma =\sigma _W`$.
Define a $`V`$-homomorphism $`\varphi [L]`$ from $`W[L]`$ to $`M`$ by $`\varphi [L]=\varphi _\alpha `$ on $`W^{(\alpha )}`$ for $`\alpha L`$. Now we show that $`\varphi [L]`$ is a $`V[L]`$-homomorphism.
Let $`wW,aV^{(\alpha )},bV^{(\beta )}`$ with $`\alpha ,\beta L`$. Let $`k_0`$ be a positive integer such that
$`(z_0+z_2)^{k_0+B(\alpha ,h)}Y_{W[L]}(a,z_0+z_2)Y_{W[L]}(b,z_2)w`$
$`=`$ $`(z_2+z_0)^{k_0+B(\alpha ,h)}Y_{W[L]}(Y(a,z_0)b,z_2)w,`$
$`(z_0+z_2)^{k_0+B(\alpha ,h)}Y_M(a,z_0+z_2)Y_M(b,z_2)\varphi (w)`$
$`=`$ $`(z_2+z_0)^{k_0+B(\alpha ,h)}Y_M(Y(a,z_0)b,z_2)\varphi (w).`$ (3.43)
Then using (3.40) we get
$`(z_0+z_2)^{k_0+B(\alpha ,h)}\varphi _{\alpha +\beta }Y_{W[L]}(a,z_0+z_2)Y_{W[L]}(b,z_2)w`$ (3.44)
$`=`$ $`(z_2+z_0)^{k_0+B(\alpha ,h)}\varphi _{\alpha +\beta }Y_{W[L]}(Y(a,z_0)b,z_2)w`$
$`=`$ $`(z_2+z_0)^{k_0+B(\alpha ,h)}Y_M(Y(a,z_0)b,z_2)\varphi (w)`$
$`=`$ $`(z_0+z_2)^{k_0+B(\alpha ,h)}Y_M(a,z_0+z_2)Y_M(b,z_2)\varphi (w)`$
$`=`$ $`(z_0+z_2)^{k_0+B(\alpha ,h)}Y_M(a,z_0+z_2)\varphi _\beta Y_M(b,z_2)w.`$
Multiplying both sides by $`(z_0+z_2)^{k_0B(\alpha ,h)}`$ we get
$`\varphi _{\alpha +\beta }Y_{W[L]}(a,z_0+z_2)Y_{W[L]}(b,z_2)w=Y_M(a,z_0+z_2)\varphi _\beta Y_M(b,z_2)w,`$ (3.45)
that is,
$`\varphi [L]Y(a,z_0+z_2)Y_{W[L]}(b,z_2)w=Y_M(a,z_0+z_2)\varphi [L]Y_M(b,z_2)w.`$ (3.46)
Since $`W^{(\beta )}`$ is linearly spanned by $`b_nW`$ for $`bV^{(\beta )},nZ`$, we have
$`\varphi [L](Y_{W[L]}(a,z)u)=Y_M(a,z)u`$ (3.47)
for $`aV^{(\alpha )},uW^{(\beta )}`$. Thus $`\varphi [L]`$ is a $`V[L]`$-homomorphism. $`\mathrm{}`$
Our next result gives a characterization of the equivalence relation on (twisted) $`V[L]`$-modules $`W[L]`$ in terms of the equivalence of $`V`$-modules:
###### Proposition 3.15
Let $`W_1`$ and $`W_2`$ be irreducible $`V`$-modules on which $`\alpha (0)`$ has rational eigenvalues for each $`\alpha L`$. Then $`\sigma _{W_1}=\sigma _{W_2}`$ and $`W_1[L]W_2[L]`$ if and only if $`W_2W_1^{(\alpha )}`$ for some $`\alpha L`$.
Proof. The “only if” part is clear. Note that $`\sigma _{W^{(\alpha )}}=\sigma _W`$ for any irreducible $`V`$-module $`W`$ and $`\alpha L`$ because $`P(W^{(\alpha )})=\alpha +P(W)`$. Assume $`W_2W_1^{(\alpha )}`$ for some $`\alpha L`$. Then $`\sigma _{W_1}=\sigma _{W_2}`$. Let $`\varphi `$ be a $`V`$-isomorphism from $`W_2`$ to $`W_1^{(\alpha )}W_1[L]`$. It follows from Proposition 3.14 that $`\varphi `$ extends to a $`V[L]`$-homomorphism $`\varphi [L]`$ from $`W_2[L]`$ into $`W_1[L]`$ with $`\varphi [L](W_2^{(\beta )})=W_1^{(\alpha +\beta )}`$ for $`\beta L`$. With each $`W^{(\beta )}`$ being an irreducible $`V`$-module, $`\varphi [L]`$ is an isomorphism. $`\mathrm{}`$
Next, we shall extend an intertwining operator $`I`$ in the category of $`V`$-modules to an intertwining operator $`I[L]`$ in the category of $`V[L]`$-modules.
###### Definition 3.16
Let $`W_1,W_2`$ and $`W_3`$ be irreducible $`V`$-modules and $`I`$ be an intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$. We define a linear map
$`I[L]:W_1[L](\mathrm{Hom}(W_2[L],W_3[L]))\{z\}`$ (3.48)
by
$`I[L](a,z)w=ϵ(\alpha ,\beta )\psi _{\alpha +\beta }E^{}(\alpha ,z)I(\psi _\alpha \mathrm{\Delta }(\beta ,z)a,z)\mathrm{\Delta }(\alpha ,z)\psi _\beta (w)`$ (3.49)
(cf. (3.13) and (3.36)) for $`aW_1^{(\alpha )},wW_2^{(\beta )}`$ with $`\alpha ,\beta L`$.
The same proof of Theorem 3.5 in \[DLM2\] gives:
$`I[L](L(1)a,z)={\displaystyle \frac{d}{dz}}I[L](a,z)\text{ for }aW_1[L],`$ (3.50)
and
$`z_0^1\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)\left({\displaystyle \frac{z_1z_2}{z_0}}\right)^{\eta ((\alpha ,h),(\beta ,h_1))}Y_W(a,z_1)I[L](b,z_2)u`$ (3.51)
$``$ $`C((\alpha ,h),(\beta ,h_1))z_0^1\delta \left({\displaystyle \frac{z_2z_1}{z_0}}\right)\left({\displaystyle \frac{z_2z_1}{z_0}}\right)^{\eta ((\alpha ,h),(\beta ,h_1))}I[L](b,z_2)Y_W(a,z_1)u`$
$`=z_2^1\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)\left({\displaystyle \frac{z_2+z_0}{z_1}}\right)^{\eta ((\alpha ,h),(\gamma ,h_2))}I[L](Y(a,z_0)b,z_2)u`$
for $`aV^{(\alpha ,h)},bW_1^{(\beta ,h_1)},uW_2^{(\gamma ,h_2)}`$, with $`\alpha ,\beta ,\gamma L,hP,h_1P(W_1),h_2P(W_2)`$. If the extensions $`W_i[L]`$ are (untwisted) $`V[L]`$-modules, then $`P(W_i)P^o`$, so that $`\eta `$ and $`C`$ have integer values. Then we conclude:
###### Proposition 3.17
Let $`W_1,W_2,W_3`$ be irreducible $`V`$-modules such that for $`\alpha L`$, $`\alpha (0)`$ has only integral eigenvalues on $`W_i`$ for $`i=1,2,3`$, or what is equivalent to, the extensions $`W_i[L]`$ are (untwisted) $`V[L]`$-modules. Let $`I`$ be an intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$ in the category of $`V`$-modules. Then $`I[L]`$ is an intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W_3[L]}{W_1[L]W_2[L]}\right)`$ in the category of $`V[L]`$-modules. $`\mathrm{}`$
Note that various $`V`$-submodules $`V^{(\alpha )}`$ of $`V[L]`$ may be $`V`$-isomorphic to each other. It was proved in \[DLM2\] that $`V[L]`$ contains an ideal $`I`$ such that each irreducible $`V`$-module $`V^{(\alpha )}`$ is of multiplicity-one in the quotient algebra $`V[L]`$ modulo $`I`$. In the following we present an abstract reformulation of this result.
Consider an (abstract) vertex operator (super)algebra $`U=_{gG}V^g`$ graded by a (finite or infinite) abelian group $`G`$ satisfying the following conditions:
(C1) $`V^0`$ is a vertex operator subalgebra and $`V^g`$ for $`gG`$ are simple currents for $`V^0`$.
(C2) For $`uV^g,vV^h`$ with $`g,hG`$, $`u_nvV^{g+h}`$ for $`nZ`$.
(C3) For $`g,hG`$, $`u_nv0`$ for some $`uV^g,vV^h,nZ`$.
It is easy to see that under these conditions, $`U`$ is a simple $`G`$-graded algebra, i.e., there is no nontrivial $`G`$-graded ideal. From Conditions (2) and (3), $`Y`$ restricted to $`V^g\times V^h`$ is a nonzero intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{V^{g+h}}{V^gV^h}\right)`$. Then from Condition (1) we have
$`[V^g][V^h]=[V^{g+h}]\text{ for }g,hG.`$ (3.52)
Set
$`G_0=\{gG|V^gV^0\text{ as }V^0\text{-modules }\}.`$ (3.53)
Using (3.52) by a routine argument we easily get (cf. \[DLM2\], Lemma 3.7):
###### Lemma 3.18
The defined subset $`G_0`$ of $`G`$ is a subgroup and for $`g,hG`$, $`[V^g]=[V^h]`$ if and only if $`ghG_0`$.$`\mathrm{}`$
For each $`g_0G_0`$, fix a $`V^0`$-isomorphism $`f_{g_0}`$ from $`V^0`$ to $`V^{g_0}`$. We particularly define $`f_0=1`$. Let $`g_0G_0,hG`$. Then $`Y(,z)f_{g_0}`$ is a nonzero intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{V^{g+h}}{V^hV^0}\right)`$. On the other hand, for any $`V^0`$-isomorphism $`\psi `$ from $`V^h`$ to $`V^{g+h}`$, $`\psi Y(,z)`$ is a nonzero intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{V^{g+h}}{V^hV^0}\right)`$. Because $`V^h,V^0,V^{g+h}`$ are simple currents and
$$[V^{g+h}]=[V^g][V^h]=[V^0][V^h]=[V^h],$$
there exists a unique $`V^0`$-isomorphism $`f_{g_0,h}`$ from $`V^h`$ to $`V^{g_0+h}`$ such that
$`f_{g_0,h}(Y(u,z)v)=Y(u,z)f_{g_0}(v)\text{ for }vV^0,uV^h.`$ (3.54)
Define a $`V^0`$-endomorphism $`\overline{f}_{g_0}`$ of $`U`$ via $`\overline{f}_{g_0}=f_{g_0,h}`$ on $`V^h`$ for $`hG`$.
Next we define $`I`$ to be the linear span of
$`\overline{f}_{g_0}(u)u\text{ for }g_0G_0,uU.`$ (3.55)
###### Lemma 3.19
The defined subspace $`I`$ is an ideal of $`U`$ with $`IV^0=0`$. Furthermore, $`I=0`$ if and only if $`G_0=0`$.
Proof. Let $`g_0G_0,h,h^{}G`$ and let $`uV^h,vV^h^{}`$. Then
$`\overline{f}_{g_0}(u)=\mathrm{Res}_{z_2}\overline{f}_{g_0}(Y(u,z_2)\mathrm{𝟏})=\mathrm{Res}_{z_2}Y(u,z_2)\overline{f}_{g_0}(\mathrm{𝟏}).`$ (3.56)
Since
$`Y(v,z)\overline{f}_{g_0}(\mathrm{𝟏})=\overline{f}_{g_0}(Y(v,z)\mathrm{𝟏})U[[z]],`$ (3.57)
we have (cf. \[Li2\])
$`Y(v,z_0+z_2)Y(u,z_2)f_{g_0,h}(\mathrm{𝟏})=Y(Y(v,z_0)u,z_2)f_{g_0,h}(\mathrm{𝟏}).`$ (3.58)
Using (3.56)-(3.58) we obtain
$`Y(v,z_0)(\overline{f}_{g_0}(u)u)`$ (3.59)
$`=`$ $`\mathrm{Res}_{z_2}Y(v,z_0+z_2)Y(u,z_2)f_{g_0,h}(\mathrm{𝟏})Y(v,z_0)u`$
$`=`$ $`\mathrm{Res}_{z_2}Y(Y(v,z_0)u,z_2)f_{g_0,h}(\mathrm{𝟏})Y(v_0,z)u`$
$`=`$ $`f_{g_0,h}(Y(v,z_0)u)Y(v,z_0)u`$
$`=`$ $`\overline{f}_{g_0}(Y(v,z_0)u)Y(v,z_0)u.`$
It follows immediately that $`I`$ is an ideal.
Clearly, $`I0`$ if $`G_0\{0\}`$, so $`I=0`$ implies $`G_0=0`$. If $`G_0=0`$, with $`f_0=1`$ from the definition of $`I`$ we have $`I=0`$. $`\mathrm{}`$
###### Proposition 3.20
The algebra $`U`$ is simple if and only if $`G_0=0`$. Furthermore, the quotient algebra $`\overline{U}=U/I`$ is simple.
Proof. From Lemma 3.19, $`U`$ is not simple if $`G_00`$. Now we prove that $`U`$ is simple if $`G_0=\{0\}`$. In view of Lemma 3.18, all $`V^g`$ for $`gG`$ are non-isomorphic irreducible $`V^0`$-modules. Then any nonzero ideal of $`U`$ as a $`V^0`$-module must be a sum of some $`V^g`$. Then it follows from the conditions (1)-(3) that any nonzero ideal of $`U`$ must be $`U`$. That is, $`U`$ is simple.
Note that $`\overline{U}=U/I`$ is a vertex operator (super)algebra graded by group $`G/G_0`$ with all the conditions (1)-(3) being satisfied. Furthermore, $`\overline{U}`$ is a direct sum of non-isomorphic simple current $`V^0`$-modules. From Part one, $`\overline{U}`$ must be simple. $`\mathrm{}`$
Applying Proposition 3.20 to $`V[L]`$ we immediately have (cf. \[DLM2\], Corollary 3.13):
###### Corollary 3.21
Let $`V,L`$ be as before. Set
$`L_0=\{\alpha L|V^{(\alpha )}V\text{ as }V\text{-modules}\}.`$ (3.60)
Then $`V[L]`$ has an ideal $`I`$ such that $`V[L]/I`$ is simple with $`V`$ as a subalgebra and such that as a $`V`$-module
$`V/I_{\alpha S}V^{(\alpha )},`$ (3.61)
where $`S`$ is a complete set of representatives of cosets of $`L_0`$ in $`L`$. $`\mathrm{}`$
### 3.3 Multiplicity-free extension $`V[L]`$
In this subsection we consider extended vertex (super)algebra $`V[L]`$ in which each $`V`$-module $`V^{(\alpha )}`$ is multiplicity-free. We shall classify all irreducible $`V[L]`$-modules in terms of irreducible $`V`$-modules and determine the fusion rules of $`V[L]`$-modules by the fusion rules of $`V`$-modules. Our concrete examples we shall construct in Section 4 is of this type, so that all the results of this subsection apply to those examples.
Throughout this subsection we assume that for any irreducible $`V`$-module $`W`$ and for $`\alpha ,\beta L`$, $`W^{(\alpha )}W^{(\beta )}`$ as $`V`$-modules if and only if $`\alpha =\beta `$.
We shall need the following result:
###### Proposition 3.22
The vertex (super)algebra $`V[L]`$ is simple. Furthermore, if $`Y_1`$ and $`Y_2`$ are two simple vertex operator (super)algebra structures on $`V[L]`$ extending the $`V`$-module structure $`Y_V`$, then vertex operator (super)algebras $`(V[L],Y_1)`$ and $`(V[L],Y_2)`$ are isomorphic.
Proof. First, we prove that $`V[L]`$ is simple. Notice that in the definition (3.13) of the vertex operator map $`Y`$, $`\psi _{\alpha +\beta },\psi _\alpha ,\psi _\beta ,E^{}(\alpha ,z)`$, $`\mathrm{\Delta }(\beta ,z)`$ and $`\mathrm{\Delta }(\alpha ,z)`$ are invertible elements and that
$$Y(a,z)bV^{(\alpha +\beta )}\text{ for }aV^{(\alpha )},bV^{(\beta )}$$
(cf. (3.13)). Since $`V`$ is simple, $`Y(u,z)v0`$ for $`0u,vV`$ (\[DL\], Proposition 11.9, or \[FHL\], Remark 5.4.6). Then it follows from Proposition 3.20 immediately that $`V[L]`$ is simple.
For $`\alpha ,\beta L`$, because $`V^{(\alpha )}`$ and $`V^{(\beta )}`$ are simple currents, there exists $`ϵ^{}(\alpha ,\beta )C^\times `$ such that
$`Y_2(a,z)b=ϵ^{}(\alpha ,\beta )Y_1(a,z)b\text{ for }aV^{(\alpha )},bV^{(\beta )}.`$ (3.62)
It follows from weak associativity of vertex operators (cf. (3.43)) that
$`ϵ^{}(\alpha ,\beta +\gamma )ϵ^{}(\beta ,\gamma )=ϵ^{}(\alpha ,\beta )ϵ^{}(\alpha +\beta ,\gamma )`$ (3.63)
for $`\alpha ,\beta L`$. That is, $`ϵ^{}`$ is a ($`C^\times `$-valued) 2-cocycle on $`L`$. We also have
$`ϵ^{}(0,\alpha )=ϵ^{}(\alpha ,0)=1.`$ (3.64)
Since $`V=V^{(0)}`$ is even for both superalgebra structures, each structure corresponds a sublattice $`L_i`$ of $`L`$ of index 2 with $`V[L_i]`$ being the even parts.
Now we claim that $`L_1=L_2`$. Otherwise, suppose $`L_1L_2\mathrm{}`$ and let $`\beta L_1L_2`$. Then we have the following skew-symmetry
$`Y_1(a,z)b=e^{zL(1)}Y_1(b,z)a,`$ (3.65)
$`Y_2(a,z)b=e^{zL(1)}Y_2(b,z)a`$ (3.66)
for $`a,bV^{(\beta )}`$. Since both $`Y_1`$ and $`Y_2`$ extend the $`V`$-module structure, the two vertex superalgebra structures have the same Virasoro vector. Consequently,
$`ϵ^{}(\beta ,\beta )=ϵ^{}(\beta ,\beta ),`$ (3.67)
which is impossible because $`ϵ^{}(\beta ,\beta )0`$.
With $`L_1=L_2`$, using the skew-symmetry we obtain
$`ϵ^{}(\alpha ,\beta )=ϵ^{}(\beta ,\alpha )\text{ for }\alpha ,\beta L.`$ (3.68)
It follows from the proof of Propositions 5.1.2 and 5.2.3 in \[FLM\] (with $`Z/sZ`$ being replaced by $`C^\times `$) that $`ϵ^{}`$ is a $`2`$-coboundary. Then the two vertex superalgebra structures on $`V[L]`$ are equivalent. $`\mathrm{}`$
###### Remark 3.23
More generally, let $`G`$ be a (finite or infinite) abelian group and let $`V`$ be a simple vertex operator algebra and $`V[G]=_{gG}V^g`$ be a $`V`$-module with each $`V^g`$ being an irreducible $`V`$-submodule. Furthermore, assume that each $`V^g`$ is a simple current of $`V`$. Then the set of equivalence class of simple vertex operator (super)algebra structures on $`V[G]`$ which extend the $`V`$-module one-to-one corresponds to the set of equivalence classes of symmetric $`C^\times `$-valued $`2`$-cocycles of $`G`$.
Similar to Proposition 3.22 we have (cf. \[DLM2\], Lemma 4.2):
###### Proposition 3.24
Let $`W`$ be an irreducible $`V`$-module. Then $`W[L]`$ is irreducible. $`\mathrm{}`$
The following theorem gives the complete reducibility for every $`V[L]`$-module under certain conditions:
###### Theorem 3.25
Assume that there is a sublattice $`L_1`$ of $`L`$ such that $`V[L_1]`$ is regular and that every irreducible $`V[L_1]`$-module is a direct sum of irreducible $`V`$-modules. Let $`\sigma `$ be an automorphism of $`V[L]`$ such that $`\sigma `$ fixes $`V[L_1]`$ point-wise. Then any $`\sigma `$-twisted weak $`V[L]`$-module is a direct sum of irreducible $`\sigma `$-twisted $`V[L]`$-modules of type $`W[L]`$ with $`\sigma =\sigma _W`$. In particular, $`V[L]`$ is regular.
Proof. Let $`M`$ be a $`\sigma `$-twisted weak $`V[L]`$-module. Since $`\sigma `$ fixes $`V[L_1]`$ point-wise, $`M`$ is a weak $`V[L_1]`$-module. With $`V[L_1]`$ being regular, $`M`$ is a direct sum of irreducible (ordinary) $`V[L_1]`$-modules. With the assumption that each irreducible $`V[L_1]`$-module is a direct sum of irreducible $`V`$-modules, $`M`$ is a direct sum of irreducible $`V`$-modules. Let $`W`$ be an irreducible $`V`$-submodule of $`M`$. By Proposition 3.14, $`\sigma =\sigma _W`$ and there exists a $`V[L]`$-homomorphism $`\varphi [L]`$ from $`W[L]`$ to $`M`$ extending the embedding $`\varphi `$ of $`W`$ into $`M`$. Since $`W[L]`$ is $`V[L]`$-irreducible (Proposition 3.24), the $`V[L]`$-submodule of $`M`$ generated by $`W`$, which is the image of $`\varphi [L]`$, is an irreducible $`\sigma `$-twisted $`V[L]`$-module. Therefore, $`M`$ is a direct sum of irreducible $`\sigma `$-twisted $`V[L]`$-modules of type $`W[L]`$. $`\mathrm{}`$
Let $`W_1,W_2,W_3`$ be irreducible $`V`$-modules such that $`\sigma _{W_i}=1`$, or equivalently, the extensions $`W_1[L],W_2[L]`$ and $`W_3[L]`$ are $`V[L]`$-modules. For each $`\alpha L`$, with $`W_3^{(\alpha )}`$ being a $`V`$-submodule of $`W_3[L]`$, by Proposition 3.14, there is a $`V[L]`$-homomorphism $`g_\alpha `$ from $`W_3^{(\alpha )}[L]`$ to $`W_3[L]`$. Then with Proposition 3.17 we obtain a natural linear map $`f_\alpha `$ from $`I_{W_1W_2}^{W_3^{(\alpha )}}`$ to $`I_{W_1[L]W_2[L]}^{W_3[L]}`$ defined by
$`f_\alpha (I)=g_\alpha I[L].`$ (3.69)
The next result gives a precise connection between the fusion rules for $`V`$-modules and the fusion rules for $`V[L]`$-modules:
###### Theorem 3.26
Let $`W_1,W_2,W_3`$ be irreducible $`V`$-modules such that $`\sigma _{W_i}=1`$, which implies that $`W_1[L],W_2[L],W_3[L]`$ are irreducible $`V[L]`$-modules. In addition, we assume that $`V`$ is quasi-rational. Then $`f=_{\alpha L}f_\alpha `$ is a linear isomorphism from $`_{\alpha L}I_{W_1W_2}^{W_3^{(\alpha )}}`$ to $`I_{W_1[L]W_2[L]}^{W_3[L]}`$. In particular,
$`N_{W_1[L]W_2[L]}^{W_3[L]}={\displaystyle \underset{\alpha L}{}}N_{W_1W_2}^{W_3^{(\alpha )}}.`$ (3.70)
Proof. Since $`W_3^{(\alpha )}W_3^{(\beta )}`$ only if $`\alpha =\beta `$, clearly, $`f_\alpha `$ is one-to-one. On the other hand, if $`𝒴`$ is an intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W_3[L]}{W_1[L]W_2[L]}\right)`$. Then by restricting $`𝒴`$ to $`W_1\times W_2`$ we have an intertwining operator $`I`$ of type $`\left(\genfrac{}{}{0pt}{}{W_3^{(\alpha )}}{W_1W_2}\right)`$ for a unique $`\alpha L`$. It is clear that $`𝒴=I[L]`$. This completes the proof. $`\mathrm{}`$
We now describe the Verlinde algebra of $`V[L]`$ in terms of the Verlinde algebra of $`V`$ explicitly. Let $`𝒜(V)`$ be the Verlinde algebra of $`V`$. The Verlinde algebra $`𝒜(V[L])`$ with a basis $`[W[L]]`$ for $`[W]𝒜(V)`$ with $`\sigma _W=1`$. First, we have:
###### Lemma 3.27
All $`[W]`$ with $`\sigma _W=1`$ linearly span a subalgebra $`A(V,L)`$ of $`𝒜(V)`$.
Proof. Suppose that $`[W_1],[W_2]\mathrm{Irr}(V)`$ with $`\sigma _{W_1}=\sigma _{W_2}=1`$. Let $`h_1P(W_1),h_2P(W_2)`$. Then $`\sigma _{W_1}=\sigma _{W_2}=1`$ amount to $`h_1,h_2L^o`$. Let $`𝒴`$ be a nonzero intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$ for some irreducible $`V`$-module $`W_3`$. Then $`h_1+h_2P(W_3)`$ because for $`hH,w_{(i)}W_i^{(0,h_i)}`$,
$`h(0)𝒴(w_{(1)},z)w_{(2)}`$ $`=`$ $`𝒴(h(0)w_{(1)},z)w_{(2)}+𝒴(w_{(1)},z)h(0)w_{(2)}`$ (3.71)
$`=`$ $`B(h,h_1+h_2)𝒴(w_{(1)},z)w_{(2)}.`$
Since $`h_1+h_2L^o`$, $`\sigma _{W_3}=1`$. Then $`[W_3]A(V,L)`$. The proof is complete. $`\mathrm{}`$
Define a subspace $`R`$ of $`A(V,L)`$ linearly spanned by
$`[W][W^{(\alpha )}]\text{ for }\alpha L.`$ (3.72)
Then $`R`$ is an two-sided ideal of $`A(V,L)`$. Indeed, let $`W_1`$ and $`W_2`$ be irreducible $`V`$-modules with $`\sigma _{W_i}=1`$ for $`i=1,2`$. For $`\alpha L`$, by Proposition 2.10,
$`I_{W_1W}^{W_2}I_{W_1W^{(\alpha )}}^{W_2^{(\alpha )}}.`$ (3.73)
Thus
$`[W_1]([W][W^{(\alpha )}])={\displaystyle \underset{[W_2]Irr(V)}{}}N_{W_1W}^{W_2}([W_2][W_2^{(\alpha )}])R.`$ (3.74)
Since $`𝒜(V)`$ is a commutative algebra, $`R`$ is a two-sided ideal. Furthermore, by Proposition 3.15, $`\overline{[W_1]}=\overline{[W_2]}`$ in the quotient algebra $`A(V,L)/R`$ if and only if $`[W_1[L]]=[W_2[L]]`$ in $`𝒜(V[L])`$. Then in view of Theorem 3.26 we immediately have:
###### Proposition 3.28
The subspace $`R`$ is a two-sided ideal of $`A[V,L]`$ and the Verlinde algebra $`𝒜(V[L])`$ is canonically isomorphic to the quotient algebra of $`A(V,L)`$ modulo $`R`$.$`\mathrm{}`$
## 4 Extended vertex operator (super)algebras of affine types
In this section we shall specialize the vertex (super)algebra $`V[L]`$ constructed in Section 3 from a pair $`(V,L)`$ to obtain extensions of vertex operator algebras associated with affine Lie algebras $`\widehat{g}`$. In the case $`g=sl(2)`$, we shall obtain Feigin-Miwa’s extended vertex operator (super)algebras $`A_k`$ \[FM\]. To apply the results of Section 3 we need to define $`V`$ and $`L`$ explicitly and check the necessary conditions. For each type, $`V`$ will be the tensor product vertex operator algebra $`L_g(k,0)M_𝐡^{}(1,0)`$ where $`𝐡^{}`$ is a 1 or 2-dimensional vector space equipped with a nondegenerate symmetric bilinear form. When defining $`𝐡^{}`$, we follow two basic principles: (1) To include all the $`L(k,0)`$-simple currents in the construction of $`V[L]`$. (2) To make $`dim𝐡^{}`$, or equivalently, to make the rank of $`V[L]`$ as small as possible. After $`𝐡^{}`$ is chosen, we still have plenty of choices for the bilinear form on $`𝐡^{}`$. Another principle to follow is to make $`V[L]`$ as large as possible. We shall define $`𝐡^{}`$ and $`L`$ type by type.
### 4.1 A complete reducibility theorem for a certain family of $`V[L]`$
In Section 3, for a general pair $`(V,L)`$, under a certain assumption we proved a complete reducibility theorem (Theorem 3.25) for extended algebra $`V[L]`$. In this section, we shall consider a family of $`V[L]`$ such that the assumption of Theorem 3.25 holds. All the extended algebras we shall construct later belong to this family, so the complete reducibility theorem holds for all of them.
Let $`g=g_1\mathrm{}g_r`$ be a semisimple Lie algebra with a Cartan subalgebra
$`𝐡=𝐡_1+\mathrm{}+𝐡_r,`$ (4.1)
where $`g_i`$ are simple factors with Cartan subalgebras $`𝐡_i`$, equipped with the normalized Killing forms. Let $`Q^{}=Q_1^{}+\mathrm{}+Q_r^{}`$ and $`P^{}=P_1^{}+\mathrm{}+P_r^{}`$ be the coroot lattice and coweight lattice of $`g`$, respectively, where $`Q_i^{}`$ and $`P_i^{}`$ are the coroot lattice and coweight lattice of $`g_i`$.
Let $`𝐤=(k_1,\mathrm{},k_r)`$ be an $`r`$-tuple of nonnegative integers. Set
$`L_g(𝐤,0)=L_{g_1}(k_1,0)\mathrm{}L_{g_r}(k_r,0),`$ (4.2)
equipped with the standard tensor product vertex operator algebra structure. Set
$`P_𝐤=\{(\lambda ^1,\mathrm{},\lambda ^r)|\lambda ^iP_{k_i}(g_i)\},`$ (4.3)
where $`P_{k_i}(g_i)`$ stands for $`P_{k_i}`$ for the Lie algebra $`g_i`$. Then $`L_g(𝐤,0)`$ is regular \[DLM1\], i.e., every weak module is a direct sum of irreducible (ordinary) modules $`L_g(𝐤,\lambda )`$, where
$`L_g(𝐤,\lambda )=L_{g_1}(k_1,\lambda ^1)\mathrm{}L_{g_r}(k_r,\lambda ^r)`$ (4.4)
for $`\lambda =(\lambda ^1,\mathrm{},\lambda ^r)P_𝐤`$.
Let $`𝐡^{}`$ be a finite-dimensional vector space equipped with a nondegenerate symmetric bilinear form $`,`$. Associated to $`𝐡^{}`$ is the vertex operator algebra $`M_𝐡^{}(1,0)`$. Set
$`V=L_g(𝐤,0)M_𝐡^{}(1,0),`$ (4.5)
equipped with the standard tensor product vertex operator algebra structure. The algebra $`V`$ can be considered as the vertex operator algebra associated to the affine algebra of the reductive Lie algebra $`g+𝐡^{}`$.
Set
$`H=𝐡+𝐡^{}g+𝐡^{}=V_{(1)}.`$ (4.6)
Then $`H`$ is a Cartan subalgebra of the reductive Lie algebra $`g+𝐡^{}`$. Clearly, (3.1) holds and
$`B(h_1,h_2)=\delta _{i,j}k_ih_1,h_2\text{ for }h_1𝐡_i,h_2𝐡_j`$ (4.7)
because $`h_1(1)h_2=h_1(1)h_2(1)\mathrm{𝟏}=\delta _{i,j}k_ih_1,h_2\mathrm{𝟏}`$ (and $`h_1(1)h_2=B(h_1,h_2)\mathrm{𝟏}`$). We also have
$`B(h_1^{},h_2^{})=h_1^{},h_2^{}\text{ for }h_1^{},h_2^{}𝐡^{}.`$ (4.8)
For $`hH`$, we write
$`h=h^1+\mathrm{}+h^r+h^{},h=h^{\prime \prime }+h^{},`$ (4.9)
where $`h^i𝐡_i,h^{}𝐡^{},h^{\prime \prime }𝐡`$.
For $`\lambda =(\lambda ^1,\mathrm{},\lambda ^r),\gamma 𝐡^{}`$ with $`\lambda ^iP_{k_i}(g_i)`$, we set
$`W(\lambda ,\gamma )=L_{g_1}(k_1,\lambda ^1)\mathrm{}L_{g_r}(k_r,\lambda ^r)M_𝐡^{}(1,\gamma ).`$ (4.10)
Then from \[FHL\] all such $`W(\lambda ,\gamma )`$ form a complete set of representatives of equivalence classes of irreducible $`V`$-modules.
Let $`L`$ be a subgroup of $`H`$ such that $`(L,B)`$ is an integral lattice of finite rank. To have a vertex (super)algebra $`V[L]`$, we shall also need the condition that $`\alpha (0)`$ has only integral eigenvalues on $`V`$ for every $`\alpha L`$. For $`\alpha L`$, since $`\alpha ^{}(0)`$ acts as zero on $`V`$, $`\alpha (0)`$ has only integral eigenvalues on $`V`$ if and only if $`\alpha ^{\prime \prime }(0)`$ has only integral eigenvalues on $`L_g(𝐤,0)`$. Then we immediately have:
###### Lemma 4.1
For $`\alpha L`$, $`\alpha (0)`$ has only integral eigenvalues on $`V`$ if and only if $`\alpha ^{\prime \prime }P^{}`$, the coweight lattice of $`g`$.$`\mathrm{}`$
Set
$`L^{}=\{\alpha ^{}|\alpha L\},L^{\prime \prime }=\{\alpha ^{\prime \prime }|\alpha L\}.`$ (4.11)
###### Lemma 4.2
Assume that the projection of $`L`$ onto $`L^{}`$ is one-to-one. Then for $`\lambda P_𝐤,\gamma 𝐡^{}`$ and for $`\alpha ,\beta L`$, $`W(\lambda ,\gamma )^{(\alpha )}W(\lambda ,\gamma )^{(\beta )}`$ if and only if $`\alpha =\beta `$.
Proof. Because $`\mathrm{\Delta }(\alpha ,z)=\mathrm{\Delta }(\alpha ^{\prime \prime },z)\mathrm{\Delta }(\alpha ^{},z)`$ and $`\mathrm{\Delta }(\alpha ^{\prime \prime },z)=1`$ on $`M_𝐡^{}(1,\gamma )`$ and $`\mathrm{\Delta }(\alpha ^{},z)=1`$ on $`L_g(𝐤,\lambda )`$ for $`\alpha L`$, we have
$`W(\lambda ,\gamma )^{(\alpha )}=L_g(𝐤,\lambda )^{(\alpha ^{\prime \prime })}M_𝐡^{}(1,\gamma )^{(\alpha ^{})}L_g(𝐤,\lambda )^{(\alpha ^{\prime \prime })}M_𝐡^{}(1,\gamma +\alpha ^{}).`$ (4.12)
We knew $`M_𝐡^{}(1,\gamma +\alpha ^{})M_𝐡^{}(1,\gamma +\beta ^{})`$ if and only if $`\alpha ^{}=\beta ^{}`$. Then it follows immediately. $`\mathrm{}`$
Now we have:
###### Proposition 4.3
Let $`V,H`$ be defined as in (4.5) and (4.6) and let $`L`$ be a subgroup of $`H`$ such that $`(L,B)`$ is an integral lattice of finite rank. Assume that $`L^{\prime \prime }P^{}`$, the projection of $`L`$ onto $`L^{}`$ is one-to-one and $`L^{}`$ is a positive definite lattice. Then $`V[L]`$ equipped with the vertex operator map $`Y`$ defined in (3.13) is a simple vertex operator (super)algebra.
Proof. Since $`L^{\prime \prime }P^{}`$, by Lemma 4.1, for $`\alpha L`$, $`\alpha (0)`$ acting on $`V`$ has only integral eigenvalues. With Lemma 4.2, in view of Corollary 3.6 and Proposition 3.22, $`V[L]`$ is a simple vertex (super)algebra with all $`L(0)`$-weights being half integers. Now we only need to verify the two grading restrictions \[FLM\].
For $`\alpha L`$, we have
$`V^{(\alpha )}=L_g(𝐤,0)^{(\alpha ^{\prime \prime })}M_𝐡^{}(1,0)^{(\alpha ^{})}=L_g(𝐤,0)^{(\alpha ^{\prime \prime })}M_𝐡^{}(1,\alpha ^{}).`$ (4.13)
With $`L_g(𝐤,0)`$ being regular, $`L_g(𝐤,0)^{(\alpha ^{\prime \prime })}`$ is an irreducible $`L_g(𝐤,0)`$-module, which is unitary. Then the $`L(0)`$-weights of $`L_g(𝐤,0)^{(\alpha ^{\prime \prime })}`$ are nonnegative. On the other hand, the lowest weight of $`M_𝐡^{}(1,\alpha ^{})`$ is $`\frac{1}{2}\alpha ^{},\alpha ^{}`$. Thus each $`V^{(\alpha )}`$ satisfies the two grading restrictions with the lowest weight at least $`\frac{1}{2}\alpha ^{},\alpha ^{}`$. Since the projection of $`L`$ onto $`L^{}`$ is one-to-one and $`L^{}`$ is positive definite, for every $`n\frac{1}{2}Z`$, $`V_{(n)}^{(\alpha )}0`$ only for finitely many $`\alpha `$. Then the two grading restrictions follows immediately. $`\mathrm{}`$
Furthermore, we have:
###### Theorem 4.4
Let $`V,L`$ be as in Proposition 4.3 with all the assumptions. In addition we assume that $`dim𝐡^{}=\text{ rank}L^{}`$. Let $`\sigma `$ be an automorphism of $`V[L]`$ of finite order which fixes $`V`$ point-wise. Then every $`\sigma `$-twisted weak $`V[L]`$-module is a direct sum of irreducible (ordinary) $`\sigma `$-twisted $`V[L]`$-modules isomorphic to $`W(\lambda ,\gamma )[L]`$ with $`\sigma =\sigma _{W(\lambda ,\gamma )}`$. In particular, every weak $`V[L]`$-module is a direct sum of irreducible (ordinary) $`V[L]`$-modules isomorphic to $`W(\lambda ,\gamma )[L]`$ for $`\lambda P_𝐤`$, $`\gamma (L^{})^o`$ (the dual of $`L^{}`$) with
$`\lambda ^1(\alpha ^1)+\mathrm{}+\lambda ^r(\alpha ^r)+\gamma ,\alpha ^{}Z\text{ for }\alpha L.`$ (4.14)
Proof. Denote by $`o(\sigma )`$ the order of $`\sigma `$. In view of Lemma 4.2, $`V^{(\alpha )}`$ for $`\alpha L`$ are non-isomorphic irreducible $`V`$-modules. Then $`\sigma (V^{(\alpha )})=V^{(\alpha )}`$ for $`\alpha L`$ and $`\sigma `$ acts on $`V^{(\alpha )}`$ as a scalar, which is an $`o(\sigma )`$-th root of unity. Therefore, $`\sigma `$ acts trivially on $`V^{(mo(\sigma )\alpha )}`$ for $`\alpha L,mZ`$.
Because $`L_g(𝐤,0)`$ has only finitely many irreducible modules up to equivalence, there exists a positive integer $`d_1`$ such that as $`L_g(𝐤,0)`$-modules,
$`L_g(𝐤,0)^{(d_1\alpha ^{\prime \prime })}L_g(𝐤,0)\text{ for all }\alpha L.`$ (4.15)
Let $`d_2`$ be another positive integer such that $`d_2L^{}`$ is an even lattice. Set $`d=o(\sigma )d_1d_2`$ and $`L_1=dL`$. Then
$`V[L_1]=_{\alpha L}V^{(d\alpha )}.`$ (4.16)
Then $`V[L_1]`$ is a simple vertex operator subalgebra of $`V[L]`$ and $`\sigma `$ fixes $`V[L_1]`$ point-wise. Furthermore, as a $`V`$-module,
$`V^{(d\alpha )}L_g(𝐤,0)^{(d\alpha ^{\prime \prime })}M_𝐡^{}(1,d\alpha ^{})L_g(𝐤,0)M_𝐡^{}(1,d\alpha ^{})`$ (4.17)
for $`\alpha L`$, hence
$`V[L_1]L_g(𝐤,0)V_{dL^{}}.`$ (4.18)
(Here we used the fact that $`dim𝐡^{}=\text{rank}L^{}`$.) Note that $`L_g(𝐤,0)V_{dL^{}}`$ is a natural simple vertex operator algebra, which is regular. It follows from Proposition 3.22 that $`V[L_1]`$ is regular. Clearly, each irreducible $`V[L_1]`$-module is a direct sum of irreducible $`V`$-modules. Then it follows from Theorem 3.25 immediately that every $`\sigma `$-twisted weak $`V[L]`$-module is a direct sum of irreducible (ordinary) $`\sigma `$-twisted $`V[L]`$-modules of type $`W(\lambda ,\gamma )`$ with $`\sigma =\sigma _{W(\lambda ,\gamma )}`$.
From Lemma 3.11, $`\sigma _{W(\lambda ,\gamma )}=1`$ if and only if for $`\alpha L`$, $`\alpha (0)`$ has only integral eigenvalues on $`W(\lambda ,\gamma )`$. With $`(\lambda ^1,\mathrm{},\lambda ^r,\gamma )`$ being an $`H`$-weight of $`W(\lambda ,\gamma )`$, we see that $`\sigma _{W(\lambda ,\gamma )}=1`$ if and only if
$`\lambda ^1(\alpha ^1)+\mathrm{}+\lambda ^r(\alpha ^r)+\gamma ,\alpha ^{}Z\text{ for }\alpha L,`$ (4.19)
which furthermore implies that $`\gamma (L^{})^o`$ because $`\lambda P_𝐤`$. This completes the proof.$`\mathrm{}`$
###### Remark 4.5
From the proof of Theorem 4.4, one can easily see that the regularity result still holds for $`V[L]`$ if we replace $`L_g(𝐤,0)`$ by any regular vertex operator algebra $`U`$.
Because
$$W(\lambda ,\gamma )^{(\alpha )}=L(𝐤,\lambda )^{(\alpha ^{\prime \prime })}M_𝐡^{}(1,\gamma +\alpha ^{})$$
and $`M_𝐡^{}(1,\gamma )M_𝐡^{}(1,\gamma ^{})`$ if and only if $`\gamma =\gamma ^{}`$, in view of Proposition 3.15, we see that $`W(\lambda ,\gamma )[L]W(\lambda ^{},\gamma ^{})[L]`$ if and only if there is $`\alpha L`$ such that
$`\gamma ^{}=\gamma +\alpha ^{},L(𝐤,\lambda ^{})^{(\alpha ^{\prime \prime })}L(𝐤,\lambda ).`$ (4.20)
To describe explicitly the equivalence relation on the set of $`W(\lambda ,\gamma )[L]`$, or to get a complete set of equivalence classes of irreducible $`V[L]`$-modules, we need to know $`L(𝐤,\lambda ^{})^{(\alpha ^{\prime \prime })}`$ as a $`\widehat{g}`$-module. Of course, from Theorem 2.13,
$`[L(𝐤,\lambda ^{})^{(\alpha ^{\prime \prime })}]=[L(𝐤,\lambda ^{})][L(𝐤,0)^{(\alpha ^{\prime \prime })}].`$ (4.21)
Nevertheless, in view of Proposition 3.28 we immediately have:
###### Proposition 4.6
The subspace $`A`$ of $`𝒱(L(𝐤,0))C[(L^{})^o]`$, linearly spanned by
$`[L(𝐤,\lambda )]e^\gamma `$ (4.22)
for $`\lambda P_𝐤,\gamma (L^{})^o`$ satisfying (4.14), is a subalgebra. Furthermore, the Verlinde algebra $`𝒱(V[L])`$ is canonically isomorphic to the quotient algebra of $`A`$ modulo the relations:
$`[L(𝐤,\lambda )]e^\gamma [L(𝐤,\lambda )^{(\alpha ^{\prime \prime })}]e^{\gamma +\alpha ^{}}`$ (4.23)
for $`\alpha L`$.$`\mathrm{}`$
### 4.2 Extended vertex operator (super)algebras $`A_k`$ of type $`sl(n+1)`$
Starting from this subsection we shall work on the setting of Section 4.1 and we shall consider a simple Lie algebra $`g`$, i.e., $`r=1`$. For $`g`$ of each type, we take
$`V=L(k,0)M_𝐡^{}(1,0)`$ (4.24)
and we define $`A_k(g)`$ to be the extended algebra $`V[L]`$ for a certain $`L`$. We shall case by case define the pair $`(𝐡^{},,)`$ and the lattice $`L`$, and then verify that $`(L,B)`$ is an integral lattice, $`L^{\prime \prime }P^{}`$, $`L^{}`$ is positive-definite and the projection of $`L`$ onto $`L^{}`$ is one-to-one, so that Proposition 4.3 and Theorem 4.4 hold.
In this subsection we shall consider the case $`g=sl(n+1)`$. For a fixed positive integer $`k`$, $`L(k,0)`$ has $`n+1`$ simple currents $`L(k,k\lambda _i)`$ for $`i=0,\mathrm{},n`$. By Corollary 2.27 the equivalence classes of the $`(n+1)`$ simple currents form a cyclic group of order $`(n+1)`$ with $`[L(k,k\lambda _1)]`$ as a generator.
Recall that $`h^{(i)}𝐡`$ with $`\alpha _j(h^{(i)})=\delta _{i,1}`$ for $`i,j=1,\mathrm{},n`$. From \[H\] (Table 1 on page 69) and Lemma 2.28, we have
$`h^{(i)}={\displaystyle \frac{1}{n+1}}\left((n+1i)\alpha _1^{}+2(n+1i)\alpha _2^{}+\mathrm{}+(i1)(n+1i)\alpha _{i1}^{}\right)`$
$`+{\displaystyle \frac{1}{n+1}}\left(i(n+1i)\alpha _i^{}+i(ni)\alpha _{i+1}^{}+\mathrm{}+i\alpha _n^{}\right),`$ (4.25)
$`h^{(i)},h^{(i)}={\displaystyle \frac{i(n+1i)}{n+1}}.`$ (4.26)
Define $`𝐡^{}=C\alpha ^{}`$ to be a one-dimensional vector space equipped with a symmetric bilinear form $`,`$ such that
$`\alpha ^{},\alpha ^{}={\displaystyle \frac{k}{n+1}}.`$ (4.27)
Set
$`L=Z\alpha ,\text{ where }\alpha =h^{(1)}+\alpha ^{}.`$ (4.28)
Because
$`B(\alpha ,\alpha )=B(h^{(1)},h^{(1)})+B(\alpha ^{},\alpha ^{})={\displaystyle \frac{kn}{n+1}}+{\displaystyle \frac{k}{n+1}}=kZ,`$ (4.29)
$`(L,B)`$ is an integral lattice. Clearly, $`L^{\prime \prime }=Zh^{(1)}P^{}`$, $`L^{}`$ is positive-definite and the projection of $`L`$ onto $`L^{}`$ is one-to-one. By Proposition 4.3, we have a simple vertex operator (super)algebra $`V[L]`$.
###### Definition 4.7
We define $`A_k(sl(n+1))`$ to be the simple vertex operator (super)algebra $`V[L]`$ with $`V`$ and $`L`$ being defined in (4.24) and (4.28).
###### Remark 4.8
There are many ways to define $`\alpha ^{},\alpha ^{}`$ such that $`V[L]`$ is a vertex operator superalgebra. For examples, one may define $`𝐡^{}`$ with $`\alpha ^{},\alpha ^{}=1\frac{\overline{nk}}{n+1}`$, where $`\overline{nk}`$ is the least nonnegative residue of $`nk`$ modulo $`n+1`$. One may also define $`𝐡^{}`$ with $`\alpha ^{},\alpha ^{}=s+\frac{k}{n+1}`$, where $`s`$ is any nonnegative integer.
For $`\lambda P_k,\gamma C`$, set
$`W(\lambda ,\gamma )=L(k,\lambda )M_𝐡^{}(1,{\displaystyle \frac{\gamma }{k}}\alpha ^{}).`$ (4.30)
Since $`L=Z(h^{(1)}+\alpha ^{})`$, (4.14) amounts to
$`\lambda (h^{(1)})+{\displaystyle \frac{\gamma }{n+1}}Z,`$ (4.31)
which from (4.25) is equivalent to
$`n\lambda (\alpha _1^{})+(n1)\lambda (\alpha _2^{})+\mathrm{}+\lambda (\alpha _n^{})+\gamma (n+1)Z.`$ (4.32)
Note that (4.32) implies $`\gamma Z`$ because $`\lambda P_k`$. We also see that in general, $`\sigma _{W(\lambda ,\gamma )}`$ is of finite order if and only if $`\gamma Q`$. Then Propositions 3.13 and 3.24 immediately give:
###### Proposition 4.9
For $`\lambda P_k,\gamma Q`$, $`\sigma _{W(\lambda ,\gamma )}`$ is of finite order and $`W(\lambda ,r)[L]`$ is an irreducible $`\sigma _{W(\lambda ,\gamma )}`$-twisted $`A_k(sl(n+1))`$-module. In particular, if (4.32) holds, $`W(\lambda ,\gamma )[L]`$ is an irreducible $`V[L]`$-module. $`\mathrm{}`$
From Theorem 4.4 we also have:
###### Proposition 4.10
Let $`\sigma `$ be an automorphism of $`A_k(sl(n+1))`$ of finite order which fixes $`V=L(k,0)M_𝐡^{}(1,0)`$ point-wise. Then every $`\sigma `$-twisted weak $`A_k(sl(n+1))`$-module is a direct sum of irreducible (ordinary) $`\sigma `$-twisted $`A_k(sl(n+1))`$-modules $`W(\lambda ,\gamma )[L]`$ for some $`\lambda P_k,\gamma Q`$ with $`\sigma =\sigma _{W(\lambda ,\gamma )}`$. In particular, every weak $`A_k(sl(n+1))`$-module is a direct sum of $`W(\lambda ,\gamma )[L]`$ for some $`\lambda P_k,\gamma Z`$ that satisfy (4.32). $`\mathrm{}`$
Let us consider the case $`n=1`$. Then we can make our results more explicit. We have $`P_k=\{0,1,\mathrm{},k\}`$ and $`h^{(1)}=\frac{1}{2}\alpha _1^{}`$. From Corollary 2.27 we have
$`L(k,0)^{(2mh^{(1)})}L(k,0),L(k,0)^{((2m+1)h^{(1)})}L(k,k)`$ (4.33)
for $`mZ`$. Since
$`V^{(m\alpha )}=L(k,0)^{(mh^{(1)})}M_𝐡^{}(1,0)^{(m\alpha ^{})}=L(k,0)^{(mh^{(1)})}M_𝐡^{}(1,m\alpha ^{}),`$ (4.34)
it follows from Proposition 3.15 that $`W(i,\gamma )[L]W(i^{},\gamma ^{})[L]`$ if and only if there exists $`mZ`$ such that
$`L(k,i^{})L(k,i)^{(mh^{(1)})},{\displaystyle \frac{\gamma ^{}}{k}}={\displaystyle \frac{\gamma }{k}}+m.`$ (4.35)
Recall that $`[L(k,k)][L(k,i)]=[L(k,ki)]`$. If $`m`$ is even, we have $`i^{}=i`$ and $`\gamma ^{}\gamma 2kZ`$. If $`m`$ is odd, we have $`i^{}=ki`$ and $`\gamma ^{}=\gamma +mk`$. Then $`W(i,\gamma )[L]=W(i^{},\gamma ^{})[L]`$ if and only if either $`i^{}=i`$ and $`\gamma ^{}\gamma \mathrm{mod}\mathrm{\hspace{0.33em}2}k`$, or $`i^{}=ki`$ and $`\gamma ^{}\gamma +k\mathrm{mod}\mathrm{\hspace{0.33em}2}k`$. Then Propositions 4.9 and 4.10 give (cf. \[FM\], Proposition 3):
###### Corollary 4.11
Every weak $`A_k(sl(2))`$-module is completely reducible and
$`W(i,j)[L]\text{ for }0ik,\mathrm{\hspace{0.33em}0}jk1\text{ with }i+j2Z`$ (4.36)
form a complete set of representatives of equivalence classes of irreducible $`A_k(sl(2))`$-modules. $`\mathrm{}`$
Note that $`W(i,j)[L]`$ was denoted by $`R(i,j)`$ in \[FM\]. Using the fusion rules for $`L(k,0)`$ we have the following relations in the Verlinde algebra $`𝒜(V)`$:
$`[W(i_1,j_1)][W(i_2,j_2)]={\displaystyle \underset{i=\text{max}(i_1i_2,i_2i_1)}{\overset{\text{min}(i_1+i_2,2ki_1i_2)}{}}}[W(i,j_1+j_2)].`$ (4.37)
Then in the Verlinde algebra of $`A_k`$, we have
$`[W(i_1,j_1)[L]][W(i_2,j_2)[L]]={\displaystyle \underset{i=\text{max}(i_1i_2,i_2i_1)}{\overset{\text{min}(i_1+i_2,2ki_1i_2)}{}}}[W(i,j_1+j_2)[L]].`$ (4.38)
Note that when $`j_1+j_2k`$, we have $`W(i,j_1+j_2)[L]=W(ki,j_1+j_2k)[L]`$.
###### Remark 4.12
Set $`L^{}=Z\alpha ^{}`$. Then, as a $`V`$-module,
$`A_kL(k,0)V_{2L^{}}+L(k,k)V_{2L^{}+\alpha ^{}}.`$ (4.39)
Furthermore, using more general fusion rules we get
$`W(i,\gamma )[L]L(k,i)V_{2L^{}+\frac{\gamma }{k}\alpha ^{}}+L(k,ki)V_{2L^{}+\frac{\gamma +k}{k}\alpha ^{}}`$ (4.40)
for $`i=0,\mathrm{},k;\gamma Q`$. With this, one can easily write down the characters of $`W(i,j)[L]`$ in terms of the characters of $`L(k,j)`$ and the theta functions of $`2L^{}`$.
###### Remark 4.13
For $`g=sl(n+1)`$, from Corollary 2.27, we have
$`L(k,0)^{(mh^{(1)})}L(k,k\lambda _{\overline{m}})\text{ for }mZ.`$ (4.41)
Then
$`A_k(sl(n+1)){\displaystyle \underset{i=0}{\overset{n}{}}}L(k,k\lambda _i)V_{2(n+1)L^{}+i\alpha ^{}}`$ (4.42)
as a $`V`$-module. More general fusion rules are needed to express $`W(\lambda ,j)[L]`$ explicitly.
### 4.3 Generating property for the extended algebras $`A_k(sl(n+1))`$
First, we review some properties for a general vertex operator superalgebra $`U`$. Recall Borcherds’ commutator formula \[B\]:
$`[u_m,v_n]_\pm ={\displaystyle \underset{i0}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{m}{i}}\right)(u_iv)_{m+ni}`$ (4.43)
for $`u,vU`$ and $`m,nZ`$, where $`[,]_\pm `$ refers to the super commutator. Thus, the super commutator $`[Y(u,z_1),Y(v,z_2)]_\pm `$ is uniquely determined by $`u_iv`$ for $`i0`$. From this we have
$`(z_1z_2)^r[Y(u,z_1),Y(v,z_2)]_\pm =0`$ (4.44)
if $`r`$ is a nonnegative integer such that $`u_iv=0`$ for $`ir`$. For homogeneous vectors $`u,vU`$ and for $`mZ`$, we have (cf. \[FLM\])
$`\mathrm{wt}(u_mv)=\mathrm{wt}u+\mathrm{wt}vm1,`$ (4.45)
where $`\mathrm{wt}u`$ stands for the $`L(0)`$-weight of $`u`$.
Let $`U=_{n\frac{1}{2}Z}U_{(n)}`$ be such that $`U_{(0)}=C(=C\mathrm{𝟏})`$ and $`U_{(n)}=0`$ for $`n<0`$. Then
$`[u_m,v_n]_+=(u_0v)_{m+n1}=\delta _{m+n,1}u_0v`$ (4.46)
for $`u,vU_{(\frac{1}{2})},m,nZ`$, where $`u_0vU_{(0)}=C`$. That is, the component operators $`u_m`$ for $`uU_{(\frac{1}{2})},mZ`$ give rise to a Clliford algebra.
It is well known (\[B\], \[FLM\]) that the weight-one subspace $`U_{(1)}`$ is a Lie algebra with $`[u,v]=u_0v`$ and with a symmetric invariant bilinear form $`(u,v)=u_1vC`$. We have
$`[u_m,v_n]=(u_0v)_{m+n}+m\delta _{m+n,0}(u,v)`$ (4.47)
for $`u,vU_{(1)},m,nZ`$. Then operators $`u_m`$ for $`uU_{(1)},mZ`$ give rise to a natural representation of affine Lie algebra $`\widehat{U_{(1)}}`$.
Now we consider $`A_k(sl(n+1))`$, which is a vertex operator algebra when $`k`$ is even and which is a vertex operator superalgebra when $`k`$ is odd. It is easy to see that vertex operator (super)algebra $`A_k(sl(n+1))`$ is generated by $`V^{(\alpha )}`$ and $`V^{(\alpha )}`$. Denote by $`V_{low}^{(\beta )}`$ the lowest $`L(0)`$-weight subspace of $`V^{(\beta )}`$ for $`\beta L`$. Because $`V`$ as a vertex operator algebra is generated by $`g+C\alpha ^{}`$ and both $`V^{(\alpha )}`$ and $`V^{(\alpha )}`$ are irreducible $`V`$-modules, $`A_k(sl(n+1))`$ is furthermore generated by
$`S:=(g+C\alpha ^{})+V_{low}^{(\alpha )}+V_{low}^{(\alpha )}.`$ (4.48)
Since (Corollary 2.27)
$$V^{(\alpha )}=L(k,k\lambda _1)M_𝐡^{}(1,\alpha ^{}),V^{(\alpha )}=L(k,k\lambda _n)M_𝐡^{}(1,\alpha ^{}),$$
we have
$`V_{low}^{(\alpha )}=L(k\lambda _1)e^\alpha ^{},V_{low}^{(\alpha )}=L(k\lambda _n)e^\alpha ^{}.`$ (4.49)
From Remark 2.23 and (2.7), we find that the lowest $`L(0)`$-weights of $`V`$-modules $`V^{(\alpha )}`$ and $`V^{(\alpha )}`$ are $`\frac{1}{2}B(\alpha ,\alpha )=\frac{k}{2}`$. Now we are ready to prove our main result of this subsection.
###### Proposition 4.14
The algebra $`A_k(sl(n+1))`$ is generated by the subspace
$$(g+C\alpha ^{})+V_{low}^{(\alpha )}+V_{low}^{(\alpha )},$$
where $`V_{low}^{(\alpha )}=L(k\lambda _1)e^\alpha ^{}`$ and $`V_{low}^{(\alpha )}=L(k\lambda _n)e^\alpha ^{}`$ are of weight $`\frac{k}{2}`$. Furthermore, the following relations hold:
$`Y(u,z_1)Y(v,z_2)=(1)^kY(v,z_2)Y(u,z_1),`$ (4.50)
$`Y(u^{},z_1)Y(v^{},z_2)=(1)^kY(v^{},z_2)Y(u^{},z_1)`$ (4.51)
for $`u,vV_{low}^{(\alpha )},u^{},v^{}V_{low}^{(\alpha )}`$ and
$`u_sv^{}V_{(ks1)},`$ (4.52)
$`(z_1z_2)^k[Y(u,z_1),Y(v^{},z_2)]_\pm =0`$ (4.53)
for $`uV_{low}^{(\alpha )},v^{}V_{low}^{(\alpha )},sZ`$.
Proof. First we calculate the lowest $`L(0)`$-weight of $`V^{(m\alpha )}`$. From Theorem 2.26 and Corollary 2.27 we have
$$V^{(m\alpha )}=L(k,0)^{(mh^{(1)})}M_𝐡^{}(1,m\alpha ^{})=L(k,k\lambda _{\overline{m}})M_𝐡^{}(1,m\alpha ^{}),$$
where $`\overline{m}`$ is the least nonnegative residue of $`m`$ modulo $`n+1`$. From Remark 2.23, we see that the lowest weight of $`L(k,k\lambda _m)`$ is $`\frac{\overline{m}(n+1\overline{m})k}{2(n+1)}`$. Then the lowest weight of $`V^{(m\alpha )}`$ is
$$\frac{\overline{m}(n+1\overline{m})k}{2(n+1)}+\frac{m^2k}{2(n+1)}=\frac{\overline{m}k}{2}+\frac{(m^2\overline{m}^2)k}{2(n+1)},$$
which is at least $`k`$ if $`|m|2`$.
Let $`u,vV_{low}^{(\alpha )}`$. Thus $`\mathrm{wt}u=\mathrm{wt}v=\frac{k}{2}`$. Then for $`i0`$, $`u_ivV^{(2\alpha )}`$ and $`\mathrm{wt}(u_iv)=ki1<k`$. Since the lowest weight of $`V^{(2\beta )}`$ is at least $`k`$, we obtain
$`u_iv=0\text{ for }i0.`$ (4.54)
Then (4.50) follows immediately from (4.43). Similarly, (4.51) holds.
(4.52) directly follows from the definition of the vertex operator map and the weight formula (4.45). Since $`u_sv^{}V^{(0)}=V`$ for $`sZ`$ and the lowest weight of $`V`$ is zero, we have
$`u_iv^{}=0\text{ for }ik.`$ (4.55)
Then (4.53) follows immediately from (4.44). $`\mathrm{}`$
###### Remark 4.15
In the case $`k=1`$, $`L(\lambda _1)`$ is the vector representation of $`sl(n+1)`$ on $`C^{n+1}`$. In this case, the algebra $`A_1(sl(n+1))`$ is generated by $`L(\lambda _1)e^\alpha ^{}+L(\lambda _1)^{}e^\alpha ^{}`$, which generates a Clliford algebra. The algebra $`A_1(sl(n+1))`$ is exactly the spinor representation of $`D_{n+1}`$ \[FFR\], which is isomorphic to $`L(1,0)+L(1,\lambda _1)`$ as a $`\widehat{D}_{n+1}`$-module.
###### Remark 4.16
When $`k=2`$,
$$L(2\lambda _1)^{}e^\alpha ^{}+(g+C\alpha ^{})+L(2\lambda _1)e^\alpha ^{}$$
is exactly the weight-one subspace of $`A_2(sl(n+1))`$ and it has a natural Lie algebra structure with the obvious $`Z`$-grading. Using the fact that $`L(2\lambda _1)^{}e^\alpha ^{}`$ and $`L(2\lambda _1)e^\alpha ^{}`$ are non-isomorphic irreducible $`(g+C\alpha ^{})`$-modules one easily shows that this Lie algebra is simple and of rank $`n+1`$. Consider the standard Dynkin diagram embedding of $`sl(n+1)`$ into $`C_{n+1}`$. Then we see
$$C_{n+1}=sl(n+1)+L(2\lambda _1)+L(2\lambda _n).$$
Thus the weight one subspace of $`A_2(sl(n+1))`$ as a Lie algebra is isomorphic to $`C_{n+1}`$. (For $`n=1`$, this was pointed out in \[FM\].) Then $`A_2(sl(n+1))`$ is a vertex operator algebra associated to the affine Lie algebra $`\widehat{C}_{n+1}`$.
###### Remark 4.17
For $`k3`$, since $`\mathrm{wt}(u_0v^{})=k12`$, $`[Y(u,z_1),Y(v^{},z_2)]_\pm `$ involves nonlinear normal ordered products of vertex operators (or fields) $`Y(a,z)`$ for $`ag+C\alpha ^{}`$. This type of algebras are commonly referred by physicists as nonlinear $`W`$-algebras.
###### Remark 4.18
The following consideration was motivated by \[GH1-2\] and \[Gun\]. In the construction of $`A_k`$, let us define $`𝐡^{}=C\alpha ^{}`$ with $`\alpha ^{},\alpha ^{}=1+\frac{k}{n+1}`$ and keep the rest unchanged. Then $`B(\alpha ,\alpha )=1+k`$. With $`(L,B)`$ being an integral lattice, $`V[L]`$ is a vertex operator (super)algebra (cf. Remark 4.8). Furthermore, $`V[L]`$ is generated by the subspace
$$V_{low}^{(\alpha )}+(g+C\alpha ^{})+V_{low}^{(\alpha )}=L(k\lambda _1)^{}e^\alpha ^{}+(g+C\alpha ^{})+L(k\lambda _1)e^\alpha ^{},$$
where $`L(k\lambda _1)e^\alpha ^{}`$ and $`L(k\lambda _1)^{}e^\alpha ^{}`$ are of weight $`\frac{k+1}{2}`$. In particular, when $`k=2`$, $`V[L]`$ is a vertex operator superalgebra and $`L(2\lambda _1)e^\alpha ^{}`$ and $`L(2\lambda _1)^{}e^\alpha ^{}`$ are of weight $`\frac{3}{2}`$. In view of this and Remark 4.16, we may view $`V[L]`$ with $`k=2`$ as a superization of the vertex operator algebra $`V[L]`$ with $`k=2`$ defined in Remark 4.16. In \[Gun\], an $`N=2`$ vertex operator superalgebra was constructed from a simple Lie algebra $`\gamma `$ equipped with a $`Z`$-grading such that $`g_m=0`$ for $`|m|>1`$. From Remark 4.16, symplectic Lie algebra $`C=C_{n+1}`$ is naturally $`Z`$-graded with only three homogeneous subspaces being nonzero. A further study on the connection between $`V[L]`$ and the $`N=2`$ vertex operator superalgebra constructed in \[Gun\] will be conducted in a future paper.
### 4.4 Extended algebras $`A_k`$ of type $`D_n`$
We consider $`g`$ of $`D_n`$ type for $`n3`$. From Corollary 2.27, the equivalence classes of simple currents $`L(k,k\lambda _1)`$, $`L(k,k\lambda _{n1})`$, $`L(k,k\lambda _n)`$ and $`L(k,0)`$ form a group which is cyclic for an odd $`n`$ and which is isomorphic to $`Z/2Z\times Z/2Z`$ for an even $`n`$. We shall define the extended algebra separately for the two cases.
From \[H\] and Lemma 2.28 we have
$`h^{(n1)}={\displaystyle \frac{1}{2}}(\alpha _1^{}+2\alpha _2^{}+\mathrm{}+(n2)\alpha _{n2}^{}+{\displaystyle \frac{1}{2}}n\alpha _{n1}^{}+{\displaystyle \frac{1}{2}}(n2)\alpha _n^{}),`$ (4.56)
$`h^{(n)}={\displaystyle \frac{1}{2}}(\alpha _1^{}+2\alpha _2^{}+\mathrm{}+(n2)\alpha _{n2}^{}+{\displaystyle \frac{1}{2}}(n2)\alpha _{n1}^{}+{\displaystyle \frac{1}{2}}n\alpha _n^{})`$ (4.57)
and
$`h^{(n1)},h^{(n1)}=h^{(n)},h^{(n)}={\displaystyle \frac{n}{4}}.`$ (4.58)
Using the relation $`h^{(n1)}=h^{(n)}+\frac{1}{2}\alpha _{n1}^{}\frac{1}{2}\alpha _n^{}`$ we get
$`h^{(n1)},h^{(n)}={\displaystyle \frac{n2}{4}}.`$ (4.59)
Case I, $`n`$ is odd.
Define $`𝐡^{}=C\alpha ^{}`$ with $`\alpha ^{},\alpha ^{}=\frac{3nk}{4}`$. Set
$`L=Z\alpha ,\text{ where }\alpha =h^{(n)}+\alpha ^{}.`$ (4.60)
Then
$`L^{}=Z\alpha ^{},L^{\prime \prime }=Zh^{(n)}.`$ (4.61)
Since
$`B(\alpha ,\alpha )=kh^{(n)},h^{(n)}+\alpha ^{},\alpha ^{}=nk,`$ (4.62)
$`(L,B)`$ is a positive definite integral lattice. By Proposition 4.3 (with the other assumptions being obvious) $`V[L]`$ is a simple vertex operator (super)algebra. We define $`A_k(g)`$ to be $`V[L]`$. Then we have the following results with the same proof as that of Propositions 4.9 and 4.10:
###### Proposition 4.19
For $`g`$ of type $`D_n`$ with an odd $`n`$, the extended algebra $`A_k(g)`$ is regular. Furthermore, for $`\lambda P_k,jQ`$, set
$$W(\lambda ,j)=L(k,\lambda )M_𝐡^{}(1,\frac{j}{3nk}\alpha ^{}).$$
Then any irreducible $`A_k(g)`$-module is isomorphic to $`W(\lambda ,j)[L]`$ for some $`\lambda P_k,jZ`$ with
$`2\lambda (\alpha _1^{})+4\lambda (\alpha _2^{})+\mathrm{}+2(n2)\lambda (\alpha _{n2}^{})+(n2)\lambda (\alpha _{n1}^{})+n\lambda (\alpha _n^{})+j4Z.`$ (4.63)
Furthermore, $`A_k(g)`$ is generated by
$$(g+C\alpha ^{})+V_{low}^{(\alpha )}+V_{low}^{(\alpha )},$$
where $`V_{low}^{(\alpha )}=L(k\lambda _n)e^\alpha ^{}`$ and $`V_{low}^{(\alpha )}=L(k\lambda _{n1})e^\alpha ^{}`$ are of weight $`\frac{nk}{2}`$, and the relations (4.50)-(4.53) with $`k`$ being replaced by $`nk`$ hold.$`\mathrm{}`$
Case II: $`n`$ is even.
Define $`𝐡^{}=C\alpha _1^{}+C\alpha _2^{}`$ to be a two-dimensional vector space with a symmetric bilinear form $`,`$ such that
$`\alpha _1^{},\alpha _1^{}=\alpha _1^{},\alpha _1^{}={\displaystyle \frac{3}{4}}nk,`$ (4.64)
$`\alpha _1^{},\alpha _2^{}={\displaystyle \frac{1}{4}}k(n2).`$ (4.65)
Define
$`L=Z\alpha _1+Z\alpha _2,`$ (4.66)
where
$`\alpha _1=h^{(n1)}+\alpha _1^{},\alpha _2=h^{(n)}+\alpha _2^{}.`$ (4.67)
Then
$`L^{}=Z\alpha _1^{}+Z\alpha _2^{},L^{\prime \prime }=Zh^{(n1)}+Zh^{(n)}.`$ (4.68)
We have
$`B(\alpha _1,\alpha _1)=B(\alpha _2,\alpha _2)=kn,`$ (4.69)
$`B(\alpha _1,\alpha _2)={\displaystyle \frac{1}{4}}k(n2)+{\displaystyle \frac{1}{4}}k(n2)={\displaystyle \frac{1}{2}}k(n2).`$ (4.70)
Since $`n`$ is even, $`(L,B)`$ is a positive-definite even lattice. Clearly, $`L^{\prime \prime }=Zh^{(n1)}+Zh^{(n)}P^{}`$, $`L^{}`$ is positive-definite, and the projection of $`L`$ onto $`L^{}`$ is one-to-one. By Proposition 4.3, $`V[L]`$ is a simple vertex operator algebra. Now we define $`A_k(g)=V[L]`$, as a simple vertex operator algebra. We just mention that this is a regular vertex operator algebra and a set of generators and relations can be worked out similarly but with some extra work.
###### Remark 4.20
Note that $`L(k,k\lambda _1)`$ is a simple current of order 2 and we have $`h^{(1)},h^{(1)}=1`$. Let $`V=L(k,0)`$ and $`L=Zh^{(1)}`$. Then in view of Corollary 3.21, $`L(k,0)+L(k,k\lambda _1)`$ has a natural simple vertex operator superalgebra structure (cf. Remark 4.15).
### 4.5 Extended vertex operator (super)algebras $`A_k(E_6)`$
Let $`g`$ be of type $`E_6`$. From Section 2.2, for any positive integer $`k`$, $`L(k,k\lambda _1)`$ and $`L(k,k\lambda _5)`$ are (the only) nontrivial simple currents for $`L(k,0)`$. From \[H\] and Lemma 2.28, we have
$`h^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left(4\alpha _1^{}+3\alpha _6^{}+5\alpha _2^{}+6\alpha _3^{}+4\alpha _4^{}+2\alpha _5^{}\right),`$ (4.71)
$`h^{(5)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left(2\alpha _1^{}+3\alpha _6^{}+4\alpha _2^{}+6\alpha _3^{}+5\alpha _4^{}+4\alpha _5^{}\right)`$ (4.72)
and
$`h^{(1)},h^{(1)}=h^{(5)},h^{(5)}={\displaystyle \frac{4}{3}}.`$ (4.73)
Define $`𝐡^{}=C\alpha ^{}`$ to be a one-dimensional vector space equipped with the bilinear form $`,`$ such that
$`\alpha ^{},\alpha ^{}={\displaystyle \frac{2k}{3}}.`$ (4.74)
Set
$`L=Z\alpha ,\text{ where }\alpha =h^{(1)}+\alpha ^{}.`$ (4.75)
Then
$`L^{\prime \prime }=Zh^{(1)}P^{},L^{}=Z\alpha ^{}.`$ (4.76)
Since
$`B(\alpha ,\alpha )={\displaystyle \frac{4k}{3}}+{\displaystyle \frac{2k}{3}}=2k,`$ (4.77)
$`(L,B)`$ is a positive-definite even lattice. By Proposition 4.3 (with the other assumptions being obvious), we have a simple vertex operator algebra $`V[L]`$. We define $`A_k(g)`$ to be the vertex operator algebra $`V[L]`$. For $`\lambda P_k,jQ`$, set
$`W(\lambda ,j)=L(k,\lambda )M_𝐡^{}(1,{\displaystyle \frac{1}{2k}}\alpha ^{}),`$ (4.78)
an irreducible $`V`$-module. Then we have:
###### Proposition 4.21
For $`g`$ of type $`E_6`$, the extended algebra $`A_k(g)`$ is regular and any irreducible module is isomorphic to $`W(\lambda ,j)[L]`$ for some $`\lambda P_k,jZ`$ with
$`4\lambda (\alpha _1^{})+3\lambda (\alpha _6^{})+5\lambda (\alpha _2^{})+6\lambda (\alpha _3^{})+4\lambda (\alpha _4^{})+2\lambda (\alpha _5^{})+j3Z.\mathrm{}`$ (4.79)
Similar to the case $`g=sl(n+1)`$, $`A_k(E_6)`$ as a vertex operator algebra is generated by
$$(g+C\alpha ^{})+V_{low}^{(\alpha )}+V_{low}^{(\alpha )}.$$
We have
$`V_{low}^{(\alpha )}=L(k,k\lambda _1)_{low}e^\alpha ^{}=L(k\lambda _1)e^\alpha ^{},`$ (4.80)
$`V_{low}^{(\alpha )}=L(k,k\lambda _1)_{low}^{}e^\alpha ^{}=L(k\lambda _1)^{}e^\alpha ^{}.`$ (4.81)
Since the $`L(0)`$-weights of $`L(k,k\lambda _1)_{low}`$ and $`L(k,k\lambda _1)_{low}^{}`$ are
$$\frac{1}{2}B(h^{(1)},h^{(1)})=\frac{k}{2}h^{(1)},h^{(1)}=\frac{2k}{3},$$
the lowest weights of $`V^{(\alpha )}`$ and $`V^{(\alpha )}`$ are $`\frac{2k}{3}+\frac{k}{3}=k`$.
For $`mZ`$, the lowest $`L(0)`$-weight of $`M_𝐡^{}(1,m\alpha ^{})`$ is $`\frac{1}{2}m\alpha ^{},m\alpha ^{}=\frac{km^2}{3}`$. Then the lowest $`L(0)`$-weight of $`V^{(m\alpha )}`$ is at least $`\frac{km^2}{3}`$. If $`|m|3`$, the lowest $`L(0)`$-weight of $`V^{(m\alpha )}`$ is at least $`3k`$. The lowest $`L(0)`$-weight of $`V^{(2\alpha )}`$ is the sum of the lowest $`L(0)`$-weight of $`L(k,0)^{(2h^{(1)})}`$ and $`\frac{4k}{3}`$. We know that $`L(k,0)^{(2h^{(1)})}L(k,k\lambda _5)`$ whose lowest $`L(0)`$-weight is $`\frac{4k}{3}`$. Then the lowest $`L(0)`$-weight of $`V^{(2\alpha )}`$ is $`\frac{8k}{3}`$, which is greater than $`2k`$. With this information, using the same proof of Proposition 4.14 we immediately have:
###### Proposition 4.22
The extended vertex operator algebra $`A_k(E_6)`$ is generated by
$`(g+C\alpha ^{})+V_{low}^{(\alpha )}+V_{low}^{(\alpha )},`$ (4.82)
where $`V_{low}^{(\alpha )}=L(k\lambda _1)e^\alpha ^{}`$ and $`V_{low}^{(\alpha )}=L(k\lambda _1)^{}e^\alpha ^{}`$ are of weight $`k`$. Furthermore, the relations (4.50)-(4.53) with $`k`$ being replaced by $`2k`$ hold.$`\mathrm{}`$
###### Remark 4.23
When $`k=1`$, $`V_{low}^{(\alpha )}+(g+C\alpha ^{})+V_{low}^{(\alpha )}`$ is exactly the weight-one subspace of $`A_1(g)`$, which is a natural Lie algebra with
$`[L(\lambda _1)e^\alpha ^{},L(\lambda _1)e^\alpha ^{}]=0,[L(\lambda _1)^{}e^\alpha ^{},L(\lambda _1)^{}e^\alpha ^{}]=0,`$ (4.83)
$`[L(\lambda _1)e^\alpha ^{},L(\lambda _1)^{}e^\alpha ^{}]g+C\alpha ^{}.`$ (4.84)
These relations give rise to a $`Z`$-grading for the Lie algebra. One can easily show that this Lie algebra is simple and of rank $`7`$. Using the standard Dynkin diagram embedding of $`E_6`$ into $`E_7`$ we can show that it is really $`E_7`$.
### 4.6 Extended vertex operator (super)algebras $`A_k(E_7)`$
Let $`g`$ be of type $`E_7`$. From Section 2.2, for any positive integer $`k`$, $`L(k,k\lambda _6)`$ is a (and the only nontrivial) simple current for $`L(k,0)`$. Using \[H\] (Table 1 on page 69) and Lemma 2.28 we have
$`h^{(6)}={\displaystyle \frac{1}{2}}\left(2\alpha _1^{}+3\alpha _7^{}+4\alpha _2^{}+6\alpha _3^{}+5\alpha _4^{}+4\alpha _5^{}+3\alpha _6^{}\right),h^{(6)},h^{(6)}={\displaystyle \frac{3}{2}}.`$ (4.85)
Define $`𝐡^{}=C\alpha ^{}`$ to be a one-dimensional vector space equipped with a bilinear form $`,`$ such that
$`\alpha ^{},\alpha ^{}={\displaystyle \frac{k}{2}}.`$ (4.86)
Set
$`L=Z\alpha ,\text{ where }\alpha =h^{(6)}+\alpha ^{}.`$ (4.87)
Then $`L^{}=Z\alpha ^{}`$ and $`L^{\prime \prime }=Zh^{(6)}`$. Since $`B(\alpha ,\alpha )=\frac{3k}{2}+\frac{k}{2}=2k`$, $`(L,B)`$ is a positive-definite even lattice. By Proposition 4.3 (with the other assumptions being obvious), $`V[L]`$ is a simple vertex operator algebra. We define $`A_k(E_7)`$ to be the simple vertex operator algebra $`V[L]`$. For $`\lambda P_k,jQ`$, we set
$`W(\lambda ,j)=L(k,\lambda )M_𝐡^{}(1,{\displaystyle \frac{j}{k}}\alpha ^{}).`$ (4.88)
In view of Theorem 4.4 we immediately have:
###### Proposition 4.24
For $`g`$ of type $`E_7`$, the extended algebra $`A_k(g)`$ is regular and any irreducible module is isomorphic to $`W(\lambda ,j)[L]`$ for $`\lambda P_k,jZ`$ with
$`2\lambda (\alpha _1^{})+3\lambda (\alpha _7^{})+4\lambda (\alpha _2^{})+6\lambda (\alpha _3^{})+5\lambda (\alpha _4^{})+4\lambda (\alpha _5^{})+3\lambda (\alpha _6^{})+j2Z.\mathrm{}`$ (4.89)
The lowest weights of $`V^{(\alpha )}`$ and $`V^{(\alpha )}`$ are $`\frac{1}{2}B(\alpha ,\alpha )=k`$. Since $`[L(k,0)^{(2h^{(6)})}]=[L(k,0)]`$, the lowest weights of $`V^{(2\alpha )}`$ and $`V^{(2\alpha )}`$ are $`\frac{1}{2}B(2\alpha ^{},2\alpha ^{})=k`$ also. For $`|m|3`$, the lowest weight of $`V^{(m\alpha )}`$ is at least $`\frac{1}{2}B(m\alpha ^{},m\alpha ^{})=\frac{m^2k}{4}`$, which is greater than $`2k`$. With this information we immediately have:
###### Proposition 4.25
The vertex operator algebra $`A_k(E_7)`$ is generated by
$$V_{low}^{(2\alpha )}+V_{low}^{(\alpha )}+(g+C\alpha ^{})+V_{low}^{(\alpha )}+V_{low}^{(2\alpha )},$$
where
$`V_{low}^{(\alpha )}=L(k\lambda _6)e^\alpha ^{},V_{low}^{(\alpha )}=L(k\lambda _6)e^\alpha ^{},`$ (4.90)
$`V_{low}^{(2\alpha )}=Ce^{2\alpha ^{}},V_{low}^{(2\alpha )}=Ce^{2\alpha ^{}}`$ (4.91)
are of weight $`k`$. $`\mathrm{}`$
###### Remark 4.26
When $`k=1`$,
$$Ce^{2\alpha ^{}}+L(\lambda _6)e^\alpha ^{}+(g+C\alpha ^{})+L(\lambda _6)e^\alpha ^{}+Ce^{2\alpha ^{}}$$
is exactly the weight-one subspace of $`A_1(g)`$. It is a $`Z`$-graded Lie algebra with the obvious grading. Similarly, we can show that it is $`E_8`$.
### 4.7 Extended algebras of types $`B_n`$ and $`C_n`$
For $`g`$ of type $`B_n`$, $`L(k,k\lambda _1)`$ is the only nontrivial simple current and for $`g`$ of type $`C_n`$, $`L(k,k\lambda _n)`$ is the only nontrivial simple current. For $`B_n`$, from \[H\] and Lemma 2.28 we have
$`h^{(1)}=\alpha _1^{}+\mathrm{}+a_{n1}^{}+{\displaystyle \frac{1}{2}}\alpha _n^{},h^{(1)},h^{(1)}=1`$ (4.92)
and for $`C_n`$ we have
$`h^{(n)}={\displaystyle \frac{1}{2}}(\alpha _1^{}+2\alpha _2^{}+\mathrm{}+n\alpha _n^{}),h^{(n)},h^{(n)}={\displaystyle \frac{n}{2}}.`$ (4.93)
###### Remark 4.27
We here correct an error of \[DLM2\] (Examples 5.12 and 5.13) where it was stated that $`L(k,k\lambda _n)`$ was the nontrivial simple current for $`g`$ of type $`B_n`$ and that $`L(k,k\lambda _1)`$ was the nontrivial simple current for $`g`$ of type $`C_n`$. For $`g`$ of type $`B_n`$, it follows from \[DLM2\] or Proposition 3.20 with $`V=L(k,0)`$ and $`L=Zh^{(1)}`$ that for any positive integral level $`k`$, $`L(k,0)+L(k,k\lambda _1)`$ has a natural simple vertex operator (super)algebra structure. However, for $`C_n`$, $`L(k,0)+L(k,k\lambda _n)`$ is a vertex operator (super)algebra only for a positive integral level $`k`$ with $`nk`$ being even.
For $`g`$ of type $`C_n`$, we define $`𝐡^{}=C\alpha ^{}`$ with $`\alpha ^{},\alpha ^{}=\frac{nk}{2}`$. Set
$`L=Z\alpha ,\text{ where }\alpha =h^{(n)}+\alpha ^{}.`$ (4.94)
Then $`L^{\prime \prime }=Zh^{(n)}`$ and $`L^{}=Z\alpha ^{}`$. Furthermore,
$`B(\alpha ,\alpha )=kh^{(n)},h^{(n)}+\alpha ^{},\alpha ^{}=nk.`$ (4.95)
Then $`(L,B)`$ is a positive definite integral lattice. Hence $`V[L]`$ is a simple vertex operator (super)algebra. Furthermore, we have:
###### Proposition 4.28
For $`g`$ of type $`C_n`$, $`A_k(g)`$ is regular and any irreducible $`A_k(g)`$-module is isomorphic to $`W(\lambda ,j)[L]`$ for $`\lambda P_k,jZ`$ with
$`\lambda (\alpha _1^{})+2\lambda (\alpha _2^{})+\mathrm{}+n\lambda (\alpha _n^{})+j2Z,`$ (4.96)
where
$`W(\lambda ,j)=L(k,\lambda )M_𝐡^{}(1,{\displaystyle \frac{j}{nk}}\alpha ^{}).`$ (4.97)
Furthermore, $`A_k(g)`$ is generated by
$$(g+C\alpha ^{})+V_{low}^{(\alpha )}+V_{low}^{(\alpha )}$$
and the the relations (4.50)-(4.53) with $`k`$ being replaced by $`nk`$ hold, where $`V_{low}^{(\alpha )}`$ and $`V_{low}^{(\alpha )}`$ are of weight $`\frac{nk}{2}`$. $`\mathrm{}`$
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# Quenched Spin Tunneling and Diabolical Points in Magnetic Molecules: II. Asymmetric Configurations
## I Introduction
### A The story so far
In a previous paper, hereafter cited as I, we studied the tunneling of a spin governed by the Hamiltonian
$$=k_2J_z^2+(k_1k_2)J_x^2g\mu _B𝐉𝐇,$$
(1)
where $`𝐉`$ is a dimensionless spin operator, and $`k_1>k_2>0`$. This Hamiltonian is the simplest descriptor of the magnetic properties of the molecule \[(tacn)<sub>6</sub>Fe<sub>8</sub>O<sub>2</sub>(OH)<sub>12</sub>\]<sup>8+</sup> (or just Fe<sub>8</sub> for short), with $`J=10`$, $`k_10.33`$ K, and $`k_20.22`$ K . Interest in this molecule arises because of its rich low temperature magnetic behavior, which include hysteresis at the level of one molecule , and more recently, the discovery of an entire lattice of diabolical points in its magnetic spectrum, a subset of which was predicted to exist earlier . It is the latter property that we wish to continue investigating in this paper.
In paper I, attention was confined to the case where the external magnetic field $`𝐇`$ is along the hard axis $`𝐱`$. In this case, the classical energy, which may be viewed as the expectation value $``$ of the Hamiltonian in a spin-coherent state, is symmetric about the $`xy`$ plane, and the problem is analogous to that of a massive particle in one dimension in a reflection symmetric double well potential. It is then natural to consider the tunneling between the symmetrically related states localized in the left- and right-hand wells. In the magnetic case, the analogous states are those with predominatly positive and negative values of $`J_z`$, at least as long as $`H_x`$ is not so large as to bring the classical minima very close to the $`x`$ axis. The surprise is that the tunnel splitting between the ground states oscillates as a function of $`H_x`$, vanishing exactly at a series of points. In fact, the splitting between higher pairs of levels also vanishes at just these points, as noted in I and earlier studies .
In addition to observing the quenching of ground state tunneling when $`𝐇𝐱`$, however, Wernsdorfer and Sessoli also performed experiments with $`H_z0`$. The reflection symmetry of the classical energy is now lost, but if the value of $`H_z`$ is chosen properly, it is possible to bring an excited state in the positive $`J_z`$ well into approximate degeneracy with the ground state of the negative $`J_z`$ well. The new discovery by them is that if $`H_x`$ is now varied, the tunnel splitting between the degenerate or quasi-degenerate levels again oscillates. It is theoretically understood that if both $`H_z`$ and $`H_x`$ are properly tuned, the splitting vanishes exactly in this case too . Experimentally, of course, one can never see a perfect zero in the splitting, and Wernsdorfer and Sessoli only see a minimum in the Landau-Zener-Stückelberg transition rate between the levels in question. The minima are so deep, however, that there is little doubt that the underlying tunneling matrix element is quenched.
When $`𝐇𝐱`$, the Hamiltonian (1) is invariant under a 180 rotation about $`𝐱`$, and the quenchings can be understood from the viewpoint of the von Neumann-Wigner theorem as allowed crossings of energy levels with different parities under this rotation. A similar argument can be made when $`𝐇𝐳`$. When $`𝐇`$ has both $`x`$ and $`z`$ components, however, the Hamiltonian has no obvious symmetry, and the above theorem states that an intersection of two energy levels is infinitely unlikely as a single parameter in the Hamiltonian is varied. For a real symmetric Hamiltonian, it is known that one must vary two parameters to obtain an intersection. Since we can choose the matrices of both $`J_z`$ and $`J_x`$ to be real in the $`J_z`$ basis, these two parameters can be taken as $`H_x`$ and $`H_z`$. The isolated points in the $`H_x`$-$`H_z`$ plane where any two levels intersect are precisely what Herzberg and Longuet-Higgins call conical intersections and what Berry and Wilkinson call diabolical points. The latter terminology originates in the resemblance of the energy surface—a double cone with a common vertex —to an Italian toy called the diavolo.
### B Content and plan of this paper; the perfect lattice hypothesis
In this paper we shall allow $`𝐇`$ to lie in the $`xz`$ plane, with a view to studying the tunneling in the asymmetrical well, and locating the diabolical points. As in paper I, our analysis is based on the discrete phase integral (DPI), or Wentzel-Kramers-Brillouin method . This method is semiclassical in character, with $`1/J`$ playing the same role as $`\mathrm{}`$ in the continuum phase integral method. To introduce this method, let us write the Schrödinger equation $`|\psi =E|\psi `$ in the $`J_z`$ eigenbasis $`\{|m\}`$. With $`J_z|m=m|m`$, $`m|\psi =C_m`$, $`m||m=w_m`$, and $`m||m^{}=t_{m,m^{}}`$ ($`mm^{}`$), we have
$$\stackrel{}{}_{n=m2}^{m+2}t_{m,n}C_n+w_mC_m=EC_m,$$
(2)
where the prime on the sum indicates omission of the $`n=m`$ term. A vivid picture of the approximation can be obtained if we think of Eq. (2) as the tight-binding model for an electron in a one-dimensional lattice with sites labelled by $`m`$, and on-site ($`w_m`$) and hopping ($`t_{m,m\pm 1}`$, $`t_{m,m\pm 2}`$) energies. If $`J1`$, these matrix elements vary slowly with $`m`$, on a length scale of order $`J`$ in fact. We may then use the approximation of semiclassical dynamics by working entirely in terms of wavepackets whose spatial extent is much less than the length scale over which the properties of the lattice vary, i.e., $`J`$, and whose spread in Bloch vectors is much less than the width of the Brillouin zone, i.e., $`2\pi `$. These ideas have close counterparts in the continuum quasiclassical method, and the DPI method is nothing but the discrete analog.
When $`𝐇𝐱`$, the problem can also be approached using instantons—indeed the oscillations in the splitting were discovered in this way. When $`𝐇`$ has other components besides $`H_x`$, however, the DPI method is, to our knowledge, the only successful one to date. Villain and Fort’s approach is also an approximate version of this method that makes artful use of some special features of the Fe<sub>8</sub> problem and works for small values of the field. Our analysis is more prosaic, and almost self-evident once one has understood how to deal with the new feature in Eq. (2)—the presence of second neighbour hopping. This gives rise to novel turning points with no continuum analogues. Connection formulas for these turning points are given in Ref. . We have quoted the results of our analysis before (b), but the details are only being presented here.
Our main result for the specific Hamiltonian (1) is for the locations of the diabolical points. We find that the $`\mathrm{}^{}`$th level in the negative $`J_z`$ well (where $`\mathrm{}^{}=0`$ denotes the lowest level) and the $`\mathrm{}^{\prime \prime }`$th level in the positive $`J_z`$ well are degenerate when (see Fig. 1)
$`{\displaystyle \frac{H_z(\mathrm{}^{},\mathrm{}^{\prime \prime })}{H_c}}`$ $`=`$ $`{\displaystyle \frac{\sqrt{\lambda }(\mathrm{}^{\prime \prime }\mathrm{}^{})}{2J}}`$ (3)
$`{\displaystyle \frac{H_x(\mathrm{}^{},\mathrm{}^{\prime \prime })}{H_c}}`$ $`=`$ $`{\displaystyle \frac{\sqrt{1\lambda }}{J}}\left[Jn\frac{1}{2}(\mathrm{}^{}+\mathrm{}^{\prime \prime }+1)\right],`$ (4)
with $`n=0,1,\mathrm{},2J(\mathrm{}^{}+\mathrm{}^{\prime \prime }+1)`$. Here, $`\lambda =k_2/k_1`$, and $`H_c=2k_1J/g\mu _B`$.
It should be stressed that Eqs. (3) and (4) only represent the first terms of a series in $`1/J`$, and that our DPI calculations give no reason to believe that the higher order terms are absent. Yet a large amount of empirical evidence suggests just this, i.e., that the formulas are exact as written! That the diabolical points lie on a perfect centered rectangular lattice, and that many pairs of levels are simultaneously degenerate, we refer to as the perfect lattice hypothesis. It has previously been made for the ground state on the $`H_z=0`$ line by Villain and Fort , and extended to include simultaneous degeneracy of the higher states by us (b). This was shown by perturbation theory in $`\lambda `$ to $`O(\lambda ^3)`$ for all $`J`$, and analytically for $`J2`$. We have since found analytically that for $`J=3/2`$, the additional diabolical point at $`(H_x,H_z)=(\sqrt{1\lambda },\sqrt{\lambda })H_c/3`$ is exact. Of course, Kramers theorem guarantees double degeneracy of all energy eigenvalues at the point $`H_x=H_z=0`$ for all half-integral $`J`$. We have also checked this hypothesis numerically for a variety of values of $`\lambda `$, and $`J`$ up to 10. We present our results for $`J=5`$ in Table I. The deviations in the locations of the diabolical points from the perfect lattice hypothesis predictions are never more than $`10^{10}`$, certainly much less than order $`1/J`$. In fact, the values listed are below the numerical tolerance that we prescribed.
Further support for the perfect lattice hypothesis comes from the fact that it is consistent with the following duality property of the Hamiltonian (1). If we set $`k_1=1`$, $``$ may be written as $`(\lambda ,H_x,H_z)`$, showing its dependence on the three parameters, $`\lambda `$, $`H_x`$ and $`H_z`$. By a $`90^{}`$ rotation about the axis $`(\widehat{𝐱}+\widehat{𝐳})/\sqrt{\mathrm{𝟐}}`$, we obtain the transformation
$$(\lambda ,H_x,H_z)(1\lambda ,H_z,H_x).$$
(5)
In particular, the spectra of the two Hamiltonians are so related, and ranking the levels is order of increasing energy, we see that if the levels with ordinal numbers $`k`$ and $`k+1`$ are degenerate when $`H_x=f_x(\lambda )`$ and $`H_y=f_y(\lambda )`$, then level numbers $`2J+2k`$ and $`2J+1k`$ are degenerate when $`H_x=f_y(1\lambda )`$ and $`H_y=f_x(1\lambda )`$. These conditions do not constrain the functions $`f_x`$ and $`f_y`$ in any real way, however, so the discovery of formulas (3) and (4) must presently be put down to serendipity.
The plan of our paper is as follows. In Sec. II, we analyze the asymmetric double well problem in completely general terms. We briefly review the DPI method and the asymmetric double-well problem in one dimensional quantum mechanics. In contrast to the situation that prevails in that case, we must deal with four independent DPI wavefunctions at every point. Consequently the matching problem is harder and more subtle. Its study forms the bulk (subsections C to G) of the analysis. The main result of this analysis is a formula for the diabolo \[Eq. (73)\], and a two level system Hamiltonian \[Eq. (74)\]. The tunneling amplitude appearing in these formulas is found in a general form \[Eq. (70)\] involving a small number of action integrals running between various turning points. Readers who are not interested in the details of the analysis should skip to subsection H and read on from Eq. (73). They may also find the last paragraph of Sec. II interesting.
In Sec. III, we apply the general analysis to the model (1) for Fe<sub>8</sub>. To keep the analysis tractable, we will assume that $`H_z/H_c`$ is small, although there is no reason why one could not apply the results of Sec. II numerically for arbitrary values of $`H_z`$. It may in fact be interesting to do this to further investigate the perfect lattice hypothesis. Further, as discussed in Sec. III, the duality property of $``$ immediately extends our work to large $`H_z`$ and small $`H_x`$.
## II General Formula for Diabolo in terms of Action Integrals
### A Summary of DPI method, critical curves, etc.
Very briefly, the DPI analysis proceeds as follows. (See Ref. and paper I for details.) Let the energy of an electron wavepacket in the equivalent tight-binding model be given by $`_{\mathrm{sc}}(q,m)`$ where $`q`$ and $`m`$ are the mean wavevector and position of the wavepacket, respectively. Holding $`m`$ fixed, we may think of $`_{\mathrm{sc}}(q,m)`$ as a dispersion relation for the local energy band. The mean velocity of the wavepacket is then given by $`v(q,m)=_{\mathrm{sc}}(q,m)/q`$. If we regard $`m`$ as a continuous variable, and approximate $`w_m`$, $`t_{m,m\pm 1}`$, and $`t_{m,m\pm 2}`$ by smooth functions $`w(m)`$, $`t_1(m)`$, and $`t_2(m)`$, then
$`_{\mathrm{sc}}(q,m)`$ $`=`$ $`w(m)+2t_1(m)\mathrm{cos}q+2t_2(m)\mathrm{cos}(2q),`$ (6)
$`v(q,m)`$ $`=`$ $`2\mathrm{sin}q(m)\left(t_1(m)+4t_2(m)\mathrm{cos}q(m)\right).`$ (7)
In close analogy with the continuum method, one can show that the quasiclassical solution to the wavefunction for a given energy $`E`$ is given by linear combinations of solutions of the form
$$C_m\frac{1}{\sqrt{v(m)}}\mathrm{exp}\left(i^mq(m^{})𝑑m^{}\right),$$
(8)
where $`q(m)`$ and $`v(m)`$ are given by
$$E=_{\mathrm{sc}}(q(m),m),v(m)=v(q(m),m).$$
(9)
It follows from Eq. (8) that the quasiclassical solution breaks down when $`v(m)=0`$—which we call a turning point, as in the continuum case. At these points, the solutions must be augmented by connection formulas. In addition to the values $`q=0`$ and $`q=\pi `$, $`v(m)`$ also vanishes at $`q=q^{}(m)=\mathrm{cos}^1\left(t_1(m)/4t_2(m)\right)`$. We will get turning points whenever $`E`$ equals $`U_0(m)`$, $`U_\pi (m)`$, or $`U_{}(m)`$, where these three functions are $`_{\mathrm{sc}}(0,m)`$, $`_{\mathrm{sc}}(\pi ,m)`$, and $`_{\mathrm{sc}}(q^{},m)`$ respectively. These three curves, which we call critical curves, collectively play the same role as the potential energy in the continuum case.
For Fe<sub>8</sub>, in the same reduced variables as in I \[$`\mu =m/\overline{J}`$, energies in units of $`k_1\overline{J}^2`$, $`\overline{J}=(J+\frac{1}{2})`$\], the on-site and hopping energies are given by
$`w(m)`$ $`=`$ $`\frac{1}{2}(1+\lambda )(1\mu ^2)2h_z\mu ,`$ (10)
$`t_1(m)`$ $`=`$ $`h_x(1\mu ^2)^{1/2},`$ (11)
$`t_2(m)`$ $`=`$ $`\frac{1}{4}(1\lambda )(1\mu ^2).`$ (12)
Here, $`\lambda =k_2/k_1`$, $`h_x=JH_x/\overline{J}H_c`$, $`h_z=JH_z/\overline{J}H_c`$, and $`H_c=2k_1J/g\mu _B`$. Thus the critical curves are given by
$`U_0(\mu )`$ $`=`$ $`1\mu ^22h_x(1\mu ^2)^{1/2}2h_z\mu ,`$ (13)
$`U_\pi (\mu )`$ $`=`$ $`1\mu ^2+2h_x(1\mu ^2)^{1/2}2h_z\mu ,`$ (14)
$`U_{}(\mu )`$ $`=`$ $`\lambda (1\mu ^2){\displaystyle \frac{h_x^2}{1\lambda }}2h_z\mu .`$ (15)
These curves are shown in Fig. 2. The energy $`U_\pi (m)`$ is the upper edge of the band $`_{\mathrm{sc}}(q,m)`$ for all $`m`$. The lower band edge is given by $`U_0(m)`$ for values of $`m`$ close to $`\pm J`$, and by $`U_{}(m)`$ for $`m`$ in the central region. In the central region, the second neighbour hopping element $`t_2(m)`$ is sufficiently large that the local energy band $`_{\mathrm{sc}}(q,m)`$ has its global minimum not at $`q=0`$, but at $`q^{}`$. Since the energy $`U_0`$ lies in the band, $`U_\pi >U_0>U_{}`$. By contrast, in the outer $`m`$ regions, $`_{\mathrm{sc}}(q,m)`$ has only one minimum (at $`q=0`$) and only one maximum (at $`q=\pi `$) for real $`q`$. The energy $`U_{}`$ lies outside the band, so that once again $`U_\pi >U_0>U_{}`$. The curves $`U_0(m)`$ and $`U_{}(m)`$ touch one another with a common tangent at $`m=\pm m^{}`$, where
$$m^{}=\overline{J}\left[1\frac{h_x^2}{(1\lambda )^2}\right]^{1/2}.$$
(16)
The turning point where $`E=U_{}`$ is special if it happens at a value of $`m`$ where $`U_{}`$ lies below the lower band edge. The wavevector $`q^{}`$ at the turning point is then complex, and the wavefunction $`C_m`$ changes from a decaying (or growing) exponential on one side to a decaying (or growing) exponential with an oscillatory envelope on the other side. These turning points are the new feature caused by second neightbor hopping that we referred to in Sec. I.
### B Nature of asymmetric double well wavefunctions
To understand how the DPI solutions are to be used to find the eigenstates, it is useful to discuss the corresponding problem for an asymmetric double well in the continuum case. Suppose the potential is as drawn in Fig. 3. The potential minima are at $`x_{0\pm }`$, and for the energy $`E`$ drawn, the turning points are at $`x_a^{}`$ and $`x_b^{}`$ in the left well, and $`x_b^{\prime \prime }`$ and $`x_a^{\prime \prime }`$ in the right well. We use the quasiclassical approximation to find a wavefunction $`\psi ^{}(x)`$ on the left hand side as follows. First, we choose the solution which decays exponentially as $`xx_a^{}\mathrm{}`$. This solution is matched via connection formulas at $`x_a^{}`$ to an oscillatory solution in the region $`x_a^{}<x<x_b^{}`$. We then use connection formulas at $`x_b^{}`$ to find the quasiclassical solution in the region $`x>x_b^{}`$. We repeat this procedure on the right hand side to find a wavefunction $`\psi ^{\prime \prime }(x)`$ that decays to zero as $`xx_a^{\prime \prime }\mathrm{}`$. The last step is to demand that the wavefunctions $`\psi ^{}(x)`$ and $`\psi ^{\prime \prime }(x)`$ be the same in the central region, i.e., in the vicinity of $`x=0`$. This demand will be unsatisfiable for an arbitrarily chosen energy $`E`$, and will provide one with the eignevalue condition.
There are two remarks that we wish make about the above procedure. The first remark concerns the basic nature of the solution. In general, in the central region, the left solution $`\psi ^{}(x)`$ will be a linear combination of two parts, $`\psi _d^{}(x)`$ and $`\psi _g^{}(x)`$, that are exponentially decaying and growing as $`xx_b^{}`$ increases, respectively. Likewise the right solution, $`\psi ^{\prime \prime }`$, will be a sum of parts that decay ($`\psi _d^{\prime \prime }`$) and grow ($`\psi _g^{\prime \prime }`$) as $`x_b^{\prime \prime }x`$ increases. The key point is that the growing parts $`\psi _g^{}`$ and $`\psi _g^{\prime \prime }`$ must be present in an eigenstate, for without them, there is no way that the values and slopes of $`\psi ^{}`$ and $`\psi ^{\prime \prime }`$ could be made to agree at $`x=0`$, say.
The second remark is technical. If the potential well is reasonably parabolic near the minimum at $`x_0`$, the Schrödinger equation can be solved directly for any choice of $`E`$ in terms of parabolic cylinder functions, and we can always find a linear combination that will decay to zero as $`xx_a^{}\mathrm{}`$. This linear combination will have both growing and decaying pieces as $`xx_b^{}`$ grows. In this way we can obtain the wavefucnction $`\psi ^{}(x)`$ on the entire left hand side without using connection formulas at $`x_a^{}`$ or $`x_b^{}`$. Once one has found $`\psi ^{}(x)`$ sufficiently far to the right of $`x_b^{}`$ in this way, one can write the parabolic cylinder functions in quasiclassical form which can then be extended in this form all the way to $`x=0`$. The right-hand wavefunction can be treated in the same way. This device leads to considerable savings in labor.
We now apply these ideas to our problem. In what follows, we will denote quantities pertaining to the left hand solution or the left hand side of the well by either a single prime or a suffix $``$, while analogous right-hand quantities will carry a double prime or a $`+`$ suffix. We consider an energy $`E`$ as drawn in Fig. 4, which leads to turning points at $`m_a^{}`$, $`m_b^{}`$, $`m_c^{}`$ on the left hand side, and $`m_c^{\prime \prime }`$, $`m_b^{\prime \prime }`$, and $`m_a^{\prime \prime }`$ on the right. The wavefunction $`C_m^{}`$ will decay away from the well bottom as $`mm_a^{}`$ decreases, oscillate in the classically allowed region $`m_a^{}<m<m_b^{}`$ . In the region just to the right of $`m_b^{}`$ it will consist of a decaying part and a growing part. The new feature will be encountered at $`m_c^{}`$ where $`E=U_{}`$. For $`m>m_c^{}`$, both the growing and decaying parts will acquire oscillatory envelopes. Similar remarks apply to the right side wavefunction $`C_m^{\prime \prime }`$.
### C DPI wavefunction in the leftmost forbidden region
We are now ready to find the wavefunction explicitly. Let us start constructing $`C_m^{}`$ from the left. The Hamilton-Jacobi equation in Eq. (9) has the general solution
$`\mathrm{cos}q(m)`$ $`=`$ $`{\displaystyle \frac{t_1(m)\pm [t_1^2(m)4t_2(m)f(m)]^{1/2}}{4t_2(m)}};`$ (17)
$`f(m)`$ $`=`$ $`w(m)2t_2(m)E.`$ (18)
This leads to four values of $`q(m)`$ for any $`E`$, since if $`q`$ is a solution, so is $`q`$. In the region $`mm_a^{}`$, since $`E<U_0`$, all four solutions are pure imaginary. We write the two which lead to decaying wavefunctions as $`\mu \mu _a^{}`$ becomes large and negative as
$$q_{1,2}=i\kappa _{1,2}(\mu ),$$
(19)
where $`\kappa _2>\kappa _1>0`$, and the corresponding DPI solutions as
$`C_{m,1}^{}`$ $`=`$ $`A^{}|v_1(m)|^{1/2}\mathrm{exp}\left(i{\displaystyle ^m}q_1(m^{})𝑑m^{}\right),`$ (20)
$`C_{m,2}^{}`$ $`=`$ $`B^{}|v_2(m)|^{1/2}\mathrm{exp}\left(i{\displaystyle _{m_c^{}}^m}q_2(m^{})𝑑m^{}\right),`$ (21)
with $`v_i(m)=v(q_i(m),m)`$. We take $`A^{}`$ and $`B^{}`$ to be real without any loss of generality. Note that we have left the lower limit of the phase integral for $`C_{m,1}^{}`$ unspecified and written it as $`m_c^{}`$ for $`C_{m,2}^{}`$. The reasons for this will become clear shortly.
To see how these two solutions behave as $`m`$ approaches $`m_a^{}`$, and continues beyond this point, let us note that
$$\mathrm{cosh}\kappa _{1,2}=\frac{|t_1|[t_1^24t_2f]^{1/2}}{4t_2}.$$
(22)
As $`mm_a^{}`$, $`\mathrm{cosh}\kappa _11`$, i.e., $`\kappa _10`$, while $`\mathrm{cosh}\kappa _21+2|t_1|/4t_2>1`$. As we cross the point $`m_a^{}`$, $`q_1`$ will become real, while $`q_2`$ will continue to be pure imaginary and large. Thus the solution $`C_{m,2}^{}`$ continues to hold at $`m_a^{}`$, while $`C_{m,1}^{}`$ breaks down at $`m_a^{}`$ , and must be related to a the solution for $`m>m_a^{}`$ by a connection formula. It is clear that the wavevector(s) for $`C_{m,2}^{}`$ will continue to be given by Eq. (22) as $`mm_b^{}`$, while those for $`C_{m,1}^{}`$ will again approach zero, necessitating the use of connection formulas to go on to $`m>m_b^{}`$. It is here that the technical remark about sidestepping the use of connection formulas that was made in connection with the continuum antisymmetric double-well is relevant. The solution $`C_{m,1}^{}`$ can be approximated by a harmonic oscillator wavefunction provided the energy $`E`$ is not very far from the minima of $`U_0`$. The asymptotic forms of this wavefunction give us the quasiclassical wavefunction in the regions $`m<m_a^{}`$ and $`m_b^{}`$ more simply. We therefore turn to this subproblem.
### D Jumping across the potential well
The assumption that $`EU_0(\mu _{0\pm })`$ is not very large, means that $`q_1(m)`$ is never very far from zero, and we may expand $`_{\mathrm{sc}}(q,m)`$ in powers of $`q`$ and $`m+m_0^{}`$. As in I, we write
$$_{\mathrm{sc}}(q,m)=E_{}+\frac{1}{2M_{}}q^2+\frac{1}{2}M_{}\omega _{}^2(mm_{0,})^2+\mathrm{}.$$
(23)
Since $`q`$ and $`m`$ are conjugate variables, we can write $`C_{m,1}^{}`$ as the solution to the differential equation
$$_{\mathrm{sc}}(i_m,m)C_{m,1}^{}=EC_{m,1}^{}.$$
(24)
Introducing two new variables $`z`$ and $`\nu ^{}`$ by the equations
$`m`$ $`=`$ $`m_0+(2M_{}\omega _0)^{1/2}z,`$ (25)
$`E`$ $`=`$ $`E_{}+(\nu ^{}+\frac{1}{2})\omega _0,`$ (26)
within the approximation (23), the differential equation becomes that for the parabolic cylinder functions:
$$\left[\frac{d^2}{dz^2}+\left(\nu ^{}+\frac{1}{2}\frac{z^2}{4}\right)\right]C_{m,1}^{}=0.$$
(27)
If we take as the two linearly independent solutions the standard forms $`D_\nu ^{}(z)`$ and $`D_\nu ^{}(z)`$ , the former must be rejected as it diverges for $`z\mathrm{}`$. We accordingly write
$$C_{m,1}^{}=A^{}(1)^{\mathrm{}^{}}D_\nu ^{}(z),$$
(28)
where $`A^{}`$ is the constant in Eq. (20), and the additional factor $`(1)^{\mathrm{}^{}}`$, where we shall define $`\mathrm{}^{}`$ shortly, is another constant introduced for later convenience.
As $`z\mathrm{}`$, $`D_\nu ^{}(z)(z)^\nu ^{}e^{z^2/4}`$, and one can show with a little work that \[modulo the factor $`(1)^{\mathrm{}^{}}`$\] this is indeed the DPI form for $`C_{m,1}^{}`$ with the approximation (23) . For $`z+\mathrm{}`$, on the other hand,
$`D_\nu ^{}(z)`$ $``$ $`\mathrm{cos}(\pi \nu ^{})z^\nu ^{}e^{z^2/4}\left(1{\displaystyle \frac{\nu ^{}(\nu ^{}1)}{2z^2}}+\mathrm{}\right)`$ (30)
$`+{\displaystyle \frac{\sqrt{2\pi }}{\mathrm{\Gamma }(\nu ^{})}}z^{\nu ^{}1}e^{z^2/4}\left(1+{\displaystyle \frac{(\nu ^{}+1)(\nu ^{}+2)}{2z^2}}+\mathrm{}\right).`$
Note that this form has both decaying and growing components. In fact, the latter component vanishes only if $`\nu ^{}`$ is a positive a integer. As explained earlier, it is essential for our DPI solution $`C_m^{}`$ to contain a growing component. We therefore allow for its presence by writing
$$\nu ^{}=\mathrm{}^{}+\frac{ϵ^{}}{\omega _0},$$
(31)
where $`\mathrm{}^{}`$ is a positive integer, and $`ϵ^{}`$ is a shift defined to lie in the interval $`(1/2,1/2)\omega _0`$. In fact, we expect $`ϵ^{}`$ to be very close to zero for any state in which there is large probaility of finding the particle in the left well. We can make the vanishing of the growing component in $`C_{m,1}^{}`$ more manifest by writing
$$\frac{1}{\mathrm{\Gamma }(\nu ^{})}=\frac{\mathrm{sin}(\pi \nu ^{})}{\pi }\mathrm{\Gamma }(1+\nu ^{})(1)^{\mathrm{}^{}+1}(\mathrm{}^{}!)\frac{ϵ^{}}{\omega _0}.$$
(32)
Combining Eqs. (28)–(32), we thus find that for $`m`$ beyond $`m_b^{}`$,
$$C_{m,1}^{}A^{}\left(\mathrm{cos}\frac{\pi ϵ^{}}{\omega _0}z^\nu ^{}e^{z^2/4}\sqrt{2\pi }(\mathrm{}^{}!)\frac{ϵ^{}}{\omega _0}z^{\nu ^{}1}e^{z^2/4}\right).$$
(33)
### E DPI form in ordinary forbidden region
The next step is to write the solution for $`C_m^{}`$ in such a way that it holds in the entire region $`m_b^{}<m<m_c^{}`$ . For $`C_{m,2}^{}`$, Eq. (21) already meets this demand, since just as at the turning point $`m_a^{}`$, $`q_2`$ stays imaginary and negative as $`m`$ passes through $`m_c^{}`$. For $`C_{m,1}^{}`$, on the other hand, Eq. (33) only holds in a region where the parabolic approximation to $`U_0`$ is good, and may not hold near $`m_c^{}`$. It is clear, however, that the wavevector $`q_1`$ associated with $`C_{m,1}^{}`$ is again given by $`i\kappa _1`$ with $`\kappa _1`$ given by Eq. (22). Hence it must be possible to write $`C_{m,1}^{}`$ in the DPI form
$$C_{m,1}^{}=|v_1(m)|^{1/2}\left[Q^{}\mathrm{exp}\left(i_{m_b^{}}^mq_1(m^{})𝑑m^{}\right)+R^{}\mathrm{exp}\left(i_{m_b^{}}^mq_1(m^{})𝑑m^{}\right)\right],$$
(34)
where $`Q^{}`$ and $`R^{}`$ are coefficients which we expect to be proportional to $`A^{}`$. To find these, let us first calculate the phase integral in the parabolic approximation. Since $`q_1=0`$ at $`m_b^{}`$, we have
$$E=_{\mathrm{sc}}(0,m_b^{})=_{\mathrm{sc}}(i\kappa _1(m),m),(m>m_b^{}).$$
(35)
Using Eq. (23), we get
$$\kappa _1(m)M_{}\omega _0[(mm_0)^2(m_b^{}m_0)^2)^{1/2}].$$
(36)
Therefore, if $`(mm_b^{})(m_b^{}m_0)`$, we obtain
$`\mathrm{exp}\left({\displaystyle _{m_b^{}}^m}\kappa _1(m^{})𝑑m^{}\right)`$ $`=`$ $`\left(2{\displaystyle \frac{mm_0}{m_b^{}m_0}}\right)^{\nu ^{}+{\scriptscriptstyle \frac{1}{2}}}e^{{\scriptscriptstyle \frac{1}{2}}\left(\nu ^{}+{\scriptscriptstyle \frac{1}{2}}\right)}`$ (38)
$`\times \mathrm{exp}\left({\displaystyle \frac{1}{2}}M_{}\omega _0(mm_0)^2\right),`$
where we have used the fact that
$$\frac{1}{2}M_{}\omega _0^2(m_b^{}m_0)^2=(\nu ^{}+\frac{1}{2})\omega _0.$$
(39)
Next we note that (a) from the definition of $`z`$, Eq. (25), the last exponential in Eq. (38) is nothing but $`\mathrm{exp}(z^2/4)`$, and that (b) $`|v_1(m)|\kappa _1(m)/M_{}\omega _0(mm_0)`$. Since $`mm_0z`$, it follows that the firat term in Eq. (34) varies as $`z^\nu ^{}\mathrm{exp}(z^2/4)`$, i.e., precisely as the first term in Eq. (33). Thus $`Q^{}`$ is indeed proportional to $`A^{}`$, and a little algebra plus the use of Eq. (39) gives
$$Q^{}=\left(\frac{\omega _0}{2M_{}}\right)^{1/4}e^{{\scriptscriptstyle \frac{1}{2}}\left(\nu ^{}+{\scriptscriptstyle \frac{1}{2}}\right)}(\nu ^{}+\frac{1}{2})^{{\scriptscriptstyle \frac{1}{2}}\left(\nu ^{}+{\scriptscriptstyle \frac{1}{2}}\right)}\mathrm{cos}\frac{\pi ϵ^{}}{\omega _0}A^{}\alpha ^{}A^{}.$$
(40)
In the same way, we can show that
$$R^{}=\left(\frac{\omega _0}{2M_{}}\right)^{1/4}e^{{\scriptscriptstyle \frac{1}{2}}\left(\nu ^{}+{\scriptscriptstyle \frac{1}{2}}\right)}(\nu ^{}+\frac{1}{2})^{{\scriptscriptstyle \frac{1}{2}}\left(\nu ^{}+{\scriptscriptstyle \frac{1}{2}}\right)}\sqrt{2\pi }(\mathrm{}^{}!)\frac{ϵ^{}}{\omega _0}A^{}\beta ^{}\frac{ϵ^{}}{\omega _0}A^{}.$$
(41)
The definitions of the factors $`\alpha ^{}`$ and $`\beta ^{}`$ which we have introduced for later convenience can be read off these equations. Note that since we took $`A^{}`$ to be real, $`Q^{}`$ and $`R^{}`$ are also real.
To summarize where we are, the complete DPI solution for $`C_m^{}`$ in the region $`m_b^{}<m<m_c^{}`$ is given by the sum of Eqs. (21) and (34). The terms in $`B^{}`$ and $`R^{}`$ are exponentially growing with increasing $`m`$, while the term in $`Q^{}`$ is exponentially decreasing. The next step is to connect this solution to the DPI form in the region $`m_c^{}<m`$. As already stated, the turning point $`m_c^{}`$ is the irregular one under the barrier, where exponentially growing and decaying wavefunctions acquire oscillatory envelopes as it is crossed.
Before we use the connection formulas at $`m_c^{}`$, it is useful to see how the quasiclassical wavevector behaves near this point. Since $`\mathrm{cos}q(m_c^{})=\mathrm{cos}q^{}=t_1(m_c^{})/4t_2(m_c^{})`$, it follows that the discriminant in Eq. (17) vanishes at $`m_c^{}`$, and both $`q_1`$ and $`q_2`$ tend to the same value $`i\kappa _c^{}`$, where
$$\mathrm{cosh}\kappa _c^{}=(|t_1|/4t_2)_{m=m_c^{}}.$$
(42)
As we cross $`m_c^{}`$, the discriminant in Eq. (17) becomes negative and $`\mathrm{cos}q(m)`$ \[and therefore $`q(m)`$\] becomes complex. We separate $`q(m)`$ into its real and imaginary parts, and write two distinct solutions as
$$q_{d,g}(m)=\pm i\kappa (m)+\chi (m),(m>m_c^{})$$
(43)
where both $`\kappa `$ and $`\chi `$ are real and positive. The subscripts $`d`$ and $`g`$ stand for ‘decaying’ and ‘growing’.
### F DPI form in oscillatory forbidden region
The connection formulas to be used at $`m_c^{}`$ were derived in Sec. IV of Ref. . The parts of $`C_m^{}`$ multiplying $`B^{}`$, $`Q^{}`$, and $`R^{}`$ correspond respectively to the cases there labelled $`(\sigma _1,\sigma _2)=(+1,1)`$, $`(1,+1)`$, and $`(1,1)`$. The $`B^{}`$ part is given by
$$C_{m,2}^{}=B^{}[\frac{1}{\sqrt{s_g(m)}}\mathrm{exp}(i_{m_c^{}}^mq_g(m^{})dm^{}\frac{\pi }{2})+\mathrm{c}.\mathrm{c}.],$$
(44)
where
$`s_g(m)`$ $`=`$ $`8t_2(m)\mathrm{sinh}\kappa (m)\mathrm{sin}\chi (m)\mathrm{sin}q_g(m)`$ (45)
$`=`$ $`8t_2\mathrm{sinh}\kappa \mathrm{sin}\chi (\mathrm{sin}\chi \mathrm{cosh}\kappa i\mathrm{cos}\chi \mathrm{sinh}\kappa ).`$ (46)
Likewise, the coefficient of $`Q^{}`$, which we call part 1a, connects to
$$C_{m,1a}^{}=e^\mathrm{\Gamma }^{}Q^{}[\frac{1}{\sqrt{s_d(m)}}\mathrm{exp}\left(i_{m_c^{}}^mq_d(m^{})dm^{}\right)+\mathrm{c}.\mathrm{c}.],$$
(47)
where
$$s_d(m)=8t_2(m)\mathrm{sinh}\kappa (m)\mathrm{sin}\chi (m)\mathrm{sin}q_d(m)=[s_g(m)]^{},$$
(48)
and $`\mathrm{\Gamma }^{}`$ is the phase integral,
$$\mathrm{\Gamma }^{}=_{m_b^{}}^{m_c^{}}\kappa (m^{})𝑑m^{},$$
(49)
which we acquire in changing the lower limits of the $`m`$ integrals from $`m_b^{}`$ to $`m_c^{}`$. Lastly, the term in $`R^{}`$, which we call part 1b, connects to
$$C_{m,1b}^{}=\frac{1}{2}e^\mathrm{\Gamma }^{}R^{}[\frac{1}{\sqrt{s_g(m)}}\mathrm{exp}\left(i_{m_c^{}}^mq_g(m^{})dm^{}\right)+\mathrm{c}.\mathrm{c}.].$$
(50)
Equations (44), (47), and (50), give us the complete wavefunction $`C_m^{}`$ in the central region near $`m=0`$. To simplify the writing, we denote
$$\mathrm{\Phi }_{\lambda _1\lambda _2}^{}(m)=_{m_c^{}}^m[\lambda _1\kappa (m^{})+i\lambda _2\chi (m^{})]𝑑m^{},$$
(51)
where $`\lambda _1`$ and $`\lambda _2`$ can be $`\pm 1`$ independently. In other words, the subscripts on $`\mathrm{\Phi }^{}`$ give the signs of the real and imaginary parts \[$`\mathrm{\Phi }_{++}^{}`$ is the integral of $`(\kappa +i\chi )`$, $`\mathrm{\Phi }_+^{}`$ that of $`(\kappa +i\chi )`$, etc.\]. The complete DPI solution for $`C_m^{}`$ can then be written as
$$C_m^{}=\left[e^\mathrm{\Gamma }^{}Q^{}\frac{e^{\mathrm{\Phi }_+^{}(m)}}{\sqrt{s_g^{}(m)}}+\left(\frac{1}{2}e^\mathrm{\Gamma }^{}R^{}iB^{}\right)\frac{e^{\mathrm{\Phi }_{++}^{}(m)}}{\sqrt{s_g(m)}}\right]+\mathrm{c}.\mathrm{c}.$$
(52)
The wavefunction from the right, $`C_m^{\prime \prime }`$ can now be written down at once. We define quantities with double primes in exact correspondence with those for $`C_m^{}`$. The analog of Eq. (52) is then
$$C_m^{\prime \prime }=\left[e^{\mathrm{\Gamma }^{\prime \prime }}Q^{\prime \prime }\frac{e^{\mathrm{\Phi }_+^{\prime \prime }(m)}}{\sqrt{s_g^{}(m)}}+\left(\frac{1}{2}e^{\mathrm{\Gamma }^{\prime \prime }}R^{\prime \prime }iB^{\prime \prime }\right)\frac{e^{\mathrm{\Phi }_{++}^{\prime \prime }(m)}}{\sqrt{s_g(m)}}\right]+\mathrm{c}.\mathrm{c}.$$
(53)
The only issue requiring any thought is what the sign suffixes in $`\mathrm{\Phi }^{\prime \prime }`$ should mean. By defining a new variable $`n=m`$, so that the problem for $`C_m^{\prime \prime }`$ becomes completely isomorphic to that for $`C_m^{}`$, and then transforming back to $`m`$, one can show that
$`\mathrm{\Gamma }^{\prime \prime }`$ $`=`$ $`{\displaystyle _{m_c^{\prime \prime }}^{m_b^{\prime \prime }}}\kappa (m^{})𝑑m^{},`$ (54)
$`\mathrm{\Phi }_+^{\prime \prime }(m)`$ $`=`$ $`{\displaystyle _m^{m_c^{\prime \prime }}}[\kappa (m^{})+i\chi (m^{})]𝑑m^{},`$ (55)
etc. Note that since $`m<m_c^{\prime \prime }`$ in the center, and $`m_c^{\prime \prime }<m_b^{\prime \prime }`$, these integrals are written so that the lower limit is less than the upper limit. Thus the suffixes on $`\mathrm{\Phi }^{\prime \prime }`$ give the true signs of its real and imaginary parts.
### G Matching of left and right wavefunctions
It remains to see if Eqs. (52) and (53) are the same function. We note that $`\mathrm{\Phi }_+^{\prime \prime }(m)`$ has the same integrand as $`\mathrm{\Phi }_+^{}(m)`$ if $`m`$ is taken to be the upper limit for both integrals. We further note that the remaining $`m`$ dependence in both terms is $`(s_g^{})^{1/2}`$. Similar remarks apply to the $`\mathrm{\Phi }_+^{}`$ and $`\mathrm{\Phi }_+^{\prime \prime }`$ terms. Thus, we conclude that $`C_m^{}`$ will equal $`C_m^{\prime \prime }`$ if the following conditions are obeyed:
$`e^{\mathrm{\Gamma }^{\prime \prime }}Q^{\prime \prime }e^{\mathrm{\Phi }_+^{\prime \prime }(m)}`$ $`=`$ $`\left(\frac{1}{2}e^\mathrm{\Gamma }^{}R^{}+iB^{}\right)e^{\mathrm{\Phi }_+^{}(m)},`$ (56)
$`e^\mathrm{\Gamma }^{}Q^{}e^{\mathrm{\Phi }_+^{}(m)}`$ $`=`$ $`\left(\frac{1}{2}e^{\mathrm{\Gamma }^{\prime \prime }}R^{\prime \prime }+iB^{\prime \prime }\right)e^{\mathrm{\Phi }_+^{\prime \prime }(m)}.`$ (57)
To simplify these conditions we note that
$`\mathrm{\Phi }_+^{}(m)\mathrm{\Phi }_+^{\prime \prime }(m)`$ $`=`$ $`{\displaystyle _{m_c^{}}^{m_c^{\prime \prime }}}[\kappa (m)i\chi (m)]𝑑m,`$ (58)
$``$ $`\mathrm{\Gamma }_ci\mathrm{\Lambda }_c,`$ (59)
where $`\mathrm{\Gamma }_c`$ and $`\mathrm{\Lambda }_c`$ are the real and imaginary parts of the integral, and the subscript ‘c’ indicates that the integrals extend over the central region of $`m`$. Equations (56) and (57) can now be written as
$`Q^{\prime \prime }`$ $`=`$ $`\left(\frac{1}{2}e^\mathrm{\Gamma }^{}R^{}+iB^{}\right)e^{\mathrm{\Gamma }^{\prime \prime }}e^{\mathrm{\Gamma }_ci\mathrm{\Lambda }_c},`$ (60)
$`Q^{}`$ $`=`$ $`\left(\frac{1}{2}e^{\mathrm{\Gamma }^{\prime \prime }}R^{\prime \prime }+iB^{\prime \prime }\right)e^\mathrm{\Gamma }^{}e^{\mathrm{\Gamma }_ci\mathrm{\Lambda }_c}.`$ (61)
If we recall that \[see Eqs. (40) and (41)\] $`Q^{}`$ and $`R^{}`$ are proportional to $`A^{}`$, and likewise for $`Q^{\prime \prime }`$ and $`R^{\prime \prime }`$, these equations are two complex equations in the four real quantities $`A^{}`$, $`B^{}`$, $`A^{\prime \prime }`$ and $`B^{\prime \prime }`$. To solve them we first note that the imaginary parts on the right hand sides must vanish. This yields
$`B^{}`$ $`=`$ $`\frac{1}{2}R^{}e^\mathrm{\Gamma }^{}\mathrm{tan}\mathrm{\Lambda }_c,`$ (62)
$`B^{\prime \prime }`$ $`=`$ $`\frac{1}{2}R^{\prime \prime }e^{\mathrm{\Gamma }^{\prime \prime }}\mathrm{tan}\mathrm{\Lambda }_c.`$ (63)
Substituting these in Eqs. (60) and (61), we obtain after some simplification
$`R^{}=2e^{\mathrm{\Gamma }_G}\mathrm{cos}\mathrm{\Lambda }_cQ^{\prime \prime },`$ (64)
$`R^{\prime \prime }=2e^{\mathrm{\Gamma }_G}\mathrm{cos}\mathrm{\Lambda }_cQ^{},`$ (65)
where $`\mathrm{\Gamma }_G`$ is the total Gamow factor
$$\mathrm{\Gamma }_G=_{m_b^{}}^{m_b^{\prime \prime }}\kappa _1(m)𝑑m,$$
(66)
and the subscript ‘1’ on $`\kappa `$ is to remind us that we must use the imaginary part of the wavevector that goes to zero at the turning points $`m_b^{}`$ and $`m_b^{\prime \prime }`$.
### H The eigenvalue condition, and the diabolo
The simplest way of solving Eqs. (64) and (65) is terms of the ratios $`\alpha ^{}`$ and $`\beta ^{}`$ defined in Eqs. (40) and (41), and the analogous ratios $`\alpha ^{\prime \prime }`$ and $`\beta ^{\prime \prime }`$. Equating the products of the left hand and right hand sides, and simplifying a little, we get
$$ϵ^{}ϵ^{\prime \prime }=4\frac{\alpha ^{}\alpha ^{\prime \prime }}{\beta ^{}\beta ^{\prime \prime }}\omega _0\omega _{0+}\mathrm{cos}^2\mathrm{\Lambda }_ce^{2\mathrm{\Gamma }_G}.$$
(67)
This is our eigenvalue condition. To understand it better, we first note that the right hand side is exponentially small on account of the square of the Gamow factor $`e^{\mathrm{\Gamma }_G}`$. Thus, ignoring for the moment the possibility that $`\mathrm{cos}\mathrm{\Lambda }_c`$ may vanish, either both $`ϵ^{}`$ and $`ϵ^{\prime \prime }`$ must be of order $`e^{\mathrm{\Gamma }_G}`$, or at the other extreme, one must be of order unity, and the other of order $`e^{2\mathrm{\Gamma }_G}`$. Suppose that $`ϵ^{}=O(e^{2\mathrm{\Gamma }_G})`$, and $`ϵ^{\prime \prime }=O(1)`$. Let us take $`A^{}=1`$. Then from Eqs. (40) and (41), we see that $`Q^{}=O(1)`$, while $`R^{}=O(e^{2\mathrm{\Gamma }_G})`$. It then follows from Eqs. (64) and (65) that $`Q^{\prime \prime }R^{\prime \prime }=O(e^{\mathrm{\Gamma }_G})`$, and in turn from the double primed analogs of Eqs. (40) and (41) that $`A^{\prime \prime }=O(e^{\mathrm{\Gamma }_G})`$. Lastly, since $`\mathrm{\Gamma }^{\prime \prime }<\mathrm{\Gamma }_G`$, Eq. (63) implies that $`B^{\prime \prime }`$ is also exponentially small . Hence, it follows that the entire wavefunction on the right hand side of the well, $`C_m^{\prime \prime }`$ is exponentially small compared to the left hand part $`C_m^{}`$. In other words, there is negligible mixing of the states in the left and right hand well. This is exactly what we expect when the energies of the two states in the absence of tunneling differ by much more than the tunneling matrix element itself.
The more interesting case, therefore, is that in which both $`ϵ^{}`$ and $`ϵ^{\prime \prime }`$ are of order $`e^{\mathrm{\Gamma }_G}`$. In the defining equations for the $`\alpha `$’s and $`\beta `$’s, Eqs. (40) and (41), we can then neglect the $`ϵ`$’s to very good approximation. This yields
$`{\displaystyle \frac{\alpha ^{}}{\beta ^{}}}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2\pi }\mathrm{}^{}!}}\left(\mathrm{}^{}+\frac{1}{2}\right)^{\mathrm{}^{}+{\scriptscriptstyle \frac{1}{2}}}e^{(\mathrm{}^{}+{\scriptscriptstyle \frac{1}{2}})},`$ (68)
$`=`$ $`{\displaystyle \frac{g_{\mathrm{}^{}}}{2\pi }},`$ (69)
where $`g_n`$ is the standard curvature correction in the phase integral expression for the tunnel splitting (see e.g., Eq. (4.10) of paper I). This quantity tends to 1 rapidly as $`n`$ gets large: $`g_0=(\pi /e)^{1/2}1.075`$, $`g_11.028`$, $`g_21.017`$, $`\mathrm{}`$. The right hand side of Eq. (67) can thus be written as $`\mathrm{\Delta }^2(\mathrm{}^{},\mathrm{}^{\prime \prime })/4`$, where
$$\mathrm{\Delta }(\mathrm{}^{},\mathrm{}^{\prime \prime })=\frac{2}{\pi }(g_{\mathrm{}^{}}g_{\mathrm{}^{\prime \prime }})^{1/2}(\omega _0\omega _{0+})^{1/2}e^{\mathrm{\Gamma }_G}\mathrm{cos}\mathrm{\Lambda }_c.$$
(70)
We further define
$`ϵ`$ $`=`$ $`\frac{1}{2}(ϵ^{}+ϵ^{\prime \prime })=E\frac{1}{2}\left(E_{}+E_++(\mathrm{}^{}+\frac{1}{2})\omega _0+(\mathrm{}^{\prime \prime }+\frac{1}{2})\omega _{0+}\right),`$ (71)
$`\delta `$ $`=`$ $`ϵ^{\prime \prime }ϵ^{}=\left(E_{}E_++(\mathrm{}^{}+\frac{1}{2})\omega _0(\mathrm{}^{\prime \prime }+\frac{1}{2})\omega _{0+}\right).`$ (72)
With these definitions, Eq. (67) can be rewritten as
$$ϵ=\pm \frac{1}{2}[\delta ^2+\mathrm{\Delta }^2(\mathrm{}^{},\mathrm{}^{\prime \prime })]^{1/2}.$$
(73)
Equation (73) is the complete formal solution to the problem of tunneling in an asymmetric double well. Along with Eqs. (70)–(72), it is the analog of the general phase integral formula for the tunnel splitting in a symmetric double well that we found in I \[see Eq. (4.38) there\]. Since there is no great need for having the final answer in simple closed form for the specfic problem (1), and since the general procedure is fully explained in Sec. IV.E of I, we do not bother to extract the singular $`\mathrm{ln}J`$ parts of the $`\mathrm{\Gamma }_G`$ integral.
It is immediately obvious that the eigenvalues in Eq. (73) are what we would get from a two level system Hamiltonian
$$_{\mathrm{TLS}}=\frac{1}{2}\left(\begin{array}{cc}\delta & \mathrm{\Delta }(\mathrm{}^{},\mathrm{}^{\prime \prime })\\ \mathrm{\Delta }(\mathrm{}^{},\mathrm{}^{\prime \prime })& \delta \end{array}\right),$$
(74)
which is of course, just what we would expect. The quantity $`ϵ`$ is the energy measured from a convenient reference point, while $`\delta `$, which depends on the fields $`h_x`$, $`h_z`$, and the quantum numbers $`\mathrm{}^{}`$ and $`\mathrm{}^{\prime \prime }`$ of the states whose mixing is being examined, is the offset between these energy levels in the absence of tunneling. Equation (70) gives the tunneling amplitude between these levels when the offset is small, i.e., when the two levels are in approximate degeneracy. Note that although this amplitude is defined even for relatively large offsets—offsets comparable to the intrawell spacings $`\omega _{0\pm }`$—and indeed is not very sensitive to the value of the offset, the concept of tunneling is physically sensible and useful only when the offset is comparable to or less than the amplitude $`\mathrm{\Delta }`$. If $`\delta \mathrm{\Delta }`$, we get $`ϵ\pm \delta /2`$, i.e., $`ϵ^{\prime \prime }\delta `$, $`ϵ^{}\mathrm{\Delta }^2/\delta `$, or the other way around. Then by the argument given after Eq. (67), the mixing between the wells is negligible.
One may wonder if the above conclusion does not invalidate the entire calculation. After all, we defined the shifts $`ϵ^{}`$ and $`ϵ^{\prime \prime }`$ assuming the wells were parabolic. Surely the corrections to the energies from cubic and higher order corrections to the potential are far larger than $`\mathrm{\Delta }`$. Since the offset between the levels must be tuned with exponential sensitivity, should not we know the energies of the levels before tunneling to the same sensitivity? The answer is no. The reason can be seen from Eq. (33). As we have seen, the amplitude of the growing part of the wavefunction plays a key role in the tunneling. From Eq. (33), we see that this amplitude is only linearly dependent on $`ϵ^{}`$. Hence, a small error in locating the absolute position of the level has little effect on the computed value of $`\mathrm{\Delta }`$. To say it another way, even though we know that the bottoms of the wells must be tuned to exponential accuracy to get significant mixing between the two wells, we cannot and need not determine the center of gravity of the two levels, $`ϵ`$, to the same accuracy to determine the tunneling amolitude itself. This feature is also present in the symmetric case studied in I, although there it is not so apparent, since Herring’s formula gives $`\mathrm{\Delta }`$ directly without making reference to the absolute energy level.
Secondly, it should be noted that Eq. (73) is nothing but the equation for the diabolo. The splitting vanishes only when $`\delta `$ and $`\mathrm{\Delta }`$ both vanish, which furnish the two conditions required to determine the diabolical point. Since both $`\delta `$ and $`\mathrm{\Delta }`$ will in general have linear terms in the deviation from the diabolical point, the energy surface is a double cone as asserted earlier.
Thirdly, let us ask if we recover the results of paper I in the symmetric case, i.e., when $`h_z=0`$. Then $`E_+=E_{}`$, $`M_+=M_{}`$, and $`\omega _0=\omega _{0+}`$. The only sensible case is $`\mathrm{}^{}=\mathrm{}^{\prime \prime }n`$, so that $`ϵ^{}=ϵ^{\prime \prime }`$, and $`\delta =0`$. The splitting is (up to an irrelevant sign) $`\mathrm{\Delta }(n,n)`$, which is precisely the tunnel splitting $`\mathrm{\Delta }_n`$ computed in paper I. In addition, however, we now have more explicit information about the wavefunction. Proceeding as before, we see that $`Q^{}A^{}1`$, $`R^{}e^{\mathrm{\Gamma }_G}`$, and $`B^{}e^{(\mathrm{\Gamma }_G\mathrm{\Gamma }^{})}`$. (The double primed quantities are equal to their single primed counterparts.) The conclusion about $`B^{}`$ is totally consistent with the approach in paper I, which was based on Herring’s formula. There one takes the wavefunctions as $`(C_{m,d}\pm C_{m,d})/\sqrt{2}`$, where $`C_{m,d}`$ is the wavefunction of a state localized in the left well, and which decays away from that well in both directions. If we equate $`C_{m,d}`$ with $`C_{m,1a}^{}`$ (and, therefore, $`C_{m,d}`$ with $`C_{m,1a}^{\prime \prime }`$) in the central region, then we see that $`B^{}`$, which by Eq. (21) gives the magnitude of the growing part of $`C_m^{}`$ at $`m=m_c^{}`$, is also the order of magnitude of $`C_{m,1a}^{\prime \prime }`$, the decaying part of $`C_m^{\prime \prime }`$ at $`m=m_c^{}`$.
### I What is the origin of the diabolical points?
Lastly, it is extremely instructive to examine the problem when $`\mathrm{\Delta }=0`$, i.e., $`\mathrm{cos}\mathrm{\Lambda }_c=0`$, without necessarily imposing the condition $`\delta =0`$, for this gives insight into what causes the quenching of the tunnel splitting. Taking the imaginary parts of Eqs. (60) and (61) we see directly that we must have $`R^{}=R^{\prime \prime }=0`$, and that
$$Q^{\prime \prime }=\pm B^{}e^{\mathrm{\Gamma }_c+\mathrm{\Gamma }^{\prime \prime }},Q^{}=\pm B^{\prime \prime }e^{\mathrm{\Gamma }_c+\mathrm{\Gamma }^{}}.$$
(75)
Going back to Eq. (41), we see that $`R^{}=0`$ requires either $`ϵ^{}=0`$ or $`A^{}=0`$, and likewise for $`R^{\prime \prime }`$, $`ϵ^{\prime \prime }`$, and $`A^{\prime \prime }`$. If $`\delta 0`$, then both $`ϵ^{}`$ and $`ϵ^{\prime \prime }`$ can not vanish, and the only solution is $`ϵ^{}0`$, $`ϵ^{\prime \prime }=0`$, $`A^{\prime \prime }=Q^{\prime \prime }=B^{}=0`$, (or the one obtained by interchanging single and double primes.) The only non-zero coefficients are $`A^{}`$, $`Q^{}`$, and $`B^{\prime \prime }`$. From Eqs. (40) and (75), we see that only one of these coefficients is independent, which can only be fixed by normalization. Thus, we see that the part $`C_{m,2}^{\prime \prime }`$ proportional to $`B^{\prime \prime }`$ should really be regarded as the extreme right-hand tail of a state localized in the left well. If we denote this state by $`|L`$, and the wavefunction $`m|L`$ by $`C_L(m)`$, then ,
$$C_L(m)=\{\begin{array}{cc}A^{}|v_1(m)|^{1/2}\mathrm{exp}\left(i^mq_1(m^{})𝑑m^{}\right),\hfill & m<m_a^{}\text{,}\hfill \\ A^{}(1)^{\mathrm{}^{}}D_{\mathrm{}^{}}(z),\hfill & mm_a^{}\mathrm{to}mm_b^{}\text{,}\hfill \\ Q^{}|v_1(m)|^{1/2}\mathrm{exp}\left(i_{m_b^{}}^mq_1(m^{})𝑑m^{}\right),\hfill & m_b^{}<m<m_b^{\prime \prime }\text{,}\hfill \\ B^{\prime \prime }|v_1(m)|^{1/2}\mathrm{exp}\left(i_{m_b^{\prime \prime }}^mq_1(m^{})𝑑m^{}\right),\hfill & m_b^{\prime \prime }<m\text{.}\hfill \end{array}$$
(76)
\[Note that in the second line, we wrote $`D_{\mathrm{}^{}}`$, not $`D_\nu ^{}`$, and that in the third line, we have not bothered to write the oscillatory exponential continuation in the region $`m_c^{}<m<m_c^{\prime \prime }`$ correctly—see Eq. (47)—since our aim now is merely to indicate the general structure.\] We can define a right-hand function $`C_R(m)`$ analogously. Indeed it is now clear that the two state Hamiltonian (74) is a truncation of the full Hamiltonian (1) in the $`|L`$, $`|R`$ basis.
The above argument shows that for any $`H_z`$ ($`\delta 0`$), we can tune $`H_x`$ so that $`\mathrm{\Delta }`$ vanishes, at which point, the energy eigenfunction is like $`C_L(m)`$ \[or $`C_R(m)`$\], which is localized in one well and does not “see” the other well at all! In ordinary one dimensional quantum mechanics, this is of course impossible, since a wavefunction like $`C_L(m)`$ which continues decaying with increasing $`m`$ in the classically allowed region of the right hand well has the wrong sign of the curvature in that well. In fact, this argument does not depend on having $`\delta 0`$. If $`\delta =0`$ in addition to $`\mathrm{\Delta }=0`$, $`C_L(m)`$ and $`C_R(m)`$ are independent solutions of Schrödingers equation, as is any linear combination, since they are degenerate.
The above point of view helps elucidate the origin of the quenching more clearly. Indeed, it is better to think about the non-symmetric situation ($`H_z0`$) than the symmetric one ($`H_z=0`$). In continuum problems with a symmetric double well, Herring’s formula gives the splitting as proportional to $`[\psi _d(x)(d\psi _d^{}/dx)]_{x=0}`$, where $`\psi _d(x)`$ is a left-well localized wavefunction . It is tempting to think that the splitting in the spin problem vanishes because the oscillatory envelope in $`C_{m,d}`$ in the central region allows the discrete analog of $`\psi _d(x)`$ or $`\psi _d^{}(x)`$ to vanish at the midpoint. This reasoning is false, as one can see from a close examination of the symmetric case wavefunction, or even more clearly, by looking at the situation when $`\mathrm{\Delta }=0`$ but $`\delta 0`$. The condition $`\mathrm{\Delta }=0`$ can not be reduced to a local property of the wavefunction such as its value or its slope at a particular point. Rather it is the global property that the phase integral $`\mathrm{\Lambda }_c`$ be an odd integer times $`\pi /2`$. From this perspective, the quenching is perhaps better visualized as a manifestation of interfering Feynman trajectories and the Berry phase, even though the value of this phase is more easily found using the DPI method.
## III Application to Fe<sub>8</sub>
Let us now apply our general formalism to Fe<sub>8</sub>. The problem of greatest interest is the location of the diabolical points, and for that we need only solve the conditions $`\delta =\mathrm{\Delta }=0`$. We have already given formulas for the matrix elements and the critical curves in Eqs. (10)–(15). The problem that remains is to use these formulas to find $`\delta `$ and $`\mathrm{\Delta }(\mathrm{}^{},\mathrm{}^{\prime \prime })`$. To keep the problem tractable and obtain answers in closed form, we will assume that $`h_z`$ is small. Specifically, we will assume that the reduced field $`h_z`$ defined above Eq. (13) is formally of order $`1/\overline{J}`$. This enables us to evaluate the turning points and action integrals as expansions in powers of $`h_z`$. Also, it is convenient to carry out all calculations in terms of the reduced variable $`\mu `$.
The first step is to obtain $`\delta `$. For this, we need to analyze the critical curve $`U_0(\mu )`$. Its minima $`\mu _{0\pm }`$ are found to be located at
$$\mu _{0\pm }=\pm \mu _0+\frac{h_zh_x^2}{1h_x^2}+O(h_z^2),$$
(77)
where
$$\mu _0=(1h_x^2)^{1/2}.$$
(78)
The quantities $`E_\pm `$, $`\omega _{0\pm }`$, and $`M_\pm `$ defined through Eq. (23) are given by
$`E_\pm `$ $``$ $`U_0(\mu _{0\pm })=\left[h_x^2\pm 2\mu _0h_z+{\displaystyle \frac{h_z^2h_x^2}{1h_x^2}}\right],`$ (79)
$`\omega _{0\pm }`$ $`=`$ $`{\displaystyle \frac{2\lambda ^{1/2}\mu _0}{\overline{J}}}\left[1\pm {\displaystyle \frac{h_z}{2\mu _0}}\left({\displaystyle \frac{1}{\lambda }}+{\displaystyle \frac{1+2h_x^2}{1h_x^2}}\right)+O(h_z^2)\right],`$ (80)
$`M_\pm `$ $`=`$ $`{\displaystyle \frac{1}{2\lambda h_x^2}}\left[1{\displaystyle \frac{h_z}{\mu _0}}\left({\displaystyle \frac{1}{\lambda }}2\right)+O(h_z^2)\right].`$ (81)
Substituting these results in Eq. (72), we obtain
$$\delta (h_z,\mathrm{}^{},\mathrm{}^{\prime \prime })=4\mu _0h_z+\frac{2\sqrt{\lambda }\mu _0}{\overline{J}}(\mathrm{}^{}\mathrm{}^{\prime \prime })\frac{\sqrt{\lambda }h_z}{\overline{J}}(\mathrm{}^{}+\mathrm{}^{\prime \prime }+1)c_1(h_x)+O(\overline{J}^3),$$
(82)
where
$$c_1(h_x)=\frac{1h_x^2+\lambda (1+2h_x^2)}{\lambda (1h_x^2)}.$$
(83)
The next step is to evaluate $`\mathrm{\Delta }(\mathrm{}^{},\mathrm{}^{\prime \prime })`$, or more precisely $`\mathrm{\Lambda }_c`$, the imaginary part of the phase integral defined in Eq. (59), since that by itself locates the diabolical points. This in turn requires expressions for $`\chi (m)`$ and the points $`m_c^{}`$ and $`m_c^{\prime \prime }`$. To obtain $`\chi (m)`$, we return to the Hamilton-Jacobi equation Eq. (9), write $`q=\kappa +i\chi `$, and separate the equation into its real and imaginary parts. Eliminating $`\kappa (m)`$ from the two equations that result, we obtain a single equation for $`\chi (m)`$, which can be written as
$`4t_2(m)X^2g(m)X+{\displaystyle \frac{t_1^2(m)}{4t_2(m)}}=0;`$ (84)
$`g(m)=w(m)+2t_2(m)E,`$ (85)
where $`X\mathrm{cos}\chi (m)`$. \[We can find the equation obeyed by $`\kappa (m)`$ similarly, and we discover that it is again Eq. (84) with $`X=\mathrm{cosh}\kappa (m)`$.\]
What value of $`E`$ should we use in Eq. (84)? As stated earlier, the value of $`\mathrm{\Delta }`$ is relatively insensitive to small changes in the absolute position of $`E`$. Thus we certainly needn’t incorporate the exponentially small shifts caused by the tunneling itself. Secondly, the tunneling is relevant only when $`\delta \mathrm{\Delta }`$. Thus it suffices to set both $`ϵ^{}`$ and $`ϵ^{\prime \prime }`$ to 0 in this part of the calculation. Then, Eqs. (71), (79), and (80), we get
$`E`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(E_{}+E_++(\mathrm{}^{}+\frac{1}{2})\omega _0+(\mathrm{}^{\prime \prime }+\frac{1}{2})\omega _{0+}\right)`$ (86)
$`=`$ $`h_x^2+{\displaystyle \frac{\sqrt{\lambda }\mu _0}{\overline{J}}}(\mathrm{}^{}+\mathrm{}^{\prime \prime }+1)+O(\overline{J}^2).`$ (87)
Substituting this along with the formulas for $`w(m)`$, $`t_1(m)`$, and $`t_2(m)`$ in Eqs. (84) and (85), yields the following equation for $`X`$:
$`(1\lambda )(1\mu ^2)X^2[1+h_x^2\mu ^2\zeta (\mu )]X+{\displaystyle \frac{h_x^2}{1\lambda }}=0;`$ (88)
$`\zeta (\mu )={\displaystyle \frac{\sqrt{\lambda }}{\overline{J}}}[\mu (\mathrm{}^{\prime \prime }\mathrm{}^{})+\mu _0(\mathrm{}^{}+\mathrm{}^{\prime \prime }+1)].`$ (89)
We have separated out the term $`\zeta (\mu )`$ in Eq. (88) as it is of order $`\overline{J}^1`$ relative to the other terms. We can solve the quadratic equation and expand the result as a power series in $`\zeta `$. Recalling that $`X=\mathrm{cos}^2\chi `$, we obtain
$$\mathrm{cos}\chi (\mu )=\frac{h_x}{\sqrt{(1\lambda )(1\mu ^2)}}\left[1+\frac{\zeta (\mu )}{2(1h_x^2\mu ^2)}+O(\zeta ^2)\right].$$
(90)
By setting $`\mathrm{cos}\chi =1`$, we can determine the points $`\mu _c^{}`$ and $`\mu _c^{\prime \prime }`$, for which it is now more convenient to write $`\mu _{c\pm }`$ instead. If the correction $`\zeta `$ were absent, it is easy to see that these points would be at $`\pm \mu _{c0}`$, where
$$\mu _{c0}=[(1\lambda h_x^2)/(1\lambda )]^{1/2}.$$
(91)
Now, with $`\zeta 0`$, we can find $`\mu _{c\pm }`$ as a series in $`\overline{J}^1`$. Up to leading corrections, we get
$$\mu _{c\pm }=\pm \mu _{c0}\frac{1}{2\sqrt{\lambda }\mu _{c0}\overline{J}}[\mu _{c0}(\mathrm{}^{\prime \prime }\mathrm{}^{})\pm \mu _0(\mathrm{}^{}+\mathrm{}^{\prime \prime }+1)].$$
(92)
The phase integral $`\mathrm{\Lambda }_c`$ is given by
$$\mathrm{\Lambda }_c=\overline{J}_{\mu _c}^{\mu _{c+}}\chi (\mu )𝑑\mu .$$
(93)
We will only evaluate this accurate to terms of order $`\overline{J}^0`$. Let us first consider the corrections entailed by replacing the limits by $`\pm \mu _{c0}`$. It is not difficult to show that without the $`\zeta `$ term in Eq. (90),
$$\chi (\mu )\frac{\sqrt{1\lambda }}{h_x}(\mu _{c0}^2\mu ^2)^{1/2}\mathrm{as}\mu \mu _{c0}.$$
(94)
Inclusion of the $`\zeta `$ correction does not change the square root approach to zero, and can not change the coefficient to leading order in $`\overline{J}^1`$. From Eq. (92), $`\mu _{c\pm }`$ differs from $`\pm \mu _{c0}`$ by terms of order $`\overline{J}^1`$, so ignoring these shifts in the limits of the integral causes an error of order $`\overline{J}^{3/2}`$ in the integral. Hence, we have
$$\mathrm{\Lambda }_c\overline{J}_{\mu _{c0}}^{\mu _{c0}}\chi (\mu )𝑑\mu +O(\overline{J}^{1/2}).$$
(95)
The error is smaller than $`O(\overline{J}^0)`$, and so may be ignored. The remaining integral may be done exactly as in paper I. We write the solution to Eq. (90) by $`\chi =\chi _0+\mathrm{\Delta }\chi `$, where $`\chi _0`$ is the solution when the $`\zeta (\mu )`$ correction is ignored, and $`\mathrm{\Delta }\chi `$ is the $`O(\zeta )`$ correction. The $`\chi _0`$ term can be integrated by parts, and yields
$`\mathrm{\Lambda }_{c0}`$ $`=`$ $`2\overline{J}{\displaystyle _0^{\mu _{c0}}}\chi _0(\mu )𝑑\mu `$ (96)
$`=`$ $`\pi \overline{J}[1h_x(1\lambda )^{1/2}].`$ (97)
The $`\mathrm{\Delta }\chi `$ term yields
$`\mathrm{\Delta }\mathrm{\Lambda }_c`$ $`=`$ $`\overline{J}{\displaystyle _{\mu _{c0}}^{\mu _{c0}}}{\displaystyle \frac{\sqrt{1\mu _{c0}^2}}{\sqrt{\mu _{c0}^2\mu ^2}}}{\displaystyle \frac{\zeta (\mu )}{2(1h_x^2\mu ^2)}}𝑑\mu `$ (98)
$`=`$ $`(\mathrm{}^{}+\mathrm{}^{\prime \prime }+1){\displaystyle \frac{\pi }{2}},`$ (99)
where the result follows after an elementary integration by substitution. Note that the first term in $`\zeta (\mu )`$ integrates to 0 as it is odd. The total result for $`\mathrm{\Lambda }_c`$ is (reverting to unscaled variables)
$`\mathrm{\Lambda }_c`$ $`=`$ $`\mathrm{\Lambda }_{c0}+\mathrm{\Delta }\mathrm{\Lambda }_c`$ (100)
$`=`$ $`{\displaystyle \frac{\pi }{2}}\left[2J(\mathrm{}^{}+\mathrm{}^{\prime \prime })2J{\displaystyle \frac{H_x}{H_c\sqrt{1\lambda }}}\right].`$ (101)
It is now elementary to find the conditions for a diabolical point. The offset $`\delta `$ becomes zero at a certain value of $`h_z`$, and the tunneling amplitude $`\mathrm{\Delta }`$ becomes zero when $`\mathrm{\Lambda }_c`$ is an odd multiple of $`\pi /2`$. Adding arguments $`\mathrm{}^{}`$ and $`\mathrm{}^{\prime \prime }`$ to indicate the level numbers becoming degenerate, the diabolicity conditions are,
$`{\displaystyle \frac{H_z(\mathrm{}^{},\mathrm{}^{\prime \prime })}{H_c}}`$ $`=`$ $`{\displaystyle \frac{\sqrt{\lambda }(\mathrm{}^{\prime \prime }\mathrm{}^{})}{2J}}\left[1+{\displaystyle \frac{\sqrt{\lambda }}{4\mu _0J}}(\mathrm{}^{}+\mathrm{}^{\prime \prime }+1)c_1(h_x)+O(J^2)\right],`$ (102)
$`{\displaystyle \frac{H_x(\mathrm{}^{},\mathrm{}^{\prime \prime })}{H_c}}`$ $`=`$ $`{\displaystyle \frac{\sqrt{1\lambda }}{J}}\left[Jn\frac{1}{2}(\mathrm{}^{}+\mathrm{}^{\prime \prime }+1)\right],n=0,1,\mathrm{},2J(\mathrm{}^{}+\mathrm{}^{\prime \prime }+1).`$ (103)
Note that what appears in these equations is $`J`$, not $`\overline{J}`$. Secondly, the restrictions on the integer $`n`$ in Eq. (103) are obtained, as explained in I, by demanding that $`\mathrm{\Lambda }_c`$ be positive.
We obtain the perfect centered rectangular lattice of diabolical points introduced in Sec. I if we ignore the $`c_1`$ term in Eq. (102), and also the restriction $`H_z/H_c1`$ used to perform the calculation. From the viewpoint of our DPI calculations, this exactness is somewhat mysterious. However, by the duality argument of Sec. I, if formulas (102) and (102) are correct for small $`H_z`$ and large $`H_x`$, then they also yield diabolical points for large $`H_z`$ and small $`H_x`$. Of course if the former set of points corresponds to low lying levels, i.e., small values of $`\mathrm{}^{}`$ and $`\mathrm{}^{\prime \prime }`$, to which our analysis applies, the latter set of points corresponds to rather highly excited levels, to which the analysis does not apply prima facie. It is nevertheless surprising that the formulas should fit together so neatly. It is also somewhat surprising that the experimentally determined Fe<sub>8</sub> diabolical points \[for the cases $`(\mathrm{}^{},\mathrm{}^{\prime \prime })=(0,0)`$, $`(0,1)`$, and $`(0,2)`$\] should lie on a centered rectangular structure so closely, since, as is known, the value of the measured $`H_x`$ period on the $`H_z=0`$ line is almost 50% different from that predicted by Eq. (103). This feature is at present understood only on the basis of numerical diagonalization of the spin Hamiltonian including fourth order terms in $`𝐉`$. An analytic approach to this aspect of the problem remains for the future.
###### Acknowledgements.
This work is supported by the NSF via grant number DMR-9616749. I am indebted to Wolfgang Wernsdorfer and Jacques Villain for useful discussions and correspondence about Fe<sub>8</sub>.
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# Cutoff dependence of the Casimir effect
## Abstract
The problem of calculating the Casimir force on two conducting planes by means of the stress tensor is examined. The evaluation of this quantity is carried out using an explicit regularization procedure which has its origin in the underlying (2+1) dimensional Poincaré invariance of the system. The force between the planes is found to depend on the ratio of two independent cutoff parameters, thereby rendering any prediction for the Casimir effect an explicit function of the particular calculational scheme employed. Similar results are shown to obtain in the case of the conducting sphere.
In 1948 Casimir first predicted that two infinite parallel plates in vacuum would attract each other. This remarkable result has its origin in the zero point energy of the electromagnetic field. While the latter is highly divergent, the change associated with this quantity for specific plate configurations has been found in numerous calculations to be finite as well as cutoff dependent and thus in principle observable. Early work to detect this small effect was characterized by relatively large experimental uncertainties which left the issue in some doubt. More recent efforts have provided quite remarkable data, but are based on a different geometry from that of Casimir. Since a rigorous theoretical calculation has never been carried out for the latter configuration, there remains room for skepticism as to whether the Casimir effect is as well established as is frequently asserted.
The most elementary calculation of the Casimir effect between two parallel conducting planes located at $`z=0`$ and $`z=a`$ employs a mode summation in the framework of a regularization which depends only on the frequency $`\omega _k=[𝐤^2+(\frac{n\pi }{a})^2]^{\frac{1}{2}}`$ where $`n=0,1,2,\mathrm{}`$. Upon combining the result obtained with the corresponding result for for the interval $`azL`$ where $`La`$ is the $`z`$-coordinate of a third conducting plane, a finite cutoff independent result $`F/A=\pi ^2/240a^4`$ is obtained for the Casimir pressure on the plate at $`z=a`$.
A considerably more elegant approach to this problem is that of Brown and Maclay who employ an image method to calculate $`0|T^{\mu \nu }(x)|0`$. Thus they showed that the photon propagator in the presence of conducting planes at $`z=0`$ and $`z=a`$ could be expressed in terms of an infinite sum over the usual (i.e., $`\mathrm{}<z<\mathrm{}`$) photon propagator with the $`z`$-coordinate of each term in the sum displaced by an even multiple of $`a`$. Since the stress tensor for the electromagnetic case is given by
$$T^{\mu \nu }(x)=F^{\mu \alpha }F_\alpha ^\nu \frac{1}{4}g^{\mu \nu }F^{\alpha \beta }F_{\alpha \beta }$$
(1)
where
$$F^{\mu \nu }(x)=^\mu A^\nu (x)^\nu A^\mu (x),$$
it follows that upon taking appropriate derivatives with respect to the propagator arguments $`x`$ and $`x^{}`$ and invoking the limit $`xx^{}`$ a formal expression can be obtained for the vacuum expectation value of the stress tensor. On the basis of covariance arguments together with the divergence and trace free property of $`T^{\mu \nu }(x)`$ it was then found in ref. 5 that
$$0|T^{\mu \nu }(x)|0=(\frac{1}{4}g^{\mu \nu }\widehat{z}^\mu \widehat{z}^\nu )(\frac{1}{2\pi ^2a^4})\underset{n=1}{\overset{\mathrm{}}{}}n^4$$
(2)
where $`\widehat{z}^\mu `$ is the unit vector (0,0,1,0) in the $`z`$-direction normal to the conducting planes.
However, there is some reason to question whether this approach has adequately dealt with the divergences which invariably occur in Casimir calculations. One notes in particular that the result (2) is obtained only after an obviously singular $`n=0`$ term has been dropped from the sum which occurs in that equation. While one can argue as in that such an $`a`$-independent term can be freely omitted since it is merely the usual subtraction of the large $`a`$ result, it is well to note that the entire sum over $`n`$ is required for a demonstration that the propagator satisfies correct boundary conditions at $`z=0,a`$. Moreover, as is shown in this work, an appropriately regularized form of (2) does not necessarily allow a separation into cutoff dependent terms and $`a`$-dependent terms, in contrast with the result found in . Of still greater import is the fact that more general regularizations than those usually considered in this calculation lead to an explicit cutoff dependence of the Casimir stress, a circumstance which would seem to deny its physical significance.
To establish the above claims one reverts from the image approach to one based on expansion of the Green’s function in terms of orthogonal functions . To this end one notes that the free field propagator in the radiation gauge can be written as
$`G^{ij}(𝐱𝐱^{},z,z^{},tt^{})={\displaystyle \underset{n\lambda }{}}{\displaystyle \frac{d𝐤d\omega }{(2\pi )^3}e^{i\omega (tt^{})}}`$ (3)
$`\times {\displaystyle \frac{A_{n\lambda }^i(𝐤,z)A_{n\lambda }^j(𝐤,z^{})}{k^2\omega ^2+(n\pi /a)^2iϵ}}e^{i𝐤(𝐱𝐱^{})}`$ (4)
where $`\lambda =1,2`$ refers to the polarization, and spatial coordinates orthogonal to the $`z`$-direction are denoted by a boldface notation. The eigenfunctions $`A_{n\lambda }^i(𝐤,z)`$ satisfy the equation
$$\left[\frac{^2}{z^2}+(n\pi /a)^2\right]A_{n\lambda }^i(𝐤,z)=0$$
and are given explicitly by
$$A_{n1}^i(𝐤,z)=\frac{\overline{k}_i}{|𝐤|}(\frac{2}{a})^{\frac{1}{2}}sin(n\pi z/a)$$
(5)
and
$$A_{n2}^i(𝐤,z)=\frac{1}{|𝐤|\omega _k}\left(\widehat{𝐳}^i\omega _k^2+\widehat{𝐳}^i\right)(\frac{2}{a})^{\frac{1}{2}}cos(n\pi z/a)$$
(6)
where $`\overline{k}_iϵ^{ij}k_j`$ with $`ϵ^{ij}`$ being the usual alternating symbol. It is important to note that each eigenfunction $`A_{n\lambda }^i(𝐤,z)`$ satisfies the boundary conditions $`\widehat{𝐳}\times 𝐄=\widehat{𝐳}𝐁=0`$ at $`z=0,a`$. This means that it is possible to introduce a regularization such that contributions from large values of $`|𝐤|`$ and/or $`n`$ are reduced without destroying the validity of the boundary conditions. This stands in marked contrast with the image method which has no mechanism for the consistent suppression of the contributions of higher order reflections.
In order to determine the regularization appropriate to this calculation one should ideally make reference to the underlying symmetry. Since the latter consists of the reflection $`zaz`$ and the (2+1) dimensional Poincaré group, it is natural to seek to classify regularization schemes according to representations of the latter. The usual cutoff method for this problem invokes a parameter which damps out the large $`\omega _k`$ contributions, an approach which makes no reference to the underlying Lorentz invariance. A far more appropriate technique is to generalize this to a cutoff based on a vector $`\sigma ^\mu `$ in (2+1) dimensions as well as a scalar cutoff $`\mathrm{\Sigma }`$ which can be used to suppress large values of the (2+1) dimensional invariant $`E^2𝐏^2`$ where $`E`$ and $`𝐏`$ are respectively the energy and momentum operators associated with this (2+1) dimensional subspace. Clearly, the credibility of the Casimir effect requires that the result be independent of the relative importance of these two competing cutoffs.
The calculation proceeds by noting that since the limit $`xx^{}`$ is to be taken symmetrically at some point, it is appropriate to use only the imaginary part of the propagator. An appropriately regularized version of this function can be inferred from Eq.(3) to be
$`I`$ $`mG_{\sigma ,\mathrm{\Sigma }}^{ij}(x,x^{})=\pi {\displaystyle \underset{n\lambda }{}}{\displaystyle \frac{d^3k}{(2\pi )^3}\delta (k^2+(n\pi /a)^2)}`$ (8)
$`\times A_{n\lambda }^i(𝐤,z)A_{n\lambda }^j(𝐤,z^{})e^{ik^\mu (xx^{})_\mu }e^{\sigma _\mu k^\mu ϵ(k^0)}e^{\mathrm{\Sigma }(k^2)^{\frac{1}{2}}}`$
where $`ϵ(k^0)`$ is the alternating function and a summation convention convention has been introduced in the Lorentz invariant subspace. Note that since both $`\sigma ^\mu `$ and $`k^\mu `$ are three vectors in that space, they satisfy the orthogonality conditions $`\widehat{z}^\mu \sigma _\mu =\widehat{z}^\mu k_\mu =0`$ . In addition it is clearly necessary to impose $`\overline{\sigma }^2\sigma ^\mu \sigma _\mu >0`$ and $`\sigma ^0>0`$ in order that this propagator exist. It will subsequently be found that its existence also requires $`\mathrm{\Sigma }<\overline{\sigma }`$.
To proceed one uses the regularization (6) and the form of the stress tensor (1). When used in conjunction with the eigenfunctions (4) and (5) the vacuum expectation value of the regularized stress tensor can be determined. With some effort this is found by straightforward calculation to yield the coordinate independent result
$`0`$ $`|`$ $`T^{\mu \nu }|0={\displaystyle \frac{2\pi }{a}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d^3k}{(2\pi )^3}\delta (k^2+(n\pi /a)^2)e^{\sigma _\mu k^\mu ϵ(k^0)}}`$
$`\times e^{\mathrm{\Sigma }n\pi /a}\left[k^\mu k^\nu +\widehat{z}^\mu \widehat{z}^\nu (n\pi /a)^2\right]`$
which is manifestly both symmetric and traceless. It can be more usefully written as
$`0`$ $`|`$ $`T^{\mu \nu }|0={\displaystyle \frac{1}{a}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}e^{\mathrm{\Sigma }n\pi /a}`$ (10)
$`\times \left({\displaystyle \frac{}{\sigma _\mu }}{\displaystyle \frac{}{\sigma _\nu }}\widehat{z}^\mu \widehat{z}^\nu {\displaystyle \frac{^2}{\sigma ^\alpha \sigma _\alpha }}\right)\mathrm{\Delta }^{n\pi /a}(i\sigma )`$
where $`\mathrm{\Delta }^{n\pi /a}(x)`$ is the (2+1) dimensional function
$$\mathrm{\Delta }^{n\pi /a}(x)=2\pi \frac{d^3k}{(2\pi )^3}e^{ikxϵ(k^0)}\delta (k^2+(n\pi /a)^2)$$
for a particle of mass $`n\pi /a`$. Since this is an $`O(2,1)`$ scalar, $`\mathrm{\Delta }^{n\pi /a}(i\sigma )`$ is a function of only the invariant $`\overline{\sigma }`$ which has the explicit form
$$\mathrm{\Delta }^{n\pi /a}(i\sigma )=\frac{1}{2\pi \overline{\sigma }}e^{\overline{\sigma }n\pi /a}.$$
The insertion of this result into Eq.(7) clearly implies that the sum over $`n`$ exists only for the case that $`\mathrm{\Sigma }<\overline{\sigma }`$ as previously stated. Upon performing the summation over $`n`$ it follows that
$$0|T^{\mu \nu }|0=(\frac{}{\sigma _\mu \sigma _\nu }\widehat{z}^\mu \widehat{z}^\nu \frac{^2}{\sigma ^\alpha \sigma _\alpha })F(\overline{\sigma },\mathrm{\Sigma })$$
where
$$F(\overline{\sigma },\mathrm{\Sigma })=\frac{1}{2\pi a\overline{\sigma }}\frac{1}{1e^{(\mathrm{\Sigma }\overline{\sigma })\pi /a}}.$$
One now performs the usual expansion of the denominator of this expression, discarding terms which give no contribution in the limit of vanishing cutoff, thereby obtaining
$$F(\overline{\sigma },\mathrm{\Sigma })[\frac{1}{2\pi ^2}\frac{1}{\overline{\sigma }}\frac{1}{\overline{\sigma }\mathrm{\Sigma }}+\frac{1}{4\pi a\overline{\sigma }}(1\frac{\mathrm{\Sigma }\pi }{6a})\frac{(\overline{\sigma }\mathrm{\Sigma })^3\pi ^2}{1440\overline{\sigma }a^4}].$$
Upon performing the derivatives and rearranging terms there finally results
$`0`$ $`|`$ $`T^{\mu \nu }|0=[g^{\mu \nu }+3{\displaystyle \frac{\sigma ^\mu \sigma ^\nu }{\overline{\sigma }^2}}\widehat{z}^\mu \widehat{z}^\nu \left]\right\{{\displaystyle \frac{1}{4\pi a\overline{\sigma }^3}}(1{\displaystyle \frac{\mathrm{\Sigma }\pi }{6a}})+`$
$`{\displaystyle \frac{(2\overline{\sigma }\mathrm{\Sigma })(\overline{\sigma }\mathrm{\Sigma })+\frac{2}{3}\overline{\sigma }^2}{2\pi ^2\overline{\sigma }^3(\overline{\sigma }\mathrm{\Sigma })^3}}+{\displaystyle \frac{\pi ^2}{1440a^4}}{\displaystyle \frac{\mathrm{\Sigma }}{\overline{\sigma }}}({\displaystyle \frac{\mathrm{\Sigma }^2}{\overline{\sigma ^2}}}1)\}`$
$`+({\displaystyle \frac{1}{4}}g^{\mu \nu }\widehat{z}^\mu \widehat{z}^\nu )\left[(1{\displaystyle \frac{\mathrm{\Sigma }}{\overline{\sigma }}}){\displaystyle \frac{\pi ^2}{180a^4}}{\displaystyle \frac{4}{3\pi ^2}}{\displaystyle \frac{1}{\overline{\sigma }}}{\displaystyle \frac{1}{(\overline{\sigma }\mathrm{\Sigma })^3}}\right].`$
If (following ) one subtracts the $`a\mathrm{}`$ result, this reduces to the more tractable form
$`0`$ $`|`$ $`\overline{T}^{\mu \nu }|0=\left[g^{\mu \nu }+3{\displaystyle \frac{\sigma ^\mu \sigma ^\nu }{\overline{\sigma }^2}}\widehat{z}^\mu \widehat{z}^\nu \right]`$
$`\times \left\{{\displaystyle \frac{1}{4\pi a\overline{\sigma }^3}}\left(1{\displaystyle \frac{\mathrm{\Sigma }\pi }{6a}}\right)+{\displaystyle \frac{\pi ^2}{1440a^4}}{\displaystyle \frac{\mathrm{\Sigma }}{\overline{\sigma }}}\left({\displaystyle \frac{\mathrm{\Sigma }^2}{\overline{\sigma }^2}}1\right)\right\}`$
$`+\left({\displaystyle \frac{1}{4}}g^{\mu \nu }\widehat{z}^\mu \widehat{z}^\nu \right)\left(1{\displaystyle \frac{\mathrm{\Sigma }}{\overline{\sigma }}}\right){\displaystyle \frac{\pi ^2}{180a^4}}`$
where an overbar notation has been used to denote this subtraction. It is noteworthy that even this removal of the large $`a`$ result does not lead to regularization independent results, a fact which has been remarked upon earlier.
Of particular interest to Casimir calculations are the stress components $`0|\overline{T}^{33}|0`$ and the energy density per unit area $`a0|\overline{T}^{00}|0`$ which are given by
$$0|\overline{T}^{33}|0=\frac{\pi ^2}{240a^4}(1\frac{\mathrm{\Sigma }}{\overline{\sigma }})$$
(11)
and
$``$ $`=`$ $`{\displaystyle \frac{\pi ^2}{720a^3}}\{(1{\displaystyle \frac{\mathrm{\Sigma }}{\overline{\sigma }}}){\displaystyle \frac{3\sigma _0^2\overline{\sigma }^2}{2\overline{\sigma }^2}}`$ (14)
$`\times [{\displaystyle \frac{\mathrm{\Sigma }}{\overline{\sigma }}}({\displaystyle \frac{\mathrm{\Sigma }^2}{\overline{\sigma }^2}}1){\displaystyle \frac{30\mathrm{\Sigma }a^2}{\pi ^2\overline{\sigma }^3}}]\}+`$
$`(a\mathrm{independent}\mathrm{terms})`$
respectively. It is striking that each of these terms retains a significant dependence on the cutoff details. In addition the usual relation assumed (as in ) to hold between $``$ and the stress components, namely
$$0|\overline{T}^{33}|0=\frac{}{a},$$
(15)
is manifestly contradicted by Eqs.(8) and (9) in agreement with results found earlier in the context of the Casimir energy of a sphere . It is significant that the relation (10) asserts an equality of the vacuum stress $`0|\overline{T}^{33}|0`$ which transforms under $`O(2,1)`$ as a scalar while the right hand side transforms as the $`\mu =\nu =0`$ component of a symmetric tensor under this group. Finally, note should be made of the fact that Eq.(9) predicts an additional Casimir force proportional to the divergent indeterminate form $`\mathrm{\Sigma }/a^2\overline{\sigma }^3`$.
To reinforce the conclusions reached here in the case of parallel plates it is useful to consider also the case of the conducting sphere, the only other geometry in three dimensions which has proved amenable to exact calculation . This case was first solved by Boyer and subsequently verified by a number of authors \[12-15\]. Following reasoning similar to that of the parallel plate case note is made of the fact that the unbroken symmetry in this case consists of time translation and rotational invariance. Thus the natural cutoff parameters in this problem should refer to the energy and angular momentum. The former is the standard one and is well known to give cutoff independent results. It will be the goal here to examine the situation which occurs when a combination of these two is considered.
This is most economically achieved by reference to which provides a useful separation of the Casimir energy into a finite part and one which requires regularization. Thus one writes for a sphere of radius $`a`$
$$E_c=E_{fin}+E_\sigma $$
where $`E_c`$, $`E_{fin}`$, and $`E_\sigma `$ are respectively the total, the regularization independent part, and the formally divergent parts of the Casimir energy. The quantity $`E_\sigma `$ is given by
$`E_\sigma `$ $`=`$ $`{\displaystyle \frac{1}{4\pi a}}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}Ree^{i\varphi }{\displaystyle _0^{\mathrm{}}}𝑑yexp(i\nu \sigma ye^{i\varphi })y{\displaystyle \frac{d}{dy}}`$
$`\times (1+y^2e^{2i\varphi })^3`$
where $`\nu =l+\frac{1}{2}`$, $`0<\varphi <\frac{\pi }{2}`$, and $`\sigma `$ is a dimensionless cutoff used to suppress the high frequency modes. Upon choosing a secondary cutoff of the form $`e^{\mathrm{\Sigma }\nu }`$ it is readily found that $`\mathrm{\Delta }E_\sigma `$ (the change induced in $`E_\sigma `$ in the limit of small cutoff) is given by
$$\mathrm{\Delta }E_\sigma =\frac{3\mathrm{\Sigma }}{2\pi a\sigma ^2}_0^{\mathrm{}}𝑑y\frac{y^2}{(1+y^2)^4}\frac{1}{y^2+(\mathrm{\Sigma }^2/\sigma ^2)}.$$
This is evaluated to yield
$$\mathrm{\Delta }E_\sigma =\frac{3}{64a\sigma }\frac{\mathrm{\Sigma }}{(\mathrm{\Sigma }+\sigma )^4}[\mathrm{\Sigma }^2+4\sigma \mathrm{\Sigma }+5\sigma ^2],$$
a result which displays yet again the cutoff dependence of the Casimir effect for a more general choice of regularization. It may be noted that aside from confirming the vanishing of $`\mathrm{\Delta }E_\sigma `$ for $`\mathrm{\Sigma }=0`$, this result shows that $`\mathrm{\Delta }E_\sigma `$ diverges for $`\sigma 0`$ with all intermediate values obtained for finite $`\mathrm{\Sigma }/\sigma `$.
In this work it has been shown that the Casimir effect is, prevailing opinion notwithstanding, highly dependent on the particular form of regularization employed for the extraction of the force. As remarked earlier as well as in ref. the recent experiments which have seemed to many to provide the long awaited precision verification of this highly subtle effect are not based upon rigorous mathematical calculation. While the parallel plate Casimir experiment is fraught with difficulties beyond the ken of this author, it would seem that the successful completion of such experiments would be invaluable for purposes of setting to rest some of the issues which have been raised in this work.
Finally, it would be remiss not to mention in some way the very extensive work on the calculation of Casimir forces using the technique of zeta function regularization . Historically, the successes of the Casimir approach in dealing with the parallel plate geometry and the sphere were obtained using conventional field theoretical subtraction procedures. Specifically, it was noted that only changes relative to the vacuum could be considered observable and it was therefore totally consistent to perform subtractions relative to the $`a\mathrm{}`$ vacuum. However, this step did not succeed in allowing one to obtain finite and observable results in more general applications. Eventually it was realized, however, that the application of zeta function regularization to such problems could yield finite results for some fairly general cases while at the same time agreeing with those obtained in the few instances in which more conventional subtractions could be successfully applied. This work makes no claim to having established any inconsistencies in the derivation of finite results for the Casimir effect when those efforts are based on the twin axioms of vacuum energy and zeta function regularization. Rather, the calculations presented here establish that the Casimir effect is generally cutoff dependent and hence incapable of being reliably determined whenever such calculations are performed using conventional (i.e., physically plausible) subtraction procedures.
###### Acknowledgements.
This work is supported in part by the U.S. Department of Energy Grant No.DE-FG02-91ER40685.
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# 1 Introduction
## 1 Introduction
The two-brane Randall-Sundrum scenario provides an appealing way to generate the electroweak gauge hierarchy as a consequence of spacetime geometry. The basic idea is to start with five dimensional anti de Sitter (AdS) space, take the region between two slices parallel to the AdS horizon, and add a 3-brane along each slice. By tuning the brane tensions, the resulting configuration can be made stable against gravitational collapse.
In this model, the ratio of the weak to the Planck scale is determined by the distance between the two branes. The distance is fixed by the expectation value of a modulus field, called the radion. The usual hierarchy problem now appears in a new guise: What fixes the radion vev, and what protects the vev against large radiative corrections?
In a recent paper, Goldberger and Wise proposed a way to stabilize the radion using five dimensional bulk matter. Supersymmetry provides another possibility. In this paper we will take some first steps in that direction, and show how to supersymmetrize the minimal Randall-Sundrum scenario.
In what follows we will use coordinates in which the five dimensional background metric takes the following form,
$$ds^2=e^{2\sigma (\varphi )}\eta _{mn}dx^mdx^n+r^2d\varphi ^2.$$
(1)
The coordinate $`x^5=r\varphi `$ parametrizes an orbifold $`S^1/Z_2`$, where the circle $`S^1`$ has radius $`r`$ and the orbifold identification is $`\varphi \varphi `$. For fixed $`\varphi `$, the coordinates $`x^m`$ ($`m=0,1,2,3`$) span flat Minkowski space, with metric $`\eta _{mn}=\mathrm{diag}(1,1,1,1)`$. We choose to work on the orbifold covering space, so we take $`\pi <\varphi \pi `$.
For the gravitational part of our action, we follow Randall and Sundrum and take the action to be the sum of bulk plus brane pieces,
$$S=S_{\mathrm{bulk}}+S_{\mathrm{brane}}.$$
(2)
The bulk action is that of pure five dimensional AdS gravity, while $`S_{\mathrm{brane}}`$ arises from the presence of two opposite tension branes.
The gravitational bulk action is given by
$$S_{\mathrm{bulk}}=\frac{\mathrm{\Lambda }}{\kappa ^2}d^5xe\left[\frac{1}{2}R+6\mathrm{\Lambda }^2\right],$$
(3)
where $`\kappa `$ is related to the four dimensional Planck constant, $`e=dete_M^A`$, and $`e_M^A`$ is the five dimensional fünfbein.<sup>1</sup><sup>1</sup>1We adopt the convention that capital letters run over the set $`\{0,1,2,3,5\}`$ and lower-case letters run from 0 to 3. Tangent space indices are taken from the beginning of the alphabet; coordinate indices are from the middle. We follow the conventions of . In this expression, $`\mathrm{\Lambda }`$ is the bulk cosmological constant and $`R`$ is the five dimensional Ricci scalar,
$$R=e_A{}_{}{}^{M}e_{B}^{}{}_{}{}^{N}R_{MN}^{}{}_{}{}^{AB}.$$
(4)
The Riemann curvature $`R_{MN}^{AB}`$ is built from the spin connection according to the following conventions,
$$R_{MN}{}_{}{}^{AB}=_M\omega _N{}_{}{}^{AB}_N\omega _M{}_{}{}^{AB}\omega _M{}_{}{}^{AC}\omega _{NC}^{}{}_{}{}^{B}+\omega _N{}_{}{}^{AC}\omega _{MC}^{}{}_{}{}^{B}.$$
(5)
The brane action serves as a source for the bulk gravitational fields. It arises from the 3-branes located at the orbifold points $`\varphi =0,\pi `$. For the case at hand, the brane action is simply
$$S_{\mathrm{brane}}=6\frac{\mathrm{\Lambda }^2}{r\kappa ^2}d^5x\widehat{e}\left[\delta (\varphi )\delta (\varphi \pi )\right],$$
(6)
where $`\widehat{e}=dete_m^a`$, and the $`e_m^a`$ are the components of the five dimensional fünfbein, restricted to the appropriate brane.
From this action it is not hard to show that the metric (1), with
$$\sigma (\varphi )=r\mathrm{\Lambda }|\varphi |,$$
(7)
is a solution to the five dimensional Einstein equations,
$$R_{MN}\frac{1}{2}g_{MN}R=6g_{MN}\mathrm{\Lambda }^2+6g_{mn}\delta _M^m\delta _N^n\left(\frac{\mathrm{\Lambda }}{r}\right)\left(\frac{\widehat{e}}{e}\right)[\delta (\varphi )\delta (\varphi \pi )].$$
(8)
Away from the branes, the bulk metric is just that of five dimensional AdS space, with cosmological constant $`\mathrm{\Lambda }`$. On the branes, the four dimensional metric is flat. As shown in , the effective theory of the gravitational zero modes is just ordinary four dimensional Einstein gravity, with a vanishing cosmological constant. The effective four dimensional squared Planck mass is $`\kappa _{\mathrm{eff}}^2=\kappa ^2(1e^{2\pi r\mathrm{\Lambda }})`$.
## 2 Supersymmetric Bulk
In what follows we will supersymmetrize the action (2). We start with the bulk action (3). Its supersymmetric extension can be found from the five dimensional supersymmetric AdS action ,
$`S_{\mathrm{bulk}}`$ $`=`$ $`\mathrm{\Lambda }{\displaystyle }d^5xe[{\displaystyle \frac{1}{2\kappa ^2}}R+\mathrm{i}ϵ^{MNOPQ}\overline{\mathrm{\Psi }}_M\mathrm{\Sigma }_{NO}D_P\mathrm{\Psi }_Q{\displaystyle \frac{1}{4}}F_{MN}F^{MN}`$ (9)
$`3\mathrm{\Lambda }\overline{\mathrm{\Psi }}_M\mathrm{\Sigma }^{MN}\mathrm{\Psi }_N+6{\displaystyle \frac{\mathrm{\Lambda }^2}{\kappa ^2}}\mathrm{i}\kappa \sqrt{{\displaystyle \frac{3}{2}}}{\displaystyle \frac{1}{2}}F_{MN}\overline{\mathrm{\Psi }}^M\mathrm{\Psi }^N`$
$`\kappa {\displaystyle \frac{1}{6\sqrt{6}}}ϵ^{MNOPQ}F_{MN}F_{OP}B_Q+\mathrm{i}\kappa \sqrt{{\displaystyle \frac{3}{2}}}{\displaystyle \frac{1}{4}}ϵ^{MNOPQ}F_{MN}\overline{\mathrm{\Psi }}_O\mathrm{\Gamma }_P\mathrm{\Psi }_Q`$
$`\kappa \mathrm{\Lambda }\sqrt{{\displaystyle \frac{3}{2}}}ϵ^{MNOPQ}\overline{\mathrm{\Psi }}_M\mathrm{\Sigma }_{NO}\mathrm{\Psi }_PB_Q+\text{ four-Fermi terms }],`$
where the $`ϵ`$ tensor is defined to have tangent-space indices, and $`ϵ^{01235}=1`$. This action contains the physical fields associated with the supergravity multiplet in five dimensions: the fünfbein $`e_M^A`$, the gravitino $`\mathrm{\Psi }_M`$, and a vector field $`B_M`$. The covariant derivative $`D_M\mathrm{\Psi }_N=_M\mathrm{\Psi }_N+\frac{1}{2}\mathrm{\Sigma }^{AB}\omega _{MAB}\mathrm{\Psi }_N`$ and the matrix $`\mathrm{\Sigma }^{AB}=\frac{1}{4}(\mathrm{\Gamma }^A\mathrm{\Gamma }^B\mathrm{\Gamma }^B\mathrm{\Gamma }^A)`$.
This action is invariant under the following supersymmetry transformations,
$`\delta e_M^A`$ $`=`$ $`\mathrm{i}\kappa (\overline{\eta }\mathrm{\Gamma }^A\mathrm{\Psi }_M\overline{\psi }_M\mathrm{\Gamma }^A\eta )`$
$`\delta B_M`$ $`=`$ $`\mathrm{i}\sqrt{{\displaystyle \frac{3}{2}}}(\overline{\eta }\psi _M\overline{\psi }_M\eta )`$
$`\delta \psi _M`$ $`=`$ $`{\displaystyle \frac{2}{\kappa }}D_M\eta +\mathrm{i}{\displaystyle \frac{\mathrm{\Lambda }}{\kappa }}\mathrm{\Gamma }_M\eta \mathrm{i}\sqrt{6}\mathrm{\Lambda }B_M\eta `$ (10)
$`\sqrt{{\displaystyle \frac{2}{3}}}(\mathrm{\Gamma }^NF_{NM}{\displaystyle \frac{1}{4}}ϵ_{MNOPQ}F^{NO}\mathrm{\Sigma }^{PQ})\eta `$
$`+\text{ three-Fermi terms.}`$
Since we work in the orbifold covering space, the spacetime manifold has no boundary, and we can freely integrate by parts. We use the 1.5 order formalism, so the spin connection obeys its own equation of motion and does not need to be varied.
For the case at hand, we must define the action of the orbifold symmetry on the AdS fields. We start by writing the five dimensional spinors in a four dimensional language, where
$$\mathrm{\Psi }_M\left(\begin{array}{c}\psi _{M\alpha }^1\\ \\ \overline{\psi }_{2M}^{\dot{\alpha }}\end{array}\right)$$
(11)
and
$$\mathrm{\Gamma }^a\left(\begin{array}{cc}0& \sigma _{\alpha \dot{\alpha }}^a\\ \\ \overline{\sigma }^{a\dot{\alpha }\alpha }& 0\end{array}\right)\mathrm{\Gamma }^5\left(\begin{array}{cc}\mathrm{i}& 0\\ \\ 0& \mathrm{i}\end{array}\right).$$
(12)
The fields $`\psi _M^i`$ (for $`i=1,2`$) are two-component Weyl spinors, in the notation of . We then define $`\psi _M^\pm =\frac{1}{\sqrt{2}}(\psi _M^1\pm \psi _M^2)`$, and likewise for $`\eta ^\pm `$.
In terms of these fields, the bulk supersymmetry transformations can be written in the following form,
$`\delta e_M^a`$ $`=`$ $`\mathrm{i}\kappa (\eta ^+\sigma ^a\overline{\psi }_M^++\eta ^{}\sigma ^a\overline{\psi }_M^{})+h.c.`$
$`\delta e_M^{\widehat{5}}`$ $`=`$ $`\kappa (\eta ^+\psi _M^{}\eta ^{}\psi _M^+)+h.c.`$
$`\delta B_M`$ $`=`$ $`\mathrm{i}\sqrt{{\displaystyle \frac{3}{2}}}(\eta ^+\psi _M^{}\eta ^{}\psi _M^+)+h.c.`$
$`\delta \psi _m^\pm `$ $`=`$ $`{\displaystyle \frac{2}{\kappa }}D_m\eta ^\pm {\displaystyle \frac{\mathrm{i}}{\kappa }}\omega _{ma\widehat{5}}\sigma ^a\overline{\eta }^{}\pm \mathrm{i}{\displaystyle \frac{\mathrm{\Lambda }}{\kappa }}e_m{}_{}{}^{a}\sigma _{a}^{}\overline{\eta }^\pm +e_{m\widehat{5}}{\displaystyle \frac{\mathrm{\Lambda }}{\kappa }}\eta ^{}`$
$`\mathrm{i}\sqrt{6}\mathrm{\Lambda }B_m\eta ^{}\sqrt{{\displaystyle \frac{2}{3}}}(e_a{}_{}{}^{N}F_{Nm}^{}\sigma ^a\overline{\eta }^{}\mathrm{i}e_{\widehat{5}}{}_{}{}^{N}F_{Nm}^{}\eta ^\pm `$
$`{\displaystyle \frac{1}{4}}ϵ_{ABCde}e_m{}_{}{}^{A}e_{}^{BN}e^{CO}F_{NO}\sigma ^{de}\eta ^\pm \pm {\displaystyle \frac{\mathrm{i}}{4}}ϵ_{abcd}e_m{}_{}{}^{a}e_{}^{bN}e^{cO}F_{NO}\sigma ^d\overline{\eta }^{})`$
$`\delta \psi _5^\pm `$ $`=`$ $`{\displaystyle \frac{2}{\kappa }}D_5\eta ^\pm {\displaystyle \frac{\mathrm{i}}{\kappa }}\omega _{5a\widehat{5}}\sigma ^a\overline{\eta }^{}+e_{5\widehat{5}}{\displaystyle \frac{\mathrm{\Lambda }}{\kappa }}\eta ^{}\pm \mathrm{i}e_5{}_{}{}^{a}{\displaystyle \frac{\mathrm{\Lambda }}{\kappa }}\sigma _a\overline{\eta }^\pm `$ (13)
$`\mathrm{i}\sqrt{6}\mathrm{\Lambda }B_5\eta ^{}+\sqrt{{\displaystyle \frac{2}{3}}}(e_a{}_{}{}^{n}F_{5n}^{}\sigma ^a\overline{\eta }^{}\mathrm{i}e_{\widehat{5}}{}_{}{}^{n}F_{5n}^{}\eta ^\pm `$
$`+{\displaystyle \frac{1}{4}}ϵ_{ABCde}e_5{}_{}{}^{A}e_{}^{BN}e^{CO}F_{NO}\sigma ^{de}\eta ^\pm \pm {\displaystyle \frac{\mathrm{i}}{4}}ϵ_{abcd}e_5{}_{}{}^{a}e_{}^{bN}e^{cO}F_{NO}\sigma ^d\overline{\eta }^{}).`$
In these expressions, all fields depend on the five dimensional coordinates. The symbol $`\widehat{5}`$ denotes the fifth tangent space index, and all covariant derivatives contain the spin connection $`\omega _{Mab}`$. Here and hereafter, we ignore all three- and four-Fermi terms.
From these transformations it is not hard to find a consistent set of $`Z_2`$ parity assignments under the orbifold transformation $`\varphi \varphi `$. The assignments must leave the action and transformation laws invariant under the $`Z_2`$ symmetry. We assign even parity to
$$e_m{}_{}{}^{a},e_{5\widehat{5}},B_5,\psi _m^+,\psi _5^{},\eta ^+$$
and odd parity to
$$e_5{}_{}{}^{a},e_m{}_{}{}^{\widehat{5}},B_m,\psi _m^{},\psi _5^+,\eta ^{}.$$
The bulk supergravity action is invariant under $`N=1`$ supersymmetry in five dimensions. The branes break all but one one four dimensional supersymmetry. To find its form, we shall study the supersymmetry transformations in the orbifold background, where $`e_{5\widehat{5}}=1`$, $`e_m{}_{}{}^{a}=e^{\sigma (\varphi )}\delta _m^a`$, and all other fields equal zero. This configuration satisfies the gravitational equations of motion when $`\sigma (\varphi )=r\mathrm{\Lambda }|\varphi |`$. Note that this background is consistent with the orbifold symmetry.
In the orbifold background, the supersymmetry variations of the bosonic fields are obviously zero. The variations of the fermions are a little trickier. In this background, the spin connection evaluates to
$$\omega _{mAM}\mathrm{\Sigma }^{AM}=\mathrm{sgn}(\varphi )\mathrm{\Lambda }\mathrm{\Gamma }_m\mathrm{\Gamma }^{\widehat{5}},$$
(14)
with all other components zero. The supersymmetry variations of the fermions reduce to the following form,
$`\delta \psi _m^\pm `$ $`=`$ $`{\displaystyle \frac{2}{\kappa }}_m\eta ^\pm \mathrm{i}\mathrm{sgn}(\varphi ){\displaystyle \frac{\mathrm{\Lambda }}{\kappa }}\sigma _m\overline{\eta }^{}\pm \mathrm{i}{\displaystyle \frac{\mathrm{\Lambda }}{\kappa }}\sigma _m\overline{\eta }^\pm `$
$`\delta \psi _5^\pm `$ $`=`$ $`{\displaystyle \frac{2}{\kappa }}_5\eta ^\pm +{\displaystyle \frac{\mathrm{\Lambda }}{\kappa }}\eta ^{}.`$ (15)
The unbroken supersymmetries are found by setting these variations to zero. The resulting Killing equations can then be solved for the Killing spinors $`\eta ^\pm `$. The solution that reduces to a flat-space supersymmetry in four dimensions is simply
$$\eta ^+=\frac{1}{\sqrt{2}}e^{\sigma (\varphi )/2}\eta (x),\eta ^{}=\frac{1}{\sqrt{2}}e^{\sigma (\varphi )/2}\mathrm{sgn}(\varphi )\eta (x),$$
(16)
where $`\mathrm{sgn}(\varphi )`$ is the step function,<sup>2</sup><sup>2</sup>2The distribution $`\mathrm{sgn}(\varphi )`$ obeys the following properties:
$$_ϵ^ϵ𝑑\varphi \mathrm{sgn}(\varphi )=0,_ϵ^ϵ𝑑\varphi \mathrm{sgn}^2(\varphi )=2ϵ,$$
when integrated against smooth functions, and
$$_ϵ^ϵ𝑑\varphi \mathrm{sgn}(\varphi )\delta (\varphi )=0,_ϵ^ϵ𝑑\varphi \mathrm{sgn}^2(\varphi )\delta (\varphi )=\frac{1}{3},$$
when integrated against $`\delta (\varphi )`$. The last relation ensures that
$$_ϵ^ϵ𝑑\varphi \frac{d}{d\varphi }\mathrm{sgn}^3(\varphi )=2.$$
We thank Jan Conrad for a discussion on this point. which evaluates to $`(1,0,1)`$, depending on the sign of $`\varphi `$. In this expression, the spinor $`\eta `$ contains four Grassmann components and is a function of $`x^0,\mathrm{},x^3`$, but not $`x^5`$. We shall see that it describes the one unbroken supersymmetry of the Randall-Sundrum scenario.
It is not hard to check that the spinors (16) are a solution to the Killing equations, for constant $`\eta `$, except for delta-function singularities at the orbifold points $`\varphi =0,\pi `$. These singularities are very important. They motivate us to change the $`\psi _5^{}`$ supersymmetry transformation so that the spinors (16) are Killing spinors everywhere. We take
$`\delta \psi _5^{}`$ $`=`$ $`{\displaystyle \frac{2}{\kappa }}D_5\eta ^{}+{\displaystyle \frac{\mathrm{i}}{\kappa }}\omega _{5a\widehat{5}}\sigma ^a\overline{\eta }^++e_{5\widehat{5}}{\displaystyle \frac{\mathrm{\Lambda }}{\kappa }}\eta ^+\mathrm{i}e_5{}_{}{}^{a}{\displaystyle \frac{\mathrm{\Lambda }}{\kappa }}\sigma _a\overline{\eta }^{}`$ (17)
$`\mathrm{i}\sqrt{6}\mathrm{\Lambda }B_5\eta ^++\sqrt{{\displaystyle \frac{2}{3}}}(e_a{}_{}{}^{n}F_{5n}^{}\sigma ^a\overline{\eta }^+\mathrm{i}e_{\widehat{5}}{}_{}{}^{n}F_{5n}^{}\eta ^{}`$
$`+{\displaystyle \frac{1}{4}}ϵ_{ABCde}e_5{}_{}{}^{A}e_{}^{BN}e^{CO}F_{NO}\sigma ^{de}\eta ^{}{\displaystyle \frac{\mathrm{i}}{4}}ϵ_{abcd}e_5{}_{}{}^{a}e_{}^{bN}e^{cO}F_{NO}\sigma ^d\overline{\eta }^+)`$
$`{\displaystyle \frac{4}{r\kappa }}[\delta (\varphi )\delta (\varphi \pi )]\eta ^+.`$
In the orbifold background, this reduces to
$$\delta \psi _5^{}=\frac{2}{\kappa }_5\eta ^{}+\frac{\mathrm{\Lambda }}{\kappa }\eta ^+\frac{4}{r\kappa }[\delta (\varphi )\delta (\varphi \pi )]\eta ^+.$$
(18)
The spinors (16) satisfy the modified Killing equations, for constant $`\eta `$, even at the orbifold points $`\varphi =0,\pi `$. Furthermore, the supersymmetry transformations still close into the $`N=1`$ supersymmetry algebra.
## 3 Supersymmetric Brane
In the previous section, we changed the gravitino supersymmetry transformations so that the Killing spinors satisfy the Killing equations at every point in $`\varphi `$. Because of this, the bulk action is no longer invariant under the supersymmetry transformations (17). In this section we will find a brane action whose variation precisely cancels that of the bulk.
We first compute the variation of the bulk action. Comparing (13) with (17), we see that the bulk variation vanishes except on the branes. Therefore, to compute the variation, we need to project the bulk fields onto the branes. For even fields, this is easy: The brane fields are just the bulk fields evaluated at the appropriate value of $`\varphi `$. For odd fields, the situation is more subtle: The brane fields must obey jump conditions across the delta function singularities and these conditions are determined by the brane action.
In what follows we will present the brane action and verify that it restores the supersymmetry of the bulk-plus-brane system. We assert that the brane action is simply
$$S_{\mathrm{brane}}=\frac{\mathrm{\Lambda }}{r\kappa ^2}d^5x\widehat{e}(3\mathrm{\Lambda }+2\kappa ^2\psi _m^+\sigma ^{mn}\psi _n^+)\left[\delta (\varphi )\delta (\varphi \pi )\right]+h.c.$$
(19)
where the fields $`e_m^a`$ and $`\psi _m^+`$ are projections of the corresponding five dimensional fields.
Given this brane action, it is easy to compute the jump conditions. From the equations of motion for $`e_m^a`$ and $`\psi _m^+`$, we find
$$[\omega _{ma5}]=\pm \mathrm{\hspace{0.17em}2}\mathrm{\Lambda }e_{ma},[\psi _m^{}]=\pm \mathrm{\hspace{0.17em}2}\psi _m^+,$$
(20)
where the square brackets denote the discontinuity across the singularity, and the $`\pm `$ applies to the brane at $`\varphi =0`$ and $`\pi `$, respectively. A consistent solution is given by
$$\omega _{ma5}=\mathrm{sgn}(\varphi )\mathrm{\Lambda }e_{ma},\psi _m^{}=\mathrm{sgn}(\varphi )\psi _m^+,$$
(21)
in the neighborhood of the branes. All other odd fields vanish on the branes.
Now that we have the solutions to the jump conditions, we are free to compute the variation of the bulk action. A small calculation gives
$`\delta S_{\mathrm{bulk}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }}{r\kappa }}{\displaystyle }d^5xee^{\widehat{5}5}[(\mathrm{\hspace{0.17em}8}\eta ^+\sigma ^{mn}D_m\psi _n^+\mathrm{i}\kappa \sqrt{6}F^{\widehat{5}m}\eta ^+\psi _m^+`$ (22)
$`+6\mathrm{i}\mathrm{\Lambda }(1\mathrm{sgn}^2(\varphi ))\eta ^+\sigma ^m\overline{\psi }_m^+)[\delta (\varphi )\delta (\varphi \pi )]]+h.c.`$
where $`\eta ^+`$ is the spinor (16). In what follows we will show that the variation of the brane action precisely cancels this term.
The supersymmetry variation of the brane action is not hard to find. The supersymmetry transformations are those of the bulk fields, as projected on the branes, subject to the jump conditions (20). From (13) and (21), we compute
$`\delta e_m^a`$ $`=`$ $`\mathrm{i}\kappa (1+\mathrm{sgn}^2(\varphi ))\eta ^+\sigma ^a\overline{\psi }_m^++h.c.`$
$`\delta \psi _m^+`$ $`=`$ $`{\displaystyle \frac{2}{\kappa }}D_m\eta ^++\mathrm{i}{\displaystyle \frac{\mathrm{\Lambda }}{\kappa }}(1\mathrm{sgn}^2(\varphi ))e_m{}_{}{}^{a}\sigma _{a}^{}\overline{\eta }^+`$ (23)
$`+\mathrm{i}\sqrt{{\displaystyle \frac{2}{3}}}F_{\widehat{5}m}\eta ^++\mathrm{i}\sqrt{{\displaystyle \frac{2}{3}}}F^{\widehat{5}n}\sigma _{mn}\eta ^+.`$
As above, $`\eta ^+`$ is given by (16). In all fields, the coordinate $`\varphi `$ is evaluated at $`\varphi =0`$ or $`\pi `$, depending on the location of the brane. Substituting (23) into (19), we find
$`\delta S_{\mathrm{brane}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }}{r\kappa }}{\displaystyle }d^5x\widehat{e}[(\mathrm{\hspace{0.17em}8}\eta ^+\sigma ^{mn}D_m\psi _n^+\mathrm{i}\kappa \sqrt{6}F^{\widehat{5}m}\eta ^+\psi _m^+`$ (24)
$`+12\mathrm{i}\mathrm{\Lambda }\mathrm{sgn}^2(\varphi )\eta ^+\sigma ^m\overline{\psi }_m^+)[\delta (\varphi )\delta (\varphi \pi )]]+h.c.`$
The variation of the brane action, (24), cancels the variation of the bulk action, (22), because $`e=e_{5\widehat{5}}\widehat{e}`$ and $`\mathrm{sgn}^2(\varphi )=1/3`$ when integrated against a delta function. This proves that the full bulk-plus-brane Randall-Sundrum action is invariant under the four dimensional supersymmetry parametrized by the Killing spinor $`\eta `$ in Eq. (16).
## 4 Minimal Effective Action
We will now derive the effective four dimensional action for the supergravity zero modes. We will see that it is nothing but the usual on-shell four dimensional flat-space supergravity action.
The zero modes of the four dimensional theory must satisfy the massless equations of motion in four dimensions. For the vierbein, the zero mode was given by Randall and Sundrum :
$$e_M{}_{}{}^{A}=\left(\begin{array}{cc}1& 0\\ 0& e^{\sigma (\varphi )}\overline{e}_m{}_{}{}^{a}(x)\end{array}\right),$$
(25)
where $`\sigma (\varphi )=r\mathrm{\Lambda }|\varphi |`$ and the vierbein $`\overline{e}_m^a`$ is a function of $`x^0,\mathrm{},x^3`$, but not $`x^5`$. The five dimensional Einstein equations, with brane sources, reduce to the usual four dimensional source-free Einstein equations for the vierbein $`\overline{e}_m^a`$.
The gravitino zero modes can be found in a similar way. One starts with the five dimensional gravitino equations of motion,
$`_5\psi _m^++{\displaystyle \frac{3}{2}}\mathrm{\Lambda }\psi _m^{}\mathrm{sgn}(\varphi )\mathrm{\Lambda }\psi _m^+`$ $`=`$ $`0`$
$`_5\psi _m^{}+{\displaystyle \frac{3}{2}}\mathrm{\Lambda }\psi _m^+\mathrm{sgn}(\varphi )\mathrm{\Lambda }\psi _m^{}`$ $`=`$ $`{\displaystyle \frac{2}{r}}\left[\delta (\varphi )\delta (\varphi \pi )\right]\psi _m^+,`$ (26)
and assumes the following ansatz,
$$\psi _m^+=\frac{1}{\sqrt{2}}\left(\frac{\kappa _{\mathrm{eff}}}{\kappa }\right)e^{\sigma (\varphi )/2}\psi _m(x),\psi _m^{}=\frac{1}{\sqrt{2}}\left(\frac{\kappa _{\mathrm{eff}}}{\kappa }\right)e^{\sigma (\varphi )/2}\mathrm{sgn}(\varphi )\psi _m(x).$$
(27)
Substituting (27) into (26), one recovers the usual four dimensional equations of motion for the gravitino field $`\psi _m`$.
In what follows, we will derive the effective four dimensional action for the supergravity zero mode fields. We start by setting all other fields to zero. This truncation is consistent with the supersymmetry transformations (10). We then substitute the zero-mode expressions into the supersymmetric bulk-plus-brane action and integrate over the coordinate $`x^5`$. We use the fact that
$$R=e^{2\sigma }\overline{R}+20\mathrm{\Lambda }^216\frac{\mathrm{\Lambda }}{r}[\delta (\varphi )\delta (\varphi \pi )]$$
(28)
and
$$\omega _{mAB}\mathrm{\Sigma }^{AB}=\mathrm{sgn}(\varphi )\mathrm{\Lambda }\mathrm{\Gamma }_m\mathrm{\Gamma }^{\widehat{5}}+\overline{\omega }_{mab}\sigma ^{ab}$$
(29)
to find
$$S_{\mathrm{eff}}=d^4x\overline{e}\left[\frac{1}{2\kappa _{\mathrm{eff}}^2}\overline{R}+ϵ^{mnpq}\overline{\psi }_m\overline{\sigma }_nD_p\psi _q\right],$$
(30)
up to four-Fermi terms. This is nothing but the on-shell action for flat-space $`N=1`$ supergravity in four dimensions.
The supersymmetry transformation laws can be found in a similar way. We start with the supersymmetry transformation parameters $`\eta ^+`$ and $`\eta ^{}`$ as above, in (16). We then substitute the zero mode expressions into the supersymmetry transformations (10). All $`x^5`$ dependent terms cancel, leaving
$`\delta e_m^a`$ $`=`$ $`\mathrm{i}\kappa _{\mathrm{eff}}\eta \sigma ^a\overline{\psi }_m+h.c.`$
$`\delta \psi _m`$ $`=`$ $`{\displaystyle \frac{2}{\kappa _{\mathrm{eff}}}}D_m\eta .`$ (31)
These are nothing but the transformations of $`N=1`$ supergravity in four dimensions (up to three-Fermi terms), with an effective four dimensional squared Planck mass, $`\kappa _{\mathrm{eff}}^2=\kappa ^2(1e^{2\pi r\mathrm{\Lambda }})`$.
## 5 Summary and Outlook
In this paper we supersymmetrized the minimal Randall-Sundrum scenario. We found the supersymmetric bulk-plus-brane action in five dimensions, as well as the corresponding supersymmetry transformations. We solved for the Killing spinor that describes the unbroken $`N=1`$ supersymmetry of the four dimensional effective theory. We derived the supergravitational zero modes, and showed that the low energy effective theory reduces to ordinary $`N=1`$ supergravity in four dimensions.
This work represents a first step towards a deeper understanding of supersymmetry in the context of warped compactifications. To study stability, one would like, of course, to include the radion multiplet, which reduces to $`N=1`$ matter in four dimensions. For phenomenology, one would also like to add supersymmetric matter on the branes and in the bulk. Work along all these lines is in progress.
We are pleased to acknowledge helpful conversations with Jan Conrad, Dan Freedman, Erich Poppitz, Raman Sundrum and Max Zucker. This work was supported by the National Science Foundation, grant NSF-PHY-9404057, and the Department of Energy, contract DE-FG03-84ER40168.
Note added: On the same day this paper was submitted to the archive, a similar paper was posted by Gherghetta and Pomarol . This paper used an $`x^5`$-dependent bulk gravitino mass to supersymmetrize the two-brane Randall-Sundrum scenario. The resulting construction can be interpreted as a truncation of a more fundamental theory with matter in the bulk. We did not take this approach because our goal was to supersymmetrize the purely gravitational case. For more on the difficulties of constructing brane-like solutions in matter-coupled five dimensional supergravity, see and .
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# I. Territory covered by 𝑁 random walkers on deterministic fractals. The Sierpinski gasket
## I INTRODUCTION
Random walk theory is a branch of statistical physics with many applications . Problems related to a single random walker have traditionally been the subject of thorough study, but their generalizations to $`N>1`$ random walkers have attracted much less attention, although there are some, generally very recent, exceptions . These multiparticle diffusion problems are characterized by the impossibility of being analyzed in terms of the single random walker theory, i.e., they can not be solved through simple averaging over the properties of a single random walker, even in the noninteracting case. The recent development of experimental techniques allowing the observation of events caused by single particles of an ensemble should give additional encouragement to the study of these of multiparticle diffusion problems.
The problem that is the subject of this paper, namely, the evaluation of the average number $`S_N(t)`$ of distinct sites visited (or territory covered) by $`N`$ random walkers up to time $`t`$, all moving from the same starting site, is a clear example of a diffusion problem that can not be solved, or even approximated, from the solution for $`N=1`$, $`S_1(t)`$. Even for independent random walkers, the overlap of the regions explored by different walkers prohibits a decomposition of $`S_N(t)`$ into single-particle contributions. The origin of the problem of evaluating $`S_N(t)`$ is usually traced back to the Proceedings of the Second Berkeley Symposium on Mathematical Statistics and Probability where the case $`N=1`$ was posed by Dvoretzky and Erdös . Since then, the quantity $`S_1(t)`$ has been studied in detail and is discussed in general references . For fractal substrates (in particular, for the two-dimensional Sierpinski gasket and the two-dimensional percolation cluster at criticality) this problem was first studied by Rammal et al. . More recently, Larralde et al. and Havlin et al. studied the problem of evaluating $`S_N(t)`$ when $`N1`$ noninteracting random walkers diffuse in Euclidean and fractal media, respectively. For fractal lattices with spectral dimension $`d_s=2d_f/d_w<2`$, it was argued that $`S_N(t)t^d_{\mathrm{}}`$ for $`t<tt_\times \mathrm{ln}N`$ and $`S_N(t)t^{d_s/2}\mathrm{ln}^{d_f/u}\left(N\right)`$ for $`t_\times t`$, where $`d_{\mathrm{}}=d_f/d_{\text{min}}`$ is the chemical dimension (or topological distance dimension), $`d_f`$ is the fractal dimension, $`d_{\text{min}}`$ is the fractal dimension of the shortest path on the fractal, $`u=d_w/(d_w1)`$ and $`d_w`$ is the diffusion exponent (or fractal dimension of the random walk) . However, Dräger and Klafter using scaling arguments have recently proposed that
$$S_N(t)t^{d_s/2}(\mathrm{ln}N)^{d_{\mathrm{}}/v}$$
(1)
for $`t_\times t`$, where $`v=d_w^{\mathrm{}}/(d_w^{\mathrm{}}1)`$ and $`d_w^{\mathrm{}}=d_w/d_{\text{min}}`$ is the chemical-diffusion exponent . Of course, the two predictions agree for the media considered, such as Sierpinski gaskets, in the present paper for which $`d_{\text{min}}=1`$. Therefore, in what follows, we will write $`d_f`$, $`u`$ and $`d_w`$ instead of $`d_{\mathrm{}}`$, $`v`$ and $`d_w^{\mathrm{}}`$, respectively. Fractals with $`d_{\text{min}}1`$ will be discussed in the following paper .
As stated above, two time regimes are observed in $`S_N(t)`$: an extremely short-time regime or regime I and a long-time regime or regime II separated by the crossover time $`t_\times \mathrm{ln}N`$. A further long-time regime, or regime III, is observed in Euclidean lattices when the movement of the independent walkers are very far from each other so that their trails (almost) never overlap and $`S_N(t)NS_1(t)`$ . In the one-dimensional lattice and fractal lattices with $`d_s2`$, the trails of the random walkers partially overlap at all times and regime III is never reached. Such is the case in this paper where we are concerned only with fractals in which $`d_s<2`$. The transition from regime I to regime II is easy to understand. In regime I, we have so many particles at every site that all nearest neighbors of the already visited sites are necessarily reached by some walker at the next time step. The minimum path length between two sites on a fractal is the chemical distance, $`\mathrm{}`$, and it is clear that after a time $`t`$ the visited zone is an hypersphere of chemical radius $`\mathrm{}=t`$ and, consequently, $`S_N(t)t^d_{\mathrm{}}`$ according to the definition of the chemical exponent . This behavior lasts until the average number of particles per site is of the order of unity, then overlapping is only partial and a new regime is established. If $`z`$ is the coordination number of the lattice (or the average coordination number in the case of stochastic fractals), the number of random walkers at every site decreases as $`N/z^t`$ in the very short-time regime because the particles are distributed between $`z`$ sites at every step. The breaking of the overlapping regime thus takes place when $`N/z^t1`$ or, equivalently, $`t_\times \mathrm{ln}N`$.
Regime II is far more interesting and difficult to analyze than regime I due to the nontrivial interplay of the walkers in their exploration of the lattice, which leads to a more complex dependence of $`S_N(t)`$ on time $`t`$ and number of particles $`N`$. In some recent work , we have shown that for independent random walks on Euclidean lattices there exist important asymptotic corrections to the main term that can not be ignored even for very large number of particles as these corrections decay only logarithmically. We will see that this also holds for the two-dimensional Sierpinski gasket.
It should be noticed that, except for the Sierpinski gasket in two dimensions when $`N=1`$ , there has never been any discussion about $`S_N(t)`$ focused on deterministic fractals, whether theoretically or numerically. Certainly, a dependence on $`t`$ and $`N`$ of the main asymptotic term of $`S_N(t)`$ for large $`N`$ has been proposed \[see Eq. (1)\], but nothing is known about the value of its amplitude or prefactor and on the relevance (if any) of the other (corrective) asymptotic terms. In this paper we present a procedure for obtaining, for a certain class of fractals, the complete asymptotic series expansion of $`S_N(t)`$ when $`N1`$. The procedure gives the main asymptotic term in full, and determines the functional form of the corrective terms, which we calculate explicitly up to second order. The fractals that we consider in this paper have to satisfy two conditions. First, the number of sites (or volume) $`V(r)`$ of the fractal inside a hypersphere of radius $`r`$ should be given by
$$V(r)=V_0r^{d_f}$$
(2)
where $`V_0`$ is a constant characteristic of the fractal substrate; and, second, the probability $`\mathrm{\Gamma }_t(𝐫)`$ that a site $`𝐫`$ has not been visited by a single random walker by time $`t`$ should decay for $`\xi =|𝐫|/Rr/R1`$ as
$$\mathrm{\Gamma }_t(𝐫)1A\xi ^{\mu u}\mathrm{exp}(c\xi ^u)\left(1+h_1\xi ^u+\mathrm{}\right),$$
(3)
where $`R^2=2Dt^{2/d_w}`$ is the mean-square displacement of a single random walker and $`D`$ is the diffusion constant. It should be clear at this point that the above two conditions can only be satisfied approximately: first, because $`V_0`$ is not strictly constant (it exhibits log-periodic oscillations of small amplitude, see Sec. III); second, because $`\mathrm{\Gamma }_t(𝐫)`$ does not solely depend on the distance $`r`$ but also (in general) on the actual location $`𝐫`$ on the lattice; and, third, because $`\mathrm{\Gamma }_t(𝐫)`$ is not continuos (this fact can be clearly seen in figure 1 of Ref. ) so that Eq. (3) can only be an approximation to the true distribution. The fluctuations in $`S_N(t)`$ associated to these effects are thus not included in our theoretical discussion. However, their importance can be gauged by resorting to simulation. For the two-dimensional Sierpinski gasket, we found that these fluctuations are indeed relevant and that they can be explained to a large extent as a consequence of the log-periodic oscillations of $`V_0`$. Finally, there is another difficulty regarding the value of $`\mathrm{\Gamma }_t(𝐫)`$: while its dominant term $`\mathrm{exp}(c\xi ^u)`$ is reasonably well established, the value (and even the form) of its subdominant factors $`A\xi ^{\mu u}`$, $`h_1\xi ^u`$, etc. is unknown (although an educated guess can be made; see Sec. II). This means that we can be fairly sure of the value of the main term of $`S_N(t)`$ because, as we will show, it depends only on the dominant term of $`\mathrm{\Gamma }_t(𝐫)`$. However, the true value of the corrective terms of $`S_N(t)`$ is more uncertain as they also depend on the subdominant factors of $`\mathrm{\Gamma }_t(𝐫)`$. Nevertheless, we will see that reasonable choices of values for these subdominant factors lead to significant improvements in the estimate of $`S_N(t)`$.
The simulation results for the two-dimensional Sierpinski gasket will reveal the great importance of the corrective terms (for example, they account for some thirty per cent of the total value of $`S_N(t)`$ even for $`N`$ as large as $`10^6`$). This is not strange because, as we will show in Sec. II, they decay only logarithmically with $`N`$. An important consequence that we will address in the following paper is that the corrective terms must be taken into account in the analyses based on “collapsing” the numerical data to determine the exponents in the main term of $`S_N(t)`$.
The paper is organized as follows. The asymptotic evaluation of $`S_N(t)`$ on fractal lattices is presented in Sec. II. The mathematical techniques involved in the calculation are very similar to those corresponding to the Euclidean case and we will only outline the main steps. Details may be found in Ref. . In Sec. III, we compare the asymptotic expansion of $`S_N(t)`$ with simulation results obtained on the Sierpinski gasket. We end in Sec. IV with some remarks on the quality of the asymptotic approximation.
## II TERRITORY COVERED BY $`N`$ RANDOM WALKERS ON A DETERMINISTIC FRACTAL SUBSTRATE
We will define the multiparticle survival probability $`\mathrm{\Gamma }_N(t,𝐫)`$ as the probability that site $`𝐫`$ has not been visited by time $`t`$ by any of the $`N`$ random walkers that start from the origin site $`𝐫=\mathrm{𝟎}`$ at time $`t=0`$. From this definition, the following relationship between the average number of distinct sites visited, $`S_N(t)`$, and the survival probability can be established :
$$S_N(t)=\underset{𝐫}{}\left\{1\mathrm{\Gamma }_N(t,𝐫)\right\},$$
(4)
where the sum is over all the sites in the fractal lattice. For independent random walkers, we have $`\mathrm{\Gamma }_N(t,𝐫)=\left[\mathrm{\Gamma }_t(𝐫)\right]^N`$, where $`\mathrm{\Gamma }_t(𝐫)\mathrm{\Gamma }_1(t,𝐫)`$ is the one-particle survival probability. Since we are interested in the behavior of $`S_N(t)`$ after a large number of steps (thus beyond of the very-short-time regime I), we replace $`\mathrm{\Gamma }_t(𝐫)`$ and $`S_N(t)`$ by their continuum approximations:
$`S_N(t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}\left[1\mathrm{\Gamma }_t^N(r)\right]𝑑V(r)`$ (5)
$`=`$ $`V_0(2D)^{d_f/2}t^{d_f/d_w}J_N(d_f),`$ (6)
where we have assumed that the fractal volume $`V(r)`$ (i.e., the number of lattice sites) of the hypersphere with radius $`r`$ is given by Eq. (2) and that, essentially, $`\mathrm{\Gamma }_t(𝐫)`$ only depends on the distance $`r`$. In Eq. (5), $`J_N(d_f)`$ is
$$J_N(d_f)=N_0^{\mathrm{}}\left[\mathrm{\Gamma }_t(\xi )\right]^{N1}\frac{d\mathrm{\Gamma }_t(\xi )}{d\xi }\xi ^{d_f}𝑑\xi .$$
(7)
In order to evaluate this integral for $`N1`$ it suffices to know $`\mathrm{\Gamma }_t(r)`$ for large $`\xi `$. For Euclidean lattices
$$\mathrm{\Gamma }_t(r)\mathrm{\Gamma }(\xi )1A\xi ^{\mu u}\mathrm{exp}(c\xi ^u)\left(1+\underset{n=1}{\overset{\mathrm{}}{}}h_n\xi ^{nu}\right)$$
(8)
when $`\xi 1`$ . In this paper we shall consider only those fractals, such as the Sierpinski gaskets, for which this equation approximately \[see discussion below Eq.(3)\] holds. Indeed, we can show that Eq. (8) is reasonable for those lattices ($`d`$-dimensional Sierpinski fractals, Given-Mandelbrot curve, one-dimensional lattice, hierarchical diamond lattice,…) where the renormalization procedure implemented in Refs. can be set up. The argument is as follows. Let $`𝐫_n^{(i)}`$, $`n=1,\mathrm{},z`$, with $`r^{(i)}=|𝐫_n^{(i)}|`$, be the position of the $`z`$ nearest neighbors of the site at $`𝐫=\mathrm{𝟎}`$ in the fractal lattice decimated $`i`$ times (see Fig. 1), and let $`h(t,r^{(i)})`$ be the probability that, in the time interval \[0,t\], a single diffusing particle that starts at $`𝐫=\mathrm{𝟎}`$ is absorbed by any of the traps located at its $`z`$ nearest neighbors placed at the sites $`𝐫_n^{(i)}`$. For large values of $`\xi r^{(i)}/R`$, i.e., for relatively short times, the event “the random walker arrives for the first time at site $`𝐫_n^{(i)}`$” and the event “the random walker arrives for the first time at site $`𝐫_m^{(i)}`$” are (almost) independent so that $`h(t,r^{(i)})`$ is (almost) the sum of the probability \[given by $`1\mathrm{\Gamma }_t(r^{(i)})`$\] of each of the $`z`$ individual events, i.e., $`h(t,r^{(i)})z\left[1\mathrm{\Gamma }_t(r^{(i)})\right]`$. Of course, the two events are not fully independent because the random walker could first arrive at site $`𝐫_m^{(i)}`$ after passing by the site $`𝐫_n^{(i)}`$. However, this is very unlikely because $`𝐫_m^{(i)}`$ and $`𝐫_n^{(i)}`$ are separated by distances of order of $`r^{(i)}`$ so that the fraction of random walkers that, after arriving at $`𝐫_n^{(i)}`$, travel to $`𝐫_m^{(i)}`$ in the short time interval \[0,t\] is of the order $`\mathrm{exp}(\xi ^u)`$, with $`\xi r^{(i)}/R1`$, because the propagator or Green’s function is of this order for large values of $`\xi `$ . Therefore, it is reasonable to assume that $`1\mathrm{\Gamma }_t(r^{(i)})=h(t,r^{(i)})/z\{1+𝒪(\mathrm{exp}(\xi ^u)]\}`$ for large $`\xi `$. But $`1h(t,r^{(i)})`$ has the form of the right-hand side of Eq. (8) with $`u=d_w/(d_w1)`$, at least for the fractals that we are considering , so that Eq. (8) for $`𝐫=𝐫^{(i)}`$ follows. In this discussion we have assumed, in order for the renormalization analysis that leads to $`\mathrm{\Gamma }_t\left(r^{(i)}\right)`$ to work, that the traps were placed at those sites (such as $`A_g`$,$`B_g`$,…in Fig. 1) that become the nearest neighbors of the starting site after several decimations. But, in the calculation of $`S_N(t)`$, the function $`\mathrm{\Gamma }_t(r)`$ is required for every pair of origin and destination sites in the lattice. At this point, we shall assume that the survival probability for $`𝐫𝐫^{(i)}`$ and $`𝐫=𝐫^{(i)}`$ are very similar, i.e., we assume that $`\mathrm{\Gamma }_t(r^{(i)})\mathrm{\Gamma }_\tau (r)=\mathrm{\Gamma }_\tau (\xi )`$ when $`r/\tau ^{1/d_w}=r^{(i)}/t^{1/d_w}=(2D)^{1/2}\xi `$ for large $`\xi `$, with $`\mathrm{\Gamma }_t(\xi )`$ given by Eq. (8). To the best of our knowledge this problem has not been studied and will require a specific and detailed simulation analysis that is not the object of the present work. Nevertheless, previous simulations on the Sierpinski gasket of other statistical quantities closely related to $`\mathrm{\Gamma }_t(r)`$, such as the propagator (or Green’s function) , enable us to affirm with confidence that the parameters $`c`$ and $`u`$ remain unchanged over the whole Sierpinski lattice whereas the subdominant ones, $`A`$, $`\mu `$, $`h_1`$,…, do not. In the following, we use $`c=0.981`$ and $`u=d_w/(d_w1)1.756`$ for the two-dimensional Sierpinski gasket (values of $`c`$ for other fractals can be found in Ref. ) but we must bear in mind that the actual values of $`A`$, $`\mu `$ and $`h_n`$, $`n=1,2,\mathrm{}`$ will likely differ from those obtained by renormalization techniques (namely, $`A=0.61`$, $`\mu =1/2`$, $`h_1=0.56`$) as these latter correspond to the special placing of the traps and origins. We must also point out that, as shown below, since the parameters $`h_n`$, $`n=1,2,3\mathrm{}`$, only contribute to the second and higher order series terms of $`S_N(t)`$ and since the real value of even the first-order asymptotic term is uncertain, the values of these parameters will not be considered here.
The evaluation of $`J_N(d_f)`$ now proceeds in analogy with the analysis for the Euclidean lattices and we shall only quote the final result for $`S_N(t)`$ in the fractal case:
$$S_N(t)\widehat{S}_N(t)(1\mathrm{\Delta })$$
(9)
with
$`\widehat{S}_N(t)`$ $`=`$ $`V_0\left(2D\right)^{d_f/2}t^{d_s/2}\left({\displaystyle \frac{\mathrm{ln}N}{c}}\right)^{d_f/u}`$ (10)
$`\mathrm{\Delta }`$ $`=`$ $`\beta {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{ln}^nN{\displaystyle \underset{m=0}{\overset{n}{}}}s_m^{(n)}\mathrm{ln}^m\mathrm{ln}N,`$ (11)
and where, up to second order ($`n=2`$),
$`s_0^{(1)}`$ $`=`$ $`\omega `$ (12)
$`s_1^{(1)}`$ $`=`$ $`\mu `$ (13)
$`s_0^{(2)}`$ $`=`$ $`(\beta 1)\left({\displaystyle \frac{\pi ^2}{12}}+{\displaystyle \frac{\omega ^2}{2}}\right)(ch_1\mu \omega )`$ (14)
$`s_1^{(2)}`$ $`=`$ $`\mu ^2+(\beta 1)\mu \omega `$ (15)
$`s_2^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\beta 1)\mu ^2`$ (16)
and $`\omega =\gamma +\mathrm{ln}A+\mu \mathrm{ln}c`$, $`\gamma 0.577215`$ is the Euler constant, and $`\beta =d_f/u=d_f(11/d_w)`$.
## III NUMERICAL RESULTS FOR THE SIERPINSKI GASKET
To check the reliability of the analysis presented in the previous section we carried out simulations of the number of distinct sites visited by $`N`$ random walkers on a two-dimensional Sierpinski lattice with $`g=11`$ generations. This means that if we take the length of any side of the smallest triangles (the zeroth decimated triangles) as the unit length, then the length of the sides of the triangle that inscribes the lattice (the $`g`$th decimated triangle) is $`2^g`$. Two different cases are analyzed: (I) random walkers are initially placed upon the center of the base of the main triangle which inscribes the lattice (point O in Fig. 1), and (II) the common starting site is randomly selected. Qualitatively and quantitatively, the results are different in case I and case II. The structure of the lattice gives rise to oscillations superimposed on the general trend of $`S_N(t)`$ in case I. This structure is smeared out in case II by the double average over experiments and over starting sites, so that $`S_N(t)`$ is now a smooth function.
We shall use the top vertex of the main triangle as the origin of the lattice and the orthogonal basis $`\{𝐞_v,𝐞_h\}`$ shown in Fig. 2, because it is convenient to divide the lattice into horizontal sets of sites with a fixed $`v`$ coordinate. The positions of the random walkers are updated at every time step by a random selection of the destination sites among the neighbors of the site occupied by the walker. In order to do that we have to take into account that the sites of the Sierpinski lattice may be classified into three types according to the relative positions of their neighbors: type $`R`$, $`L`$ and $`C`$ as shown in Fig. 2. Since the Sierpinski lattice has many holes on every scale, it is not efficient to identify the lattice sites and the random walkers positions by two coordinates $`\{v,h\}`$ because there are many coordinate pairs corresponding to no actual lattice site. Much memory can be saved if the sites are numbered from top to bottom and from left to right, so that the top vertex of the main triangle is site number $`1`$ and the right vertex is site number $`3(3^g+1)/2`$.
### A Case I
First, we will discuss the simulation results for the territory covered by $`N`$ random walkers placed initially at site O in Fig. 1. In order to compare with the zeroth- and first-order asymptotic expression, cf. Eq. (9), we must know the values of $`V_0`$, $`D`$, $`c`$, $`u`$, $`A`$ and $`\mu `$. In Fig. 3(a) we have plotted the fractal volume of a circle of radius $`r`$, $`V(r)`$, centered upon the privileged site O. The observed log-periodic structure is a consequence of the empty and filled triangular areas that repeat periodically as $`r`$ is increased, but the general trend is well represented by a term of the form $`V_0r^{d_f}`$ with $`V_03.0\pm 0.1`$. In Fig. 3(b), in which the quotient $`V(r)/r^{d_f}`$ is plotted versus $`\mathrm{log}_2r`$, one clearly sees the log-periodic oscillations of $`V_0`$. As our theory assumes a constant value for $`V_0`$, we take the average value over the last period (from maximum to maximum), $`V_0=2.93`$, as a reasonable criterion for comparison with the simulation results. In order to find the diffusion coefficient $`D`$ of a random walker starting at O, we performed $`10^6`$ simulations up to $`t=400`$. The linear numerical fit between $`t=50`$ and $`t=400`$ gives $`d_w2.32`$ (the exact value is $`d_w=\mathrm{ln}5/\mathrm{ln}22.322`$) and $`2D1.05`$. Numerical fits using other time intervals (excluding short times, of course) lead to similar values, and we take $`2D1.05\pm 0.02`$ as a reliable estimate. For the parameters $`c`$ and $`u`$ we take the values (see Sec. II) $`0.981`$ and $`d_w/(d_w1)=1.756`$, respectively. As discussed in Sec. II, the values of $`\mu `$ and $`A`$ are much less certain and we will use here two pairs of values: those obtained by renormalization, i.e., $`\mu =1/2`$ and $`A=0.61`$, and these same values increased by one, i.e., $`\mu =3/2`$ and $`A=1.61`$. Of course, this last pair of parameters are arbitrary (other values could also be used) and are mainly given to show the relevance of the corrective terms.
Simulation results for $`S_N(t)`$ until $`t=1000`$ are shown in Fig. 4 for $`N=10^3`$ and $`N=10^6`$. Overall good agreement is obtained in the comparison with the theoretical prediction of Eq. (9), especially when the values $`\mu =3/2`$ and $`A=1.61`$ are used. This last finding should only be taken as the manifestation of the importance of the corrective terms which can lead to such dramatic changes and improvements in $`S_N(t)`$ after modifying some of the subdominant factors of the survival probability. Of course, additional independent study will be necessary to check the form of $`\mathrm{\Gamma }_t(r)`$ given by Eq. (8) and to find out whether the values for $`\mu `$ and $`A`$ considered here are good estimates of the real values . Nevertheless, it should be noted that the decrease in $`1\mathrm{\Gamma }_t(r)`$ when averaging over the whole lattice with respect to its renormalization value, as is implied by the corresponding increment of $`\mu `$ ($`\mu =1/2\mu =3/2`$), is analogous to the decrease of the propagator when this same averaging is performed. Given that the two statistical quantities (propagator and survival probability) are closely related, one is inclined to accept that, at least, the proposed increment in the value of $`\mu `$ captures the right tendency. The subsequent improvement in the prediction of $`S_N(t)`$ supports this supposition.
The theoretical expression was not able to give a perfect account of the log-periodic oscillations superimposed on the general trend of $`S_N(t)`$ shown in Fig. 4. Notice the self-affinity of the numerical $`S_N(t)`$ plot (the analytical lines are obviously self-affine): if the segment of horizontal axis of Fig. 4 between $`t=0`$ and $`t=200`$ is expanded by a factor $`5`$ and the corresponding segment of the vertical axis is expanded by a factor $`3=5^H`$ ($`H=d_f/d_w=d_s/2`$), we get a figure indistinguishable from Fig. 4. The origin of these scaling factors is clear: the enlargement of the sides of the Sierpinski gasket by a factor $`2`$ implies that its fractal volume increases by a factor $`3`$ and the time that a random walker takes to traverse it increases by $`5`$ .
The structure of $`S_N(t)`$ is more clearly perceived in Fig. 5 where the quotient between the theoretical prediction and the simulation results is plotted. It is remarkable how relatively poor is the performance of the zeroth-order approximation (or main asymptotic term) in predicting the value of $`S_N(t)`$: it accounts for hardly eighty per cent of $`S_N(t)`$ for values of $`N`$ as large as $`10^6`$. However, the inclusion of the first corrective asymptotic term (especially for some suitable selections of the subdominant parameters $`A`$ and $`\mu `$) leads to a noticeable improvement. The log-periodic structure is observed both for $`N=10^3`$ and $`N=10^6`$ but in the latter case this structure is richer and strikingly similar to that of $`V(r)`$ as shown in Fig. 3(b). We attribute this fact to a better mapping of the lattice structure as more and more random walkers are involved in the exploration. We have plotted the solid line in Fig. 4 with the aim of showing to what extent the oscillatory behavior of $`S_N(t)`$ as shown in Fig. 5 can be interpreted as a consequence of the oscillatory behavior of $`V_0`$ shown in Fig. 3(b). The line is generated in the same way as the dotted line, i.e., by means of the first-order approximation of Eq. (9) with $`A=1.61`$ and $`\mu =3/2`$, but, instead of using the averaged value $`V_0=2.93`$ (as in the dotted line), we use the actual oscillatory value of $`V_0`$ taken from the last oscillation (from maximum to maximum) shown in Fig. 3(b). The way in which the solid line runs alongside the simulation results supports the view that the log-oscillatory behavior of $`S_N(t)`$ mainly comes from the log-oscillatory behavior of $`V_0`$.
### B Case II
We also study the effect on $`S_N(t)`$ of choosing other lattice sites, besides the point O, as starting sites for the random walkers. To this end we performed simulations where all the $`N`$ random walkers start on a site randomly selected within the shaded area of Fig. 1 in order to avoid the finite size effects. As expected, the fractal volume $`V(r)`$ and the average number of distinct sites visited $`S_N(t)`$ are smooth functions in this case. An estimate of $`V_0`$ by numerical fitting gives $`V_03.6`$. The analysis of the simulation results ($`10^4`$ runs for $`10^3`$ randomly selected starting sites) for the mean-square displacement of a single random walker is compatible with $`2D0.8`$ when the fit is carried out inside the time interval ($`t=50,t=400`$). Simulation results for $`S_N(t)`$ (five runs for $`10^3`$ randomly selected starting sites) until $`t=200`$ for $`N=1024`$ are shown in Fig. 6. They are compared with the theoretical prediction of the zeroth- and first-order approximations of Eq. (9) for $`\mu =3/2`$ and $`A=1.61`$, and with the corresponding simulation results when the origin is at O (case II). Again, we find a relatively poor performance of the zeroth-order approximation, as well as substantial improvement when the fisrt-order approximation is used, although there is still room for further enhancement. It should be noted that the performance of the two asymptotic approximations is completely analogous to that obtained for Euclidean lattices . In these Euclidean media we found that the second-order asymptotic approximation gives rise to a significant improvement in the estimate of $`S_N(t)`$ even for relatively small values of $`N`$. It thus seems natural to conjecture that the same will occur for the Sierpinski gasket, although definitive confirmation of this guess must wait until reliable values for $`A`$, $`\mu `$ and $`h_1`$ are calculated.
Finally, in Fig. 7 we show the dependence on $`N`$ of $`S_N(t)`$ for case II. We have plotted two first-order asymptotic curves: for the first curve we take the usual values $`\mu =3/2`$ and $`A=1.61`$, and the new values $`\mu =1.75`$ and $`A=1.75`$ are used for the second curve. Again, one sees the great importance of the asymptotic corrective terms as they substantially improve the zeroth-order (main term) asymptotic prediction. We have used the new pair of parameters simply as another example to illustrate the gross effect of the subdominant factors of the survival probability $`\mathrm{\Gamma }_t(r)`$ on the theoretical prediction of $`S_N(t)`$. The excellent agreement reached with $`\mu =1.75`$ and $`A=1.75`$ should not, however, be considered as an indication that they are the correct subdominant parameters of the survival probability .
Notice that the simulation results and theoretical predictions for $`S_N(t)`$ differ very little from case I to case II despite $`V_0`$ and $`D`$ being clearly different in the two situations. The reason for this coincidence is that the quantity $`V_0(2D)^{d_f/2}`$ that appears in the amplitude of the main asymptotic term of $`S_N(t)`$ is almost invariant with respect to translations of the origin site: its value was approximately $`3.05`$ when the origin is placed at O, and $`3.02`$ when the origin sites were randomly selected.
## IV SUMMMARY
We addressed the problem of calculating the average number $`S_N(t)`$ of distinct sites visited (or territory covered) by $`N`$ independent random walkers that diffuse on deterministic fractal lattices. The validity of the main result of this paper, Eq. (9), is based upon two conditions: (a) the volume (number of sites) of the fractal substrate grows as $`V(r)=V_0r^{d_f}`$ and (b) the asymptotic form of the survival probability, $`\mathrm{\Gamma }_t(r)`$, for large $`r`$ and small $`t`$ has essentially the same form as that corresponding to the Euclidean lattice case. The mathematical method used to derive such a result had already been successfully applied to Euclidean lattices and the fractal case is a fairly straightforward generalization when the previous conditions are fulfilled. In order to check the goodness of the approximation, we carried out numerical simulations on a standard deterministic substrate (the two-dimensional Sierpinski gasket) obtaining reasonable agreement with the theoretical results, especially when theoretical first-order asymptotic corrective terms are considered. The performance of the theoretical expressions discussed closely resembles that attained for Euclidean lattices. However, a more definitive check of the theoretical expressions for $`S_N(t)`$ that include corrective terms is hindered by the uncertainty in the value of the subdominant parameters $`A`$, $`\mu `$, $`h_1`$, …that appear in the survival probability $`\mathrm{\Gamma }_t(r)`$. Unfortunately, this function cannot be completely determined by the rigorous renormalization scheme mentioned in Sec. II, so that the definitive determination of its subdominant terms by numerical (or other analytical procedures) is a problem for future work which will surely be beset by with the technical difficulties associated with the identification of these faint subdominant terms .
###### Acknowledgements.
This work has been supported by the DGICYT (Spain) through Grant No. PB97-1501 and by the Junta de Extremadura-Fondo Social Europeo through Grant No. IPR99C031.
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# Irreducibility and Compositeness in q-Deformed Harmonic Oscillator Algebras
## I Introduction
In the last decades $`q`$-deformed algebras have been object of interest in the literature and a great effort has been devoted to its understanding and development. In particular, the interest in q-deformed algebras resides in the fact that they are deformed versions of the standard Lie algebras, and give them back as the deformation parameter $`q`$ goes to unity. Furthermore, since it is known that the deformed algebras encompass a set of symmetries that is richer than that of the Lie algebras, one is tempted to recognize that quantum algebras can be the appropriate tool to be dealt with in describing symmetries of physical systems which cannot be properly treated within the Lie algebras, although the direct interpretation of the deformation in these cases is sometimes incomplete or even completely lacking. For instance, in some cases like the XXZ-model, where the ferromagnetic/antiferromagnetic nature of a spin $`\frac{1}{2}`$ chain of length $`N`$ can be simulated through the introduction of a $`q`$-deformed algebra, or the rotational bands in deformed nuclei and molecules which can be fitted via a $`q`$-rotor Hamiltonian, instead of using the variable moment of inertia (VMI model), the physical meaning of the deformation parameter is established. Notwithstanding this interpretation difficulty, from the original studies which appeared in connection with problems related to solvable statistical mechanics models and quantum inverse scattering theory, a solid development has emerged which encompass nowadays various branches of mathematical problems related to physical applications, such as deformed superalgebras, knot theories, noncommutative geometries and so on. The introduction of a $`q`$-deformed bosonic harmonic oscillator is a subject of great interest in this context and, as a tool for providing a boson realisation of the quantum algebra $`su_q(2)`$, brought to light new commutation relations which have been extensively discussed in the literature.
On the other hand, some concepts directly related to the arithmetical foundations of deformed algebras were well known to mathematicians since the last century. For instance, the Gauss polynomials appearing in restricted partition theory can be directly interpreted as a $`q`$-generalization of the standard binomials; as such, the Gauss polynomials, or the $`q`$-binomials, as they are sometimes known, also generalize the concept of number as well. In that form, the Gauss extension of the number concept, sometimes known as $`Q`$-number, is also related to the usual $`q`$-bracket of extensive use in deformed algebras. In this connection, if, in general, the surprising effectiveness of number theory seems not to be completely realized, the success of recent examples pervading several areas can be credited to the use of that branch of science: solvable models in statistical mechanics benefited from Rogers-Ramanujan-Baxter relations, computation and cryptography, the fourth test on general relativity, dynamical systems, and primitive-roots-of-unity-based reflecting gratings in concert halls have their very foundations on basic number theory and algorithms .
In this paper we want to address the question whether the extension of the number concept proposed by Gauss, namely the $`Q`$-numbers, can farther help us in the study of the $`q`$-deformed harmonic oscillator. To this aim we are directly guided by the central role played by the number concept in this context. Based on this, we introduce the $`Q`$-numbers as our starting point to define the action of the creation/annihilation operators on the Fock space states. From this we show how we obtain a version of the $`q`$-deformed harmonic oscillator algebra already discussed in the literature , the $`A_q`$ algebra. We also discuss how a second set of operators obeying the $`q`$-deformed harmonic oscillator algebra can be introduced, the $`\overline{A_q}`$ algebra, such that they satisfy the conjugate relations with respect to the $`A_q`$ algebra, and discuss some possible reductions of the algebra when we choose the allowed values of the deformation parameter $`q`$. The cases for real and roots-of-unity values of $`q`$ are analysed. Furthermore, we also show how the reducibility of the algebra space representation appears for the different values of $`q.`$
Using the algebras $`A_q`$ and $`\overline{A_q}`$ we introduce a self-adjoint $`q`$-deformed harmonic oscillator Hamiltonian, akin to that proposed by Floratos and Tomaras and related to a system of two anyons, which allows us to test the reducibility criteria discussed before. This allows us to separate the physical systems according to the different algebras obtained for the different values of the deformation parameter $`q`$. In this form, we show how it is possible to distinguish different subsystems within the original oscillator Hamiltonian when $`q`$ assume nonprimitive roots of unity values. In this sense, we discuss the possibility of uncovering the compositeness character of the so-called $`k`$-fermions when discussing the reducibility of the representation space for $`q`$-deformed oscillator algebra at the roots of unity.
This paper is organized as follows: Section II is devoted to a brief review of the Gauss polynomials ($`Q`$-numbers) and their basic properties. In section III we derive and discuss the $`q`$-oscillator algebras from the extended number concept and in section IV we present the conditions for the reducibility of the Fock space representation. $`q`$-Deformed oscillator Hamiltonians are discussed in Section V, where examples of how the reducibility conditions sieve the space representation into subspaces are also exhibited. Finally the conclusions are presented in Section VI.
## II Gauss Polynomials: $`Q`$-Numbers
The generating function of restricted partitions of a positive integer $`N`$ into at most $`m`$ parts, each $`n,`$ is written as
$$G(n,m;q)=\frac{\left(1q^{n+m}\right)\left(1q^{n+m1}\right)\mathrm{}\left(1q^{m+1}\right)}{\left(1q\right)\left(1q^2\right)\mathrm{}\left(1q^n\right)},$$
(1)
$`q1`$, and the Gauss polynomials are defined through the relation
$$\left[\begin{array}{c}n\hfill \\ m\hfill \end{array}\right]=G(nm,m;q),$$
(2)
which is valid for $`0mn,`$ and zero otherwise. The Gauss polynomial is a polynomial of degree $`m(nm)`$ in $`q`$ that presents a very important property, namely
$$\underset{q1}{lim}\left[\begin{array}{c}n\hfill \\ m\hfill \end{array}\right]=\left(\genfrac{}{}{0pt}{}{n}{m}\right),$$
(3)
where $`\left(\genfrac{}{}{0pt}{}{n}{m}\right)`$ is the standard binomial. Thus, we conclude that the Gauss polynomials generalize the concept of binomials and, furthermore, as a special and important case, with $`m=1`$, the Gauss polynomial, that is now denoted $`Q`$-number, extends the concept of number since
$$\underset{q1}{lim}\left[\begin{array}{c}n\hfill \\ 1\hfill \end{array}\right]=\left(\genfrac{}{}{0pt}{}{n}{1}\right)=n.$$
(4)
On the other hand, this polynomial also allows us to establish inner contact with some aspects of number theory, since when $`q`$ is a $`n`$th root of unity, we have
$$\left[\begin{array}{c}n\hfill \\ 1\hfill \end{array}\right]=1+q+q^2+\mathrm{}+q^{n1}=\frac{1q^n}{1q}=0.$$
(5)
This is the fundamental equation whose $`n`$ solutions are roots of unity; furthermore, for $`n`$ prime, $`n1`$ of these are primitive roots .
Besides those important properties, the Gauss polynomials also satisfy the additional following relations
$$\left[\begin{array}{c}n\hfill \\ 0\hfill \end{array}\right]=\left[\begin{array}{c}n\hfill \\ n\hfill \end{array}\right]=1,$$
(6)
$$\left[\begin{array}{c}n\hfill \\ m\hfill \end{array}\right]=\left[\begin{array}{c}n\hfill \\ nm\hfill \end{array}\right],$$
(7)
[nm]=[n1m]+qnm[
-n1-m1],delimited-[]𝑛𝑚delimited-[]fragmentsn1𝑚superscript𝑞𝑛𝑚delimited-[]
-n1-m1\left[\begin{tabular}[]{l}$n$\\
$m$\end{tabular}\right]=\left[\begin{tabular}[]{l}$n-1$\\
$m$\end{tabular}\right]+q^{n-m}\left[\begin{tabular}[]{l}$n-1$\\
$m-1$\end{tabular}\right], (8)
$$\left[\begin{array}{c}n\hfill \\ m\hfill \end{array}\right]=\left[\begin{array}{c}n1\hfill \\ m1\hfill \end{array}\right]+q^m\left[\begin{array}{c}n1\hfill \\ m\hfill \end{array}\right].$$
(9)
## III $`q`$-Oscillator Algebras
Let us consider, as our starting point, the standard Fock space generated by $`\left\{|n\right\},`$
$$a0=0,$$
(10)
$$a^{}n=\sqrt{n+1}n+1,an=\sqrt{n}n1,$$
(11)
and
$$\widehat{N}n=nn,$$
(12)
where the creation and annihilation operators obey the following commutation relations
$$aa^{}a^{}a=1;[\widehat{N},a^{}]=a^{};[\widehat{N},a]=a,$$
(13)
from which it follows that
$$\widehat{N}=a^{}a.$$
(14)
Since the number concept is inherent to the Fock description, we are strongly motivated by the results of the previous section to construct a new pair of creation and annihilation operators in such a form to deal with that generalized number concept. To this aim we introduce new operators, whose matrix elements in the Fock space involve the Gauss polynomials
$$a_{}0=0$$
(15)
$$a_+n=\sqrt{\left\{n+1\right\}_q}n+1$$
(16)
$$a_{}n=\sqrt{\left\{n\right\}_q}n1$$
(17)
$$\widehat{N}n=nn$$
(18)
$$[\widehat{N},a_+]=a_+$$
(19)
$$[\widehat{N},a_{}]=a_{},$$
(20)
although $`\widehat{N}`$ $``$ $`a_+a_{}`$. Here we have adopted the notation
$$\left\{n\right\}_q\left[\begin{array}{c}n\hfill \\ 1\hfill \end{array}\right].$$
(21)
We can pose now the question: what is the algebra satisfied by $`a_+`$ and $`a_{}`$? Since
$$a_{}a_+n=\left\{n+1\right\}_qn,$$
(22)
$$a_+a_{}n=\left\{n\right\}_qn,$$
(23)
and considering relation (9), we conclude that
$$a_{}a_+qa_+a_{}=1,$$
(24)
that is a q-deformed commutation relation as already exhibited in the literature. Let us denote relations (19, 20, 24) by $`A_q`$ algebra. We can construct an $`\overline{A}_q`$ algebra out of the relations conjugated to those defining the $`A_q`$ algebra (19, 20, 24):
$$[\widehat{N},a_+^{}]=a_+^{},$$
(25)
$$[\widehat{N},a_{}^{}]=+a_{}^{},$$
(26)
$$a_+^{}a_{}^{}q^{}a_{}^{}a_+^{}=1.$$
(27)
In principle, these operators act on the dual space (bra space) to the considered Fock (ket) space. However, we can infer the action of these operators onto the ket space just by using the orthonormality of the states $`|n`$. It yields
$$a_{}^{}n=\left(\sqrt{\left\{n+1\right\}_q}\right)^{}n+1,$$
(28)
and
$$a_+^{}n=\left(\sqrt{\left\{n\right\}_q}\right)^{}n1.$$
(29)
With the above results, it is possible to examine if there is an algebra relating the creation/anihilation operators of the $`A_q`$ algebra and their respective Hermitian conjugates, constituents of the $`\overline{A}_q`$ algebra. Using the action of these operators over the ket space, it is possible to obtain the following relations:
$$a_{}a_{}^{}=\left|\left\{\widehat{N}+1\right\}_q\right|,$$
(30)
$$a_{}^{}a_{}=\left|\left\{\widehat{N}\right\}_q\right|,$$
(31)
and similarly
$$a_+^{}a_+=\left|\left\{\widehat{N}+1\right\}_q\right|,$$
(32)
$$a_+a_+^{}=\left|\left\{\widehat{N}\right\}_q\right|.$$
(33)
Here we shall only consider cases when $`q`$ is real valued or a root of unity, which are the most commonly found cases in the literature.
For real $`q`$ it is possible to verify that
$$\left|\left\{\widehat{N}\right\}_q\right|=\left\{\widehat{N}\right\}_q,$$
(34)
which together with Eqs. (32) and (33), and the recurrence relation of the Gauss polynomials, Eq. (9), yields:
$$a_{}a_{}^{}qa_{}^{}a_{}=1.$$
(35)
Similarly
$$a_+^{}a_+qa_+a_+^{}=1.$$
(36)
In this case (real $`q`$), through Eqs. (16, 29), it is possible to identify $`a_{}^{}a_+`$.
When $`q`$ is the fundamental root of unity it can be, by its turn, verified that
$$\left|\left\{\widehat{N}\right\}_q\right|=\left[\widehat{N}\right]_{q^{1/2}},$$
(37)
where
$$\left[X\right]_q=\frac{q^Xq^X}{qq^1}$$
(38)
defines the $`q`$-bracket of $`X`$. Equation (37), together with Eqs. (32) and (33) yields
$$a_{}a_{}^{}q^{\frac{1}{2}}a_{}^{}a_{}=q^{\frac{\widehat{N}}{2}}.$$
(39)
Similarly
$$a_+^{}a_+q^{\frac{1}{2}}a_+a_+^{}=q^{\frac{\widehat{N}}{2}}.$$
(40)
The last two equations characterize the $`q`$-oscillator algebra introduced by Biedenharn and McFarlane .
On the other hand, when $`q`$ is a root of unity, except the fundamental one, Eq. (37) is no longer valid, instead
$$\left|\left\{\widehat{N}\right\}_q\right|=\left|\left[\widehat{N}\right]_{q^{1/2}}\right|.$$
(41)
Using the definition of the $`q`$-bracket, Eq. (38), when $`q`$ is a general root of unity, $`q_j=\mathrm{exp}(\frac{2\pi i}{m}j)`$, a relation between $`\left[k\right]_{q_j^{\frac{1}{2}}}`$ and its $`m`$complementar $`\left[mk\right]_{q_j^{\frac{1}{2}}}`$, can be directly obtained
$`\left[mk\right]_{q_j^{\frac{1}{2}}}={\displaystyle \frac{e^{i\frac{\pi }{m}j(mk)}e^{i\frac{\pi }{m}j(mk)}}{e^{i\frac{\pi }{m}j}e^{i\frac{\pi }{m}j}}}=\left(1\right)^{j1}{\displaystyle \frac{e^{i\frac{\pi }{m}jk}e^{i\frac{\pi }{m}jk}}{e^{i\frac{\pi }{m}j}e^{i\frac{\pi }{m}j}}}`$
$$\left[mk\right]_{q_j^{\frac{1}{2}}}=\left(1\right)^{j1}\left[k\right]_{q_j^{\frac{1}{2}}}.$$
(42)
When $`q`$ is the fundamental root of unity then $`j=1`$, and we have
$$\left[mk\right]_{q_1^{\frac{1}{2}}}=\left[k\right]_{q_1^{\frac{1}{2}}}.$$
(43)
Furthermore, for the case of the inverse of such root of unity, $`q_j^1=\mathrm{exp}\left(\frac{2\pi i}{m}j\right)=\mathrm{exp}\left[\frac{2\pi i}{m}(mj)\right]`$, we can verify in exactly the same way that
$$\left[k\right]_{q_j^{\frac{1}{2}}}=\left(1\right)^{k1}\left[k\right]_{q_j^{\frac{1}{2}}}.$$
(44)
Now, using Eqs. (42-44) we obtain the following additional relation
$$\left[mk\right]_{q_j^{\frac{1}{2}}}=\left(1\right)^{mk1}\left[mk\right]_{q_j^{\frac{1}{2}}}.$$
(45)
These relations will be shown to be useful when we deal with $`q`$-oscillator Hamiltonians in finite-dimensional spaces.
## IV Reducibility of the Fock Representation
Now, considering the actions of $`a_+`$ and $`a_{}`$ on the Fock representation, Eqs. (16) and (17), we will analyse its reducibility properties. The various possibilities are studied below.
### A First case: $`\left\{n\right\}_q0,n>0`$
All states of the $`\{|n\}`$ representation can be obtained through successive applications of $`a_+`$ over the vacuum. In that case, $`\{|n\}`$ is irreducible with respect to the algebra $`\{a_{},a_+,\widehat{N},I\}.`$
### B Second case: $`\left\{m\right\}_q=0,\left\{n\right\}_q0n,0<n<m`$
In that case
$$a_+|m1=0,$$
(46)
and also
$$a_{}|m=0.$$
(47)
From these results it follows that the subspace generated by $`\{|0,|1,\mathrm{},|m1\}`$ is invariant under the action of the set $`\{a_{},a_+,\widehat{N},I\}`$, and it is then an irrep of dimension $`m`$ of the deformed algebra. For all $`q1`$, i.e., deformed cases, the hypothesis $`\left\{m\right\}_q=0,`$ $`\left\{n\right\}_q0,n,0<n<m`$ can be written as
$$\frac{q^m1}{q1}=0,\frac{q^n1}{q1}0,n,0<n<m$$
(48)
$$q^m=1,q^n1,n,0<n<m,$$
(49)
which is the definition of the primitive $`m`$th roots of unity. Therefore there will always be irreps of dimension $`m`$ whenever $`q`$ is a primitive $`m`$th root of unity.
### C Third case: $`l,0<l<m\mathrm{}\left\{m\right\}_q=0,\left\{l\right\}_q=0`$
This is the equivalent to state that $`q`$ is a non-primitive root of unity. Let us also suppose that $`l`$ is the smallest integer satisfying the hypothesis above, i.e., it is the smallest number for which $`q^l=1.`$ Then
$$q^k1,k,0<k<l,$$
(50)
and therefore the subspace generated by $`\{|0,|1,\mathrm{},|m1\}`$ is reducible in irreps of dimension $`l.`$
Labelling the $`m1`$ roots of unity as
$$q_j=e^{2\pi i\frac{j}{m}},j=1,2,\mathrm{},m1,$$
(51)
and for $`r`$ the greatest common divisor (GCD) of $`m`$ and $`j`$, i.e.,
$`j`$ $`=`$ $`rs`$ (52)
$`m`$ $`=`$ $`rl,`$ (53)
where $`s/l`$ is an irreducible fraction, then
$$q_j=e^{2\pi i\frac{s}{l}}.$$
(54)
In this way, the subspace generated by $`\{|0,|1,\mathrm{},|m1\}`$ is reducible, as we saw, in irreps of dimension $`l`$, which is the smallest value for which $`q_j^l=1`$.
Then, for each $`m`$th root of unity labelled by $`j`$, the dimension of the irreps will be $`l=m/r`$, where $`r`$ is the GCD of $`j`$ and $`m`$, and the $`m`$dimensional representation breaks into $`r`$ irreps of dimension $`l=m/r`$.
## V $`q`$-Deformed Oscillator Hamiltonian
We can obtain a $`q`$-deformed Hermitian oscillator Hamiltonian from the deformed operator algebra presented above through the direct construction
$$H=\frac{1}{2}\mathrm{}\omega \left(a_{}a_{}^{}+a_{}^{}a_{}\right)=\frac{1}{2}\mathrm{}\omega \left(a_+a_+^{}+a_+^{}a_+\right),$$
(55)
that can be written, in general, as
$$H=\frac{1}{2}\mathrm{}\omega \left(\sqrt{\left\{\widehat{N}+1\right\}_q\left\{\widehat{N}+1\right\}_q^{}}+\sqrt{\left\{\widehat{N}\right\}_q\left\{\widehat{N}\right\}_q^{}}\right),$$
(56)
which is equivalent to
$$H=\frac{1}{2}\mathrm{}\omega \left(\left|\left\{\widehat{N}+1\right\}_q\right|+\left|\left\{\widehat{N}\right\}_q\right|\right).$$
(57)
As was discussed in the preceding sections
$$\left|\left\{\widehat{N}\right\}_q\right|=\{\begin{array}{c}\left\{\widehat{N}\right\}_q,\text{for }q\text{ a real number}\hfill \\ \left|\right[\widehat{N}\left]_{q^{1/2}}\right|,\text{for }q\text{ a root of unity.}\hfill \end{array}$$
(58)
So, for real $`q`$, the Hamiltonian, Eq. (55), is written as
$$H=\frac{1}{2}\mathrm{}\omega \left(\left\{\widehat{N}+1\right\}_q+\left\{\widehat{N}\right\}_q\right),$$
(59)
and, for $`q`$ being a root of unity, it can be easily seen to reduce to
$$H=\frac{1}{2}\mathrm{}\omega \left(\left|\left[\widehat{N}+1\right]_{q_j^{1/2}}\right|+\left|\left[\widehat{N}\right]_{q_j^{1/2}}\right|\right).$$
(60)
However, when $`q`$ is furthermore singled out as the fundamental primitive root of unity, the above expression, according to Eq. (37), reduces to
$$H=\frac{1}{2}\mathrm{}\omega \left(\left[\widehat{N}+1\right]_{q_1^{1/2}}+\left[\widehat{N}\right]_{q_1^{1/2}}\right),$$
(61)
which is the usual proposal for the $`q`$-deformed oscillator . This last expression has the symmetry $`qq^1`$, since this is a symmetry of the bracket itself.
Since the deformed oscillator Hamiltonian is written directly in terms of the brackets of the operator $`\widehat{N}`$, the reducibility properties of the Fock representation space will appear in the spectrum of that operator as well. In this sense, the spectrum will be broken into blocks associated to subspaces of prime dimension whenever the initial space dimension is a composite integer number and we work with the non-primitive roots of unity.
As an application of what has been presented above we will discuss some simple cases. In this connection, we need not to work with all the roots of unity due to the properties presented for the brackets in the previous sections. In fact, using relations (42-45), we need not calculate the matrices representing the Hamiltonian for some primitive roots.
First, let us consider $`m=2`$. In this case, the matrix representing the $`q`$-deformed oscillator Hamiltonian is directly written since we only have to work with the fundamental primitive root of unity, $`q_1^{1/2}=\mathrm{exp}\left(i\frac{\pi }{2}\right)`$. Using the fact that $`\left[2\right]_{q_1^{1/2}}=0`$, we get
$$H=\frac{1}{2}\mathrm{}\omega \left(\begin{array}{cc}\left[1\right]_{q_1^{1/2}}& 0\\ 0& \left[1\right]_{q_1^{1/2}}\end{array}\right)=\frac{1}{2}\mathrm{}\omega \left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(62)
For $`m=3`$, which is the next prime number, and also using Eq. (43), we get for the fundamental primitive root of unity, $`q_1^{1/2}=\mathrm{exp}\left(i\frac{\pi }{3}\right),`$
$$H=\frac{1}{2}\mathrm{}\omega \left(\begin{array}{ccc}\left[1\right]_{q_1^{1/2}}& 0& 0\\ 0& \left[1\right]_{_{q_1^{1/2}}}+\left[2\right]_{q_1^{1/2}}& 0\\ 0& 0& \left[2\right]_{q_1^{1/2}}\end{array}\right)=\frac{1}{2}\mathrm{}\omega \left(\begin{array}{ccc}1& 0& 0\\ 0& 2& 0\\ 0& 0& 1\end{array}\right).$$
(63)
If we now consider the case $`m=6`$, we can verify how the matrix representing the Hamiltonian breaks into blocks, each with a prime dimension, as occurs in the representation space of the $`q`$-deformed algebra. To this end, let us first of all consider the Hamiltonian associated to the fundamental primitive root of unity, $`q_1^{1/2}=\mathrm{exp}\left(i\frac{\pi }{6}\right)`$. In this case, using Eq. (42), we see that the matrix is also symmetric and has the form
$$H=\frac{1}{2}\mathrm{}\omega \left(\begin{array}{cccccc}1& & & & & \\ & 1+\left[2\right]_{q_1^{1/2}}& & & \text{0}& \\ & & \left[2\right]_{q_1^{1/2}}+\left[3\right]_{q_1^{1/2}}& & & \\ & & & \left[3\right]_{q_1^{1/2}}+\left[2\right]_{q_1^{1/2}}& & \\ & \text{0}& & & \left[2\right]_{q_1^{1/2}}+1& \\ & & & & & 1\end{array}\right).$$
(64)
Now, if we consider the non-primitive roots of $`m=6`$, we see that there are three of them, namely, $`q_2^{1/2}=\mathrm{exp}\left(i\frac{\pi }{6}2\right)`$, $`q_3^{1/2}=\mathrm{exp}\left(i\frac{\pi }{6}3\right)`$ and $`q_4^{1/2}=\mathrm{exp}\left(i\frac{\pi }{6}4\right)`$ respectively. In fact, $`q_2`$ is the inverse of $`q_4`$ and $`q_3`$ is its own inverse. For the first root the Hamiltonian matrix will be
$$H=\frac{1}{2}\mathrm{}\omega \left(\begin{array}{cccccc}1& & & & & \\ & 2& & & \text{0}& \\ & & 1& & & \\ & & & 1& & \\ & \text{0}& & & 2& \\ & & & & & 1\end{array}\right),$$
(65)
which breaks into two blocks, each one being the matrix associated to a $`m=3`$ $`q`$-deformed oscillator. On the other hand, for the second non-primitive root of unity, $`q_3^{1/2}=\mathrm{exp}\left(i\frac{\pi }{6}3\right)`$, we get
$$H=\frac{1}{2}\mathrm{}\omega \left(\begin{array}{cccccc}1& & & & & \\ & 1& & & \text{0}& \\ & & 1& & & \\ & & & 1& & \\ & \text{0}& & & 1& \\ & & & & & 1\end{array}\right).$$
(66)
The three blocks associated to the $`m=2q`$-deformed oscillator are readily seen in this case. Since the Hamiltonian is given by (60), and using Eq. (44), we can conclude that the same matrices would be obtained if the inverse roots of unity were used. Therefore, for prime dimension spaces the matrices representing the deformed oscillator Hamiltonian are irreducible for any primitive root of unity. For nonprime integer space dimension and deformations at the non-primitive roots of unity, the $`q`$-deformed oscillator represents in fact a composite system with as many irreducible constituents (diagonal blocks) as are the number of prime factors of the starting space dimension.
## VI Conclusions
In the present paper, starting from the Gauss extension of the number concept, we have reobtained the $`q`$-deformed harmonic oscillator algebra discussed in for general deformation parameter $`q`$. For the particular case of $`q`$ being a fundamental root of unity, we recover the deformed harmonic oscillator algebra satisfied by $`a_{(+)}`$ and $`a_{(+)}^{}`$ as introduced by Biedenharn and MacFarlane. On the other hand, some useful relations between the Gauss polynomials and the standard $`q`$-bracket have also been discussed for $`q`$ a root of unity.
A discussion on the dimensions of the Fock state space representation for $`q`$ real or a root of unity shows that they can be infinite as well as finite-dimensional depending on $`q`$ being real or a root of unity, as it has been already pointed out. However, as a further conclusion, it is also shown that, for the particular case of $`q`$ being selected as a non-primitive root of unity, the representation space, besides being finite-dimensional, also breaks into subspaces, the dimension of each block being clearly defined by the prime decomposition of the number characterizing the original space dimension. In this form, it is shown that the use of non-primitive roots of unity allows one to verify the reducibility character of the Fock space representation, which, by its turn, shows that the subspaces characterized by prime dimensions play the role of fundamental blocks within the full space.
A $`q`$-deformed harmonic oscillator Hamiltonian was presented which allowed us to fully exploit the Fock space reducibility discussed previously. For the cases when $`q`$ was a primitive root of unity, the Hamiltonian matrix only exhibited the usual symmetries of the Q-numbers, while for $`q`$ being a non-primitive root of unity (which will occur only when the space dimension is a composite number) the matrices reduced to submatrices along the diagonal, thus indicating that the original $`q`$-deformed oscillator is in fact made up of irreducible subsystems, each one of them of a prime dimension. This result strongly suggests the conclusion that the so called k-fermions, or quons, discussed in the context of roots of unity deformation, are not necessarily fundamental entities, but they may be, in some cases, composite systems made up of entities of more fundamental character. This characterization can be directly verified by studying the degree of reducibility of the space representation through the prime decomposition of the space dimension from which we started.
Acknowledgement: D.G. and B.M.P are partially supported by Conselho Nacional de Desenvolvimento Científico e Tecnológico, CNPq; J.T.L. is partially supported by PICDT/CAPES, and M.R. has a fellowship from Fundação de Amparo à Pesquisa do Estado de São Paulo, FAPESP.
## A Polychronakos Realization
We start from the fundamental relation (24) which clearly reduces to the classical oscillator algebra when $`q1`$. We recall the Polychronakos realization :
$`a_{}`$ $`=`$ $`U_{}(q,\widehat{N})a`$ (A1)
$`a_+`$ $`=`$ $`U_+(q,\widehat{N})a^{},`$ (A2)
where $`\widehat{N}`$ is the usual nondeformed number operator. Using Eq. (A2) in Eq. (24) we obtain:
$$F(q,\widehat{N}+1)qF(q,\widehat{N})=1,$$
(A3)
where
$$F(q,\widehat{N})=U_+(q,\widehat{N})U_{}(q,\widehat{N}1)\widehat{N}.$$
(A4)
Representing Eq. (A4) on the Fock space $`\{|n\}`$ we get:
$$F(q,n+1)qF(q,n)=1,$$
(A5)
which is the recurrence relation (9) for the Gauss polynomials. We then may infer
$`a_+a_{}`$ $`=`$ $`g(q,\widehat{N})`$ (A6)
$`a_{}a_+`$ $`=`$ $`g(q,\widehat{N}+1),`$ (A7)
or in terms of $`F`$
$$F(q,\widehat{N})=g(q,\widehat{N}).$$
(A8)
In order to fulfil the deformed algebra, it is enough that
$$U_+(q,\widehat{N})U_{}(q,\widehat{N}1)\widehat{N}=g(q,\widehat{N}).$$
(A9)
Choosing
$$U_{}(q,\widehat{N}1)=\sqrt{\frac{\left\{\widehat{N}\right\}_q}{\widehat{N}}},U_{}(q,\widehat{N})=\sqrt{\frac{\left\{\widehat{N}+1\right\}_q}{\widehat{N}+1}},$$
(A10)
we obtain
$$U_+(q,\widehat{N})=\sqrt{\frac{\left\{\widehat{N}\right\}_q}{\widehat{N}}}.$$
(A11)
Therefore
$`a_{}`$ $`=`$ $`a\sqrt{{\displaystyle \frac{\left\{\widehat{N}\right\}_q}{\widehat{N}}}}`$ (A13)
$`a_+`$ $`=`$ $`a^{}\sqrt{{\displaystyle \frac{\left\{\widehat{N}+1\right\}_q}{\widehat{N}+1}}.}`$ (A14)
This choice guarantees the unitarity $`(a_+=a_{}^{})`$ when $`q`$ is a real parameter. Otherwise, the representation turns out to be non-unitary.
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# Energy distribution analysis of the wavepacket simulations of CH4 and CD4 scattering
## 1 Introduction
The dissociation of methane on transition metals is an important reaction in catalysis. It is the rate limiting step in steam reforming to produce syngas. It is also prototypical for C–H activation in other processes. A large number of molecular beam experiments in which the dissociation energy was measured as a function of translational energy have already been done on this system. These experiments have contributed much to our understanding of the mechanism of the dissociation. Some of them observed that vibrationally hot $`\mathrm{CH}_4`$ dissociates more readily than cold $`\mathrm{CH}_4`$, with the energy in the internal vibrations being about as effective as the translational energy in inducing dissociation. A more detailed assessment of the importance of the internal vibrations could not be made, because of the large number of internal vibrations. A recent molecular beam experiment with laser excitation of the $`\nu _3`$ mode succeeded in measuring a dramatical enhancement of the dissociation on a Ni(100) surface, but it is still much too low to account for the vibrational activation observed in previous studies and indicates that other vibrationally excited modes contribute significantly to the reactivity of thermal samples.
Wavepacket simulations are being used more and more to study the dynamics of this kind of molecule surface reactions. The published wavepacket simulations on the methane dissociation reaction on transition metals have treated the methane molecule always as a diatomic up to now. Besides the C–H bond and molecule surface distance, a combination of other coordinates were included, like (multiple) rotations and some lattice motion. None of them have looked at the role of the internal vibrations. Various theoretical studies have obtained reaction pathways and barriers for dissociation by DFT calculations, but they cannot explain the role of the vibrational modes in the reaction dynamics either.
A nice way to study reaction dynamics is the use of isotopes. The most recent wavepacket simulation on the dissociation probability of CH<sub>4</sub> and CD<sub>4</sub> showed a semiquantitative agreement with the molecular beam experiments of Ref., except for the isotope effect and the extracted vibrational efficacy. The molecular beam study with laser excitation of the $`\nu _3`$ asymmetrical stretch mode shows that the incorrect vibrational efficacy is caused by the assumptions in the fit procedure that both stretch modes behaves identical. One of the possible explanation of the incorrect isotope effect can be the role played by the non-included intramolecular vibrations.
In a previous paper we reported on wavepacket simulations that we have done to determine which and to what extent internal vibrations are important for the dissociation of CH<sub>4</sub>. We were not able yet to simulate the dissociation including all internal vibrations. Instead we simulated the scattering of methane, for which all internal vibrations can be included, and used the results to deduce consequences for the dissociation. We used model potential energy surfaces (PESs) that have been developed with Ni(111) in mind, but our results should hold for other surfaces as well. At a translational energy up to 96 kJ/mol we found that the scattering is almost completely elastic. Vibrational excitations when the molecule hits the surface and the corresponding deformation depend on generic features of the potential energy surface. In particular, our simulations indicate that for methane to dissociate the interaction of the molecule with the surface should lead to an elongated equilibrium C–H bond length close to the surface.
We have been using the multiconfigurational time-dependent Hartree (MCTDH) method for our wavepacket simulation, because it can deal with a large number of degrees of freedom and with large grids. This method has been applied successfully to gas phase reactions and reactions at surfaces.
In this paper we report wavepacket simulations of CD<sub>4</sub> scattering including all internal vibrations for fixed orientations, performed on the same model PESs as in our previous paper. Translational motion parallel to the surface and all rotational motion was neglected. No degrees of freedom of the surface were included. Experiments show that coupling with these degrees of freedom is dependent on the metal surface. For example, the observed surface temperature effect are small on Ni, but quite large on Pt. As we are only interested in the role of internal vibrations, we have not included degrees of freedom of the surface to keep the simulations as simple as possible. We will discuss the vibrational excitation and the deformation of the CD<sub>4</sub> molecule when it hits the surface and compare it with CH<sub>4</sub>. Later on we will look at the energy distribution of the kinetic energy per mode and the potential energy in some terms of the PES with the elongated equilibrium bond length close to the surface for both isotopes. The transfer of translational kinetic energy towards vibrational kinetic energy gives an indication about the dissociation probability, since vibrational kinetic energy helps in overcoming the dissociation barrier. It gives a better idea too about which modes are essential to include in a more accurate wavepacket simulation of methane dissociation. After that we will discuss the implications of this for the dissociation and give a summary with some general conclusions.
## 2 Computational details
### 2.1 The Potential Energy Surfaces
We used for the scattering of CD<sub>4</sub> the same model PESs as we did for CH<sub>4</sub>. Since we expressed the PES in mass-weighted coordinated the parameters in the PESs for CD<sub>4</sub> differs from CH<sub>4</sub>. We will now give an overview of our model PESs and the corresponding parameters for CD<sub>4</sub>. The parameters of CH<sub>4</sub> for these PESs were already given in Ref. , where also more detailed information about our assumptions and contour plots of some cross-section of the model PESs can be found.
The PESs we used can all be written as
$$V_{\mathrm{total}}=V_{\mathrm{intra}}+V_{\mathrm{surf}},$$
(1)
where $`V_{\mathrm{intra}}`$ is the intramolecular PES and $`V_{\mathrm{surf}}`$ is the repulsive interaction with the surface. For the $`V_{\mathrm{intra}}`$ we looked at four different types of PESs. The $`V_{\mathrm{intra}}`$ include also for two types changes in the intramolecular potential due to interactions with the surface.
#### 2.1.1 A harmonic potential
The first one is completely harmonic. We have used normal mode coordinates for the internal vibrations, because these are coupled only very weakly. In the harmonic approximation this coupling is even absent so that we can write $`V_{\mathrm{intra}}`$ as
$$V_{\mathrm{intra}}=V_{\mathrm{harm}}=\frac{1}{2}\underset{i=2}{\overset{10}{}}k_iX_i^2,$$
(2)
the summation is over the internal vibrations, $`X_i`$’s are mass-weighted displacement coordinates and $`k_i`$ are mass-weighted force constants. (see Table 1 for definitions and values); ($`X_1`$ is the mass-weighted overall translation along the surface normal). The force constants have been obtained by fitting them on the experimental vibrational frequencies of CH<sub>4</sub> and CD<sub>4</sub>.
We have assumed that the repulsive interaction with the surface is only through the deuterium atoms that point towards the surface. We take the $`z`$-axis as the surface normal. In this case the surface PES is given by
$$V_{\mathrm{surf}}=\frac{A}{N_D}\underset{i=1}{\overset{N_D}{}}e^{\alpha z_i},$$
(3)
where $`N_D`$ is the number of deuteriums that points towards the surface, $`\alpha `$=1.0726 atomic units and $`A`$=6.4127 Hartree. These parameters are chosen to give the same repulsion as the PES that has been used in an MCTDH wavepacket simulation of CH<sub>4</sub> dissociation.
If we write $`V_{\mathrm{surf}}`$ in terms of normal mode coordinates, then we obtain for one deuterium pointing towards the surface
$$V_{\mathrm{surf}}=Ae^{\alpha _1X_1}e^{\alpha _2X_2}e^{\alpha _3X_3}e^{\alpha _4X_4},$$
(4)
where $`A`$ as above, and $`\alpha `$’s as given in Table 2. $`X_2`$, $`X_3`$ and $`X_4`$ correspond all to $`a_1`$ modes of the C<sub>3v</sub> symmetry (see Fig. 1). There is no coupling between the modes $`X_5`$ to $`X_{10}`$ in the $`V_{\mathrm{surf}}`$ part of the PES, which are all $`e`$ modes of the C<sub>3v</sub> symmetry.
For two deuteriums we obtain
$`V_{\mathrm{surf}}=A`$ $`e^{\alpha _1X_1}e^{\alpha _2X_2}e^{\alpha _3X_3}e^{\alpha _4X_4}e^{\alpha _5X_5}`$ (5)
$`\times {\displaystyle \frac{1}{2}}[`$ $`e^{\beta _3X_7}e^{\beta _3X_8}e^{\beta _5X_9}e^{\beta _5X_{10}}`$
$`+`$ $`e^{\beta _3X_7}e^{\beta _3X_8}e^{\beta _5X_9}e^{\beta _5X_{10}}],`$
with $`A`$ again as above, $`\alpha `$’s and $`\beta `$’s as given in Table 2. $`X_2`$, $`X_3`$, $`X_4`$ and $`X_5`$ correspond all to $`a_1`$ modes of C<sub>2v</sub>. $`X_7`$, $`X_8`$, $`X_9`$ and $`X_{10}`$ correspond to $`b_1`$ and $`b_2`$ modes of C<sub>2v</sub>. $`X_6`$ corresponds to the $`a_2`$ mode of C<sub>2v</sub> and has no coupling with the other modes in $`V_{\mathrm{surf}}`$.
For three deuteriums we obtain
$`V_{\mathrm{surf}}=A`$ $`e^{\alpha _1X_1}e^{\alpha _2X_2}e^{\alpha _3X_3}e^{\alpha _4X_4}`$ (6)
$`\times {\displaystyle \frac{1}{3}}[`$ $`e^{\beta _1X_5}e^{\beta _2X_6}e^{\beta _3X_7}e^{\beta _4X_8}e^{\beta _5X_9}e^{\beta _6X_{10}}`$
$`+`$ $`e^{\beta _1X_5}e^{\beta _2X_6}e^{\beta _3X_7}e^{\beta _4X_8}e^{\beta _5X_9}e^{\beta _6X_{10}}`$
$`+`$ $`e^{2\beta _1X_5}e^{2\beta _3X_7}e^{2\beta _5X_9}],`$
with $`A`$ again as above, $`\alpha `$’s and $`\beta `$’s as given in Table 2. $`X_2`$, $`X_3`$ and $`X_4`$ corresponds to $`a_1`$ modes in the C<sub>3v</sub> symmetry (see Fig. 1). Because these last six coordinates correspond to degenerate $`e`$ modes of the C<sub>3v</sub> symmetry, the $`\beta `$ parameters are not unique.
#### 2.1.2 An anharmonic intramolecular potential
Even though we do not try to describe the dissociation of methane in this and our previous paper, we do want to determine which internal vibration might be important for this dissociation. The PES should at least allow the molecule to partially distort as when dissociating. The harmonic PES does not do this. A number of changes have therefor been made. The first is that we have describe the C–D bond by a Morse PES.
$$V_{\mathrm{Morse}}=D_e\underset{i=1}{\overset{4}{}}\left[1e^{\gamma \mathrm{\Delta }r_i}\right]^2,$$
(7)
where $`D_e=0.1828`$ Hartree (the dissociation energy of methane in the gas-phase) and $`\mathrm{\Delta }r_i`$ the change in bond length from the equilibrium distance. $`\gamma `$ was calculated by equating the second derivatives along one bond of the harmonic and the Morse PES. If we transform Eq. (7) back into normal mode coordinates, we obtain
$$V_{\mathrm{Morse}}=D_e\underset{i=1}{\overset{4}{}}\left[1e^{\gamma _{i2}X_2}e^{\gamma _{i3}X_3}e^{\gamma _{i4}X_4}e^{\gamma _{i7}X_7}e^{\gamma _{i8}X_8}e^{\gamma _{i9}X_9}e^{\gamma _{i,10}X_{10}}\right]^2,$$
(8)
with $`D_e`$ as above. $`\gamma `$’s are given in Tables 3 and 4. Note that, although we have only changed the PES of the bond lengths, the $`\nu _4`$ umbrella modes are also affected. This is because these modes are not only bending, but also contain some changes of bond length.
The new intramolecular PES now becomes
$$V_{\mathrm{intra}}=V_{\mathrm{harm}}+V_{\mathrm{Morse}}V_{\mathrm{corr}},$$
(9)
where $`V_{\mathrm{harm}}`$ is as given in Eq. (2) and $`V_{\mathrm{corr}}`$ is the quadratic part of $`V_{\mathrm{Morse}}`$, which is already in $`V_{\mathrm{harm}}`$.
#### 2.1.3 Intramolecular potential with weakening C–D bonds
When the methane molecule approach the surface the overlap of substrate orbitals and anti-bonding orbitals of the molecule weakens the C–D bonds. We want to include this effect for the C–D bonds of the deuteriums pointing towards the surface. We have redefined the $`V_{\mathrm{Morse}}`$ given in Eq. (8) and also replace it in Eq. (9). A sigmoidal function is used to switch from the gas phase C–D bond to a bond close to the surface. We have used the following, somewhat arbitrary, approximations. (i) The point of inflection should be at a reasonable distance from the surface. It is set to the turnaround point for a rigid methane molecule with translation energy 93.2 kJ/mol plus twice the fall-off distance of the interaction with the surface. (ii) The depth of the PES of the C–D bond is 480 kJ/mol in the gas phase, but only 93.2 kJ/mol near the surface. The value 93.2 kJ/mol corresponds to the height of the activation barrier used in our dissociation. (iii) The exponential factor is the same as for the interaction with the surface.
If we transform to normal-mode coordinates for the particular orientations, we then obtain
$$V_{\mathrm{weak}}=D_e\underset{i=1}{\overset{4}{}}W_i\left[1e^{\gamma _{i2}X_2}e^{\gamma _{i3}X_3}e^{\gamma _{i4}X_4}e^{\gamma _{i7}X_7}e^{\gamma _{i8}X_8}e^{\gamma _{i9}X_9}e^{\gamma _{i,10}X_{10}}\right]^2,$$
(10)
where $`W_i=1`$ for non-interacting bonds and
$$W_i=\frac{1+\mathrm{\Omega }e^{\alpha _1X_1+\omega }}{1+e^{\alpha _1X_1+\omega }}$$
(11)
for the interacting bonds pointing towards the surface. $`\alpha _1`$ is as given in Table 2, $`\gamma `$’s are given in Tables 3 and 4, $`\mathrm{\Omega }=1.94210^1`$ and $`\omega =7.197`$.
#### 2.1.4 Intramolecular potential with elongation of the C–D bonds
A weakened bond generally has not only a reduced bond strength, but also an increased bond length. We include the effect of the elongation of the C–D bond length of the deuteriums that point towards the surface due to interactions with the surface. We have redefined the $`V_{\mathrm{Morse}}`$ given in Eq. (8) and also replace it in Eq. (9) for this type of PES. We have used the following approximations: (i) The transition state, as determined by Ref. and , has a C–H bond that is 0.54 Å longer than normal. This elongation should occur at the turn around point for a rigid methane molecule with a translation energy of 93.2 kJ/mol. (ii) The exponential factor is again the same as for the interaction with the surface.
If we transform to normal-mode coordinates for the particular orientations, then we obtain
$$V_{\mathrm{shift}}=D_e\underset{i=1}{\overset{4}{}}\left[1e^{\gamma _{i2}X_2}e^{\gamma _{i3}X_3}e^{\gamma _{i4}X_4}e^{\gamma _{i7}X_7}e^{\gamma _{i8}X_8}e^{\gamma _{i9}X_9}e^{\gamma _{i,10}X_{10}}\mathrm{exp}[S_ie^{\alpha _1X_1}]\right]^2,$$
(12)
where $`\alpha _1`$ is as given in Table 2, $`\gamma `$’s are given in Tables 3 and 4. For orientation with one deuterium towards the surface we obtain; $`S_1=2.94210^2`$ and $`S_2=S_3=S_4=0`$, with two deuteriums; $`S_1=S_2=0`$ and $`S_3=S_4=1.69810^2`$, and with three deuteriums; $`S_1=0`$ and $`S_2=S_3=S_4=2.94210^2`$.
### 2.2 Initial States
The exact wave-function of a $`D`$-dimensional system, is expressed in the MCTDH approximation by the form
$$\mathrm{\Psi }_{\mathrm{MCTDH}}(q_1,\mathrm{},q_D;t)=\underset{n_1\mathrm{}n_D}{}c_{n_1\mathrm{}n_D}(t)\psi _{n_1}^{(1)}(q_1;t)\mathrm{}\psi _{n_D}^{(D)}(q_D;t).$$
(13)
All initial states in the simulations start with the vibrational ground state. The initial translational part $`\psi ^{(\mathrm{tr})}`$ is represented by a Gaussian wave-packet,
$$\psi ^{(\mathrm{tr})}(X_1)=(2\pi \sigma ^2)^{1/4}\mathrm{exp}\left[\frac{(X_1X_0)^2}{4\sigma ^2}+iP_1X_1\right],$$
(14)
where $`\sigma `$ is the width of the wave-packet (we used $`\sigma =320.248`$ atomic units), $`X_0`$ is the initial position (we used $`X_0=11\sigma `$, which is far enough from the surface to observe no repulsion) and $`P_1`$ is the initial momentum. Since we used mass-weighted coordinates the Gaussian wavepacket are identical for CD<sub>4</sub> and CH<sub>4</sub>. We performed simulations in the energy range of 32 to 96 kJ/mol. We here present only the results of 96 kJ/mol (equivalent to $`P_1=0.2704`$ atomic units), because they showed the most obvious excitation probabilities for $`V_{\mathrm{Morse}}`$. We used seven natural single-particle states, 512 grid points and a grid-length of $`15\sigma `$ for the translational coordinate. With this grid-width we can perform simulation with a translational energy up to 144 kJ/mol.
Gauss-Hermite discrete-variable representations (DVR) were used to represent the wavepackets of the vibrational modes. We used for all simulations of CD<sub>4</sub> the same number of DVR points as for CH<sub>4</sub>, which was 5 DVR points for the $`\nu _2`$ bending modes and 8 DVR points for the $`\nu _4`$ umbrella, $`\nu _3`$ asymmetrical stretch, and $`\nu _1`$ symmetrical stretch mode for an numerical exact integration, except for the simulations with $`V_{\mathrm{shift}}`$, where we used 16 DVR points for the $`\nu _1`$ symmetrical stretch mode, because of the change in the equilibrium position.
Also the same configurational basis was used for both isotopes. We did the simulation with one bond pointing towards the surface in eight dimensions, because the $`\nu _2`$ bending modes $`X_5`$ and $`X_6`$ do not couple with the other modes. We needed four natural single-particle states for modes $`X_2`$, $`X_3`$ and $`X_4`$, and just one for the others. So the number of configurations was $`7^14^31^4=448`$. The simulation with two bonds pointing towards the surface was performed in nine dimensions. One of the $`\nu _2`$ bending mode ($`X_6`$) does not couple with the other modes, but for the other mode $`X_5`$ we needed four natural single-particle states. The number of configurations was $`7^14^41^4=1792`$, because we needed the same number of natural single-particle states as mentioned above for the other modes. We needed ten dimensions to perform the simulation with three bonds pointing towards the surface. We used here one natural single-particle state for the modes $`X_5`$ to $`X_{10}`$ and four natural single-particle states for $`X_2`$ to $`X_4`$, which gave us $`7^14^31^6=448`$ configurations.
## 3 Results and Discussion
### 3.1 Excitation probabilities and structure deformation of CD<sub>4</sub>
The scattering probabilities for CD<sub>4</sub> are predominantly elastic, as we also found in our previous simulations of CH<sub>4</sub> scattering. The elastic scattering probability is larger than 0.99 for all orientation of the PESs with $`V_{\mathrm{Morse}}`$ and $`V_{\mathrm{weak}}`$ at a translational energy of 96 kJ/mol. For the PES with $`V_{\mathrm{shift}}`$ we observe an elastic scattering probability of 0.981 for the orientation with one, 0.955 with two and 0.892 with three deuteriums pointing towards the surface. This is lower than we have found for CH<sub>4</sub>, which is 0.956 for the orientation with three hydrogens pointing towards the surface and larger than 0.99 for the others. The higher inelastic scattering probabilities of CD<sub>4</sub> was expected, because the force constants $`k_i`$ of CD<sub>4</sub> are decreased up to 50% with respect to those of CH<sub>4</sub> and the translational surface repulsion fall-off differs only little.
When we look at the excitation probabilities at the surface for the PES with $`V_{\mathrm{Morse}}`$ and $`V_{\mathrm{weak}}`$, then we observe generally an increase for CD<sub>4</sub> compared with CH<sub>4</sub>, except for the $`\nu _4`$ umbrella mode in the orientation with two bond pointing towards the surface. Relevant differences in the structure deformations are observed only in the bond angles, which are increased for CD<sub>4</sub> in the orientations with one and three bonds pointing towards the surface. The bond angle deformation of the angle between the bonds pointing towards the surface in the orientation with two bonds pointing towards the surface is decreased for CD<sub>4</sub>. We observe again that the PES with $`V_{\mathrm{weak}}`$ gives larger structure deformations than the PES with $`V_{\mathrm{Morse}}`$, but the differences are smaller for CD<sub>4</sub> than CH<sub>4</sub>.
For the PES with $`V_{\mathrm{shift}}`$ we do not observe this effect on the bond angle deformation. The bond angle deformation for the orientation with two and three deuteriums pointing towards the surface is the same as for CH<sub>4</sub> and it is just $`0.1^{}`$ less for the bond angle at the surface side in the orientation with one deuterium pointing towards the surface. The excitation probabilities (see Table 5) for the $`\nu _2`$ bending and $`\nu _4`$ umbrella modes become higher for all orientations for CD<sub>4</sub>, which is necessary for getting the same bond angle deformations as CH<sub>4</sub>.
The changes in the bond distances for the orientations with one and two bonds pointing towards the surface is for CD<sub>4</sub> almost the same as for CH<sub>4</sub>. For the orientation with three bonds pointing towards the surface, we found that the maximum bond lengthening of the bonds on the surface side was $`0.032`$ Å less for CD<sub>4</sub> than CH<sub>4</sub>. We also found that the bond shortening of the bond pointing away from the surface is $`0.010`$ Å more for CD<sub>4</sub>. These are only minimal differences, which also only suggest that the bond deformation for CD<sub>4</sub> has been influenced slightly more by the $`\nu _3`$ asymmetrical stretch mode than the $`\nu _1`$ symmetrical stretch mode. The observed excitation probabilities for these modes do not contradict this, but are not reliable enough for hard conclusions because of their high magnitude. It is also not clear, beside of this problem, what they really represent. Is the excitation caused by a different equilibrium position of the PES at the surface in a mode or is it caused by extra energy in this mode? To answer these questions we decided to do an energy distribution analysis during the scattering for both isotopes.
### 3.2 Energy distribution in CH<sub>4</sub> and CD<sub>4</sub>
The energy distribution analysis is performed by calculating the expectation values of the important term of the Hamiltonian $`H`$ for the wave-function $`\mathrm{\Psi }(t)`$ at a certain time $`t`$ during the scattering of CD<sub>4</sub> and CH<sub>4</sub> for all presented orientations in this and our previous paper. We will present here only the results of the PES with $`V_{\mathrm{shift}}`$, because it is the only model PES for which the energy distribution analysis is relevant for the discussion of the dissociation hypotheses later on.
We can obtain good information about the energy distribution per mode by looking at the kinetic energy expectation values $`\mathrm{\Psi }(t)|T_j|\mathrm{\Psi }(t)`$ per mode $`j`$ (see Table 6), because the kinetic energy operators $`T_j`$ have no cross terms like the PESs have. When we discuss the kinetic energy of a mode we normally refer to the $`a_1`$ mode of the C<sub>3v</sub> or C<sub>2v</sub> symmetry, because in these modes we have observed the highest excitation probabilities and the change in kinetic energy in the other modes is generally small.
By looking at the expectation values of some terms of the PES $`\mathrm{\Psi }(t)|V_{\mathrm{term}}|\mathrm{\Psi }(t)`$ (see Table 7), we obtain information about how the kinetics of the scattering is driven by the PES. The $`V_{\mathrm{surf}}`$ PES \[see Eqs. (4), (5) and (6)\] is the surface hydrogen/deuterium repulsion for a given orientation. $`V_{\mathrm{harm}}(\nu _2)`$ and $`V_{\mathrm{harm}}(\nu _4)`$ \[see Eq. (2)\] are the pure harmonic terms of the intramolecular PES of the $`a_1`$ modes in the C<sub>3v</sub> and C<sub>2v</sub> symmetry corresponding to a $`\nu _2`$ bending and $`\nu _4`$ umbrella modes, respectively. The pure harmonic correction terms of $`V_{\mathrm{corr}}`$ \[see Eq. (9)\] are included in them. $`V_{\mathrm{bond}}(R_{\mathrm{up}})`$ and $`V_{\mathrm{bond}}(R_{\mathrm{down}})`$ are the potential energy in a single C–H or C–D bond pointing respectively towards and away from the surface, and they give the expectation value of one bond term of $`V_{\mathrm{shift}}`$ \[see Eq. (12)\]. All given expectation values are the maximum deviation of the initial values, which effectively means the values at the moment the molecule hits the surface.
The largest changes in expectation values are, of course, in the kinetic energy of the translational mode. The translational kinetic energy does not become zero as we should expect in classical dynamics. The loss of translational kinetic energy is primary absorbed by the $`V_{\mathrm{surf}}`$ terms of the PESs. The expectation values of the $`V_{\mathrm{surf}}`$ terms show the ability of the hydrogens or deuteriums to come close to the metal surface, since in real space their exponential fall-offs are the same for both isotopes. For a rigid molecule the sum of the translational kinetic energy and $`V_{\mathrm{surf}}`$ should be constant, so all deviations of this sum have to be found back in the intramolecular kinetic energy and other PES terms.
We observe that both the minimum in the translational kinetic energy and the maximum in the $`V_{\mathrm{surf}}`$ terms were higher for CH<sub>4</sub> than CD<sub>4</sub>, so we have to find more increase in energy in the intramolecular modes and PES terms for CD<sub>4</sub> than CH<sub>4</sub>. We indeed do so and that can be one of the reasons we found higher inelastic scatter probabilities for CD<sub>4</sub> for the PES with $`V_{\mathrm{shift}}`$ .
For the orientations with one and two bonds pointing towards the surface we observe a large increase of the kinetic energy in the $`\nu _3`$ asymmetrical stretch mode. If we compare this with the excitation probabilities, we find that the kinetic energy analysis gives indeed a different view on the dynamics. For the orientation with two bond pointing towards the surface we have found for both isotopes very high excitation probabilities in the $`\nu _1`$ and $`\nu _3`$ stretch modes. We know now from the kinetic energy distribution that for the $`\nu _1`$ symmetrical stretch mode the high excitation probability is caused by the change of the equilibrium position of the $`\nu _1`$ mode in the PES and that for the $`\nu _3`$ stretch mode probably the PES also has become narrower.
For the orientation with three bonds pointing towards the surface we also obtain an large increase of the kinetic energy of the $`\nu _3`$ asymmetrical stretch mode, but we also find an even larger increase in the kinetic energy of the $`\nu _1`$ symmetrical stretch mode. The total kinetic energy was extremely large, because the kinetic energy of the translational mode becomes also much larger than for the other orientations. Because of this the $`V_{\mathrm{surf}}`$ terms had to be around twice as low as for the other orientations.
All $`V_{\mathrm{bond}}(R_{\mathrm{up}})`$ terms become lower compared to the initial value, especially in the orientation with two bond pointing towards the surface. In the orientation with one bond pointing towards the surface, the $`V_{\mathrm{bond}}(R_{\mathrm{down}})`$ term became higher. This is caused by the repulsion of $`V_{\mathrm{surf}}`$ in the direction of the bond. The increase of this PES term value is higher for CD<sub>4</sub> than CH<sub>4</sub>.
In the orientation with three bond pointing towards the surface we also observe a higher $`V_{\mathrm{bond}}(R_{\mathrm{down}})`$ value, with also the highest increase for CD<sub>4</sub>. In relation with the somewhat shorter bond distance for the $`R_{\mathrm{down}}`$ of CD<sub>4</sub> compared with CH<sub>4</sub>, which was also a bit lower compared with the other orientations, we know now that the hydrogens and especially the deuterium have problems in following the minimum energy path of the PES with $`V_{\mathrm{shift}}`$ during the scattering dynamics. This leads to higher kinetic energy in the vibrational modes, which results in more inelastic scattering.
The $`V_{\mathrm{harm}}(\nu _2)`$ term increases in respect to the initial value, but not as much as the increase of the $`V_{\mathrm{harm}}(\nu _4)`$ is for the orientation with two bonds pointing towards the surface. The values of $`V_{\mathrm{harm}}(\nu _4)`$ for CD<sub>4</sub> are even higher than for CH<sub>4</sub>. We observe also a larger increase of the kinetic energy in the $`\nu _4`$ umbrella mode for CD<sub>4</sub> than for CH<sub>4</sub>. So although there is somewhat more energy transfer to the vibrational modes for CD<sub>4</sub> than CH<sub>4</sub>, this extra vibrational energy is absorbed especially in the $`\nu _4`$ umbrella mode of CD<sub>4</sub>.
### 3.3 Dissociation hypotheses
We like now to discuss some possible implications of the scattering simulation for the isotope effect on the dissociation of methane. In our previous paper we have already drawn some conclusions about the possible reaction mechanism and which potential type would be necessary for dissociation. We found the direct breaking of a single C–H bond in the initial collision more reasonable than the splats model with single bond breaking after an intermediary Ni–C bond formation as suggested by Ref. , because the bond angle deformations seems to small to allow a Ni–C to form. From the simulations with CD<sub>4</sub> we can draw the same conclusions. The PES with $`V_{\mathrm{shift}}`$ gives the same angle deformations for both isotopes, which is not sufficient for the splats model. The other potentials give higher bond angle deformations for the orientation with three deuteriums pointing towards the surface. If the Ni–C bond formation would go along this reaction path, then the dissociation of CD<sub>4</sub> should be even more preferable than CH<sub>4</sub>, which is not the case. So we only have to discuss the implication of the scattering simulation for the dissociation probabilities of CH<sub>4</sub> and CD<sub>4</sub> for a direct breaking of a single bond reaction mechanism. This reaction mechanism can be influenced by what we will call a direct or an indirect effect.
A direct effects is the expected changes in the dissociation probability between CH<sub>4</sub> and CD<sub>4</sub> for a given orientation. Since we expect that we need for dissociation a PES with an elongation of the bonds pointing towards the surface, we only have to look at the isotope effect in the simulation for the PES with $`V_{\mathrm{shift}}`$ for different orientations to discuss some direct effect. It is clear from our simulations that the bond lengthening of CD<sub>4</sub> is smaller than CH<sub>4</sub> for the orientation with three bonds pointing towards the surface. If this orientation has a high contribution to the dissociation of methane, then this can be the reason of the higher dissociation probability of CH<sub>4</sub>. In this case our simulations also explain why Ref. did not observe a high enough isotope effect in the dissociation probability of their simulation with CH<sub>4</sub> and CD<sub>4</sub> modelled by a diatomic, because we do not observe a change in bond lengthening between the isotopes for the orientation with one bond pointing towards the surface.
The orientation with three bonds pointing towards the surface is also the orientation with the highest increase of the total vibrational kinetic energy for the PES with $`V_{\mathrm{shift}}`$, because the energy distribution analysis shows besides an high increase of the kinetic energy in the $`\nu _3`$ asymmetrical stretch mode also an high increase in the $`\nu _1`$ symmetrical stretch mode. Since vibrational kinetic energy can be used effectively to overcome the dissociation barrier, the orientation with three bonds indicates to be a more preferable orientation for dissociation. Moreover the relative difference in kinetic energy between both isotopes is for the $`\nu _1`$ stretch mode larger than for the $`\nu _3`$ stretch mode. If the kinetic energy in the $`\nu _1`$ stretch mode contributes significantly to overcoming the dissociation barrier, then it is another explanation for the low isotope effect in Ref..
An indirect effect is the expected changes in the dissociation probability between CH<sub>4</sub> and CD<sub>4</sub> through changes in the orientations distribution caused by the isotope effect in the vibrational modes. This can be the case if the favourable orientation for dissociation is not near the orientation with three bonds pointing towards the surface, but more in a region where one or two bonds pointing towards the surface. These orientations do not show a large difference in deformation for the PES with $`V_{\mathrm{shift}}`$. We can not draw immediate conclusion about the indirect effect from our simulations, since we did not include rotational motion, but our simulation show that an indirect isotope effect can exist. For the PES with $`V_{\mathrm{Morse}}`$ in the orientation with three bonds pointing towards the surface, we observe that CD<sub>4</sub> is able to come closer to the surface than CH<sub>4</sub>. So this rotational orientation should be more preferable for CD<sub>4</sub> than for CH<sub>4</sub>. On the other hand, if the PES is in this orientation more like $`V_{\mathrm{shift}}`$ the dissociation probability in other orientation can be decreased for CD<sub>4</sub> through higher probability in inelastic scattering channels.
So for both effects the behaviour of the orientation with three bonds pointing towards the surface seems to be essential for a reasonable description of the dissociation mechanism of methane. A wavepacket simulation of methane scattering including one or more rotational degrees of freedom and the vibrational stretch modes will be a good starting model to study the direct and indirect effects, since most of the kinetic energy changes are observed in the stretch modes and so the bending and umbrella modes are only relevant with accurate PESs. Eventually dissociation paths can be introduced in the PES along one or more bonds.
Beside of our descriptions of the possible isotopes effect for the dissociation extracted of the scatter simulations we have to keep in mind that also a tunneling mechanism can be highly responsible for the higher observed isotope effect in the experiment and that a different dissociation barrier in the simulations can enhance this effect of tunneling.
## 4 Conclusions
The scattering is in all cases predominantly elastic. However, the observed inelastic scattering is higher for CD<sub>4</sub> compared with previous simulation on CH<sub>4</sub> for the PES with an elongated equilibrium bond length close to the surface. When the molecule hits the surface, we observe in general a higher vibrational excitation for CD<sub>4</sub> than CH<sub>4</sub>. The PES with an elongated equilibrium bond length close to the surface gives for both isotopes almost the same deformations, although we observe a somewhat smaller bond lengthening for CD<sub>4</sub> in the orientation with three bonds pointing towards the surface. The other model PESs show differences in the bond angle deformations and in the distribution of the excitation probabilities of CD<sub>4</sub> and CH<sub>4</sub>, especially for the PES with only an anharmonic intramolecular potential.
Energy distribution analysis contributes new information on the scattering dynamics. A higher transfer of translational energy towards vibrational kinetic energy at the surface results in higher inelastic scattering. The highest increase of vibrational kinetic energy is found in the $`\nu _3`$ asymmetrical stretch modes for all orientations and also in the $`\nu _1`$ symmetrical stretch mode for the orientation with three bonds pointing towards the surface, when the PES has an elongated equilibrium bond length close to the surface.
Our simulations give an indication that the isotope effect in the methane dissociation is caused mostly by the difference in the scattering behaviour of the molecule in the orientation with three bonds pointing towards the surface. At least multiple vibrational stretch modes should be included for a reasonable description of isotope effect in the methane dissociation reaction.
## Acknowledgments
This research has been financially supported by the Council for Chemical Sciences of the Netherlands Organization for Scientific Research (CW-NWO). This work has been performed under the auspices of NIOK, the Netherlands Institute for Catalysis Research, Lab Report No. TUE-99-5-02.
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# Macroscopic material properties from quasi-static, microscopic simulations of a two-dimensional shear-cell
## 1 Introduction
Macroscopic continuum equations for the description of the behavior of granular media rely on constitutive equations for stress, strain, and other physical quantities describing the state of the system. One possible way of obtaining an observable like the stress is to perform discrete particle simulations herrmann98 ; cundall79 and to average over the “microscopic” quantities in the simulation, in order to obtain an averaged macroscopic quantity. In the literature, slightly different definitions for stress and strain averaging procedures can be found drescher72 ; rothenburg81 ; savage81 ; cundall82 ; goddard86 ; bathurst88 ; kruyt96 ; liao97b ; goldhirsch99 ; kuhl99 .
The outline of this study is as follows. In section 2 the discrete element simulation method is discussed and, in section 3 some averaging methods are introduced and applied to a scalar quantity, namely the volume fraction. Section 4 contains the definitions and averaging strategies for fabric, stress, and elastic strain and in section 5 some material properties are extracted from the results obtained in section 4.
## 2 Modelling discrete particles
The elementary units of granular materials are mesoscopic grains. In order to account for the excluded volume, one can assume that the grains are impenetrable but deform under stress. Since the realistic modelling of the deformations of the particles in the framework of a continuum theory walton93 ; lian96 would be much too complicated, we relate the interaction force to the overlap $`\delta `$ of two particles. Note that the evaluation of the inter particle forces based on the overlap may not be sufficient to account for the nonlinear stress distribution inside the particles. Consequently, our results presented below are of the same quality as this basic assumption.
The force laws used are material dependent, involving properties such as Young’s modulus of elasticity, and have to be validated by comparison with experimental measurements foerster94 ; labous97 ; falcon98 .
When all forces $`\stackrel{}{f}_i`$ acting on the particle $`i`$, either from other particles, from boundaries or from external forces, are known, the problem is reduced to the integration of Newton’s equations of motion for the translational and the rotational degrees of freedom
$$m_i\frac{\mathrm{d}^2}{\mathrm{d}t^2}\stackrel{}{r}_i=\stackrel{}{f}_i,\mathrm{and}I_i\frac{\mathrm{d}^2}{\mathrm{d}t^2}\stackrel{}{\mathrm{\Phi }}_i=\stackrel{}{M}_i.$$
(1)
The mass of particle $`i`$ with diameter $`d_i`$ is $`m_i`$, and its moment of inertia is $`I_i=q_im_i(a_i)^2`$, with the radius $`a_i=d_i/2`$ and the dimensionless shape factor $`q_i`$. The vectors $`\stackrel{}{r}_i`$ and $`\stackrel{}{\mathrm{\Phi }}_i`$ give the position and the orientation angles of particle $`i`$, respectively. In our model attractive forces and the presence of other phases are neglected, we focus on “dry granular media”. Particle-particle interactions are short range and active on contact only. The total force (torque) due to other particles is thus $`\stackrel{}{f}_i=_c\stackrel{}{f}_i^c`$ ($`\stackrel{}{M}_i=_c\stackrel{}{M}_i^c`$), where the sum runs over all contacts of particle $`i`$. The torque $`\stackrel{}{M}_i^c=\stackrel{}{l}_i^c\times \stackrel{}{f}_i^c`$ is related to the force $`\stackrel{}{f}_i^c`$ via the cross product with the branch vector $`\stackrel{}{l}_i^c`$ from the particle center to the point of contact $`c`$. Eq. (1) consists of six scalar equations in three dimensions and reduces to three equations in two dimensions (2D), two for the linear and one for the rotational degree of freedom. In the following the force laws for $`\stackrel{}{f}_i^c`$ accounting for excluded volume, dissipation, and friction are introduced.
### 2.1 Force laws
The particles $`i`$ and $`j`$ interact only when they are in contact so that their overlap $`\delta =\frac{1}{2}(d_i+d_j)(\stackrel{}{r}_i\stackrel{}{r}_j)\stackrel{}{n}`$ is positive, with the unit vector $`\stackrel{}{n}=(\stackrel{}{r}_i\stackrel{}{r}_j)/|\stackrel{}{r}_i\stackrel{}{r}_j|`$ that points from $`j`$ to $`i`$. The symbol ‘$``$’ denotes the scalar product of vectors or, more generally, the order-reduction by one for each of two tensors.
The first contribution to the force acting on particle $`i`$ from $`j`$ is an elastic repulsive force
$$\stackrel{}{f}_{\mathrm{n},\mathrm{el}}=k_n\delta \stackrel{}{n},$$
(2)
where $`k_n`$ is proportional to the material’s modulus of elasticity with units $`[N/m]`$. Since we are interested in disks rather than spheres, we use a linear spring that follows Hooke’s law, whereas in the case of elastic spheres, the Hertz contact law would be more appropriate hertz82 ; landau75 .
The second contribution, a viscous dissipation, is given by the damping force in the normal direction
$$\stackrel{}{f}_{\mathrm{n},\mathrm{diss}}=\gamma _n\dot{\delta }\stackrel{}{n},$$
(3)
where $`\gamma _n`$ is a phenomenological normal viscous dissipation coefficient with units \[kg s<sup>-1</sup>\] and $`\dot{\delta }=\stackrel{}{v}_{ij}\stackrel{}{n}`$ the relative velocity in the normal direction $`\stackrel{}{v}_{ij}=\stackrel{}{v}_i\stackrel{}{v}_j`$.
The third contribution to the contact force – accounting for tangential friction – can be chosen in the simplest case, according to Coulomb, as $`\stackrel{}{f}_{\mathrm{t},\mathrm{friction}}\mu |\stackrel{}{f_n}|\stackrel{}{t}`$, where $`\mu `$ is the friction coefficient and $`\stackrel{}{t}=\dot{\stackrel{}{\xi }}/|\dot{\stackrel{}{\xi }}|`$ is the tangential unit-vector parallel to the tangential component of the relative velocity $`\dot{\stackrel{}{\xi }}=\stackrel{}{v}_{ij}(\stackrel{}{v}_{ij}\stackrel{}{n})\stackrel{}{n}`$. Because this non-smooth ansatz leads to numerical problems for small $`\dot{\stackrel{}{\xi }}`$ a regularizing viscous force $`\stackrel{}{f}_{\mathrm{t},\mathrm{viscous}}=\gamma _t\dot{\stackrel{}{\xi }}`$ is added. The two forces are combined by taking the minimum value
$$\stackrel{}{f}_\mathrm{t}=\mathrm{min}(|\gamma _t\dot{\stackrel{}{\xi }}|,|\mu \stackrel{}{f_n}|)\stackrel{}{t}.$$
(4)
The effect of a more realistic tangential force law according to Cundall and Strack cundall83 will be reported elsewhere latzel99a . Due to the boundary conditions introduced below, it is also necessary to account for the friction with the bottom
$$\stackrel{}{f}_\mathrm{b}=\mu _\mathrm{b}mg\widehat{\stackrel{}{v}},$$
(5)
with the unit vector in the direction of the particles velocity $`\widehat{\stackrel{}{v}}=\stackrel{}{v}/|\stackrel{}{v}|`$. The effect of a bottom friction was discussed also in schollmann99 . In summary, we combine the forces and use
$$\stackrel{}{f}_i=\underset{c}{}(\stackrel{}{f}_{\mathrm{n},\mathrm{el}}+\stackrel{}{f}_{\mathrm{n},\mathrm{diss}}+\stackrel{}{f}_\mathrm{t})+\stackrel{}{f}_\mathrm{b}$$
(6)
for the forces acting on particle $`i`$ at its contact $`c`$ with particle $`j`$.
### 2.2 Model system
In the simulations presented in this study, a two-dimensional Couette shear-cell is used, filled with a bidisperse packing of disks, as sketched in Fig. 1. The system undergoes slow shearing introduced by turning the inner ring. The properties of the particles, used for the force laws of subsection 2.1, are summarized in table 1.
The boundary conditions are based on an experimental realization howell99 ; veje98b ; veje99 . For more details on other simulations, see veje98b ; schollmann99 .
In the simulations different global volume fractions
$$\overline{\nu }=\frac{1}{V_{\mathrm{tot}}}\underset{p=1}{\overset{N}{}}V^p$$
(7)
of the shear-cell are examined. The sum in Eq. (7) runs over all particles $`p`$ with volume $`V^p`$ in the cell with $`V_{\mathrm{tot}}=\pi (R_o^2R_i^2)`$. In this study $`\overline{\nu }`$ is varied between $`0.8084`$ and $`0.8194`$, by varying the particle number, see table 2. For the calculation of the global volume fraction, the small particles glued to the wall are counted with half their volume only, and thus contribute with $`\overline{\nu }_{\mathrm{wall}}=0.0047`$ to $`\overline{\nu }`$.
Note that we use the phrases “volume” and “volume fraction” even if, strictly speaking, the unfamiliar terms “disk area” and “area fraction” could be used. The reasons for this choice are: (i) The methods discussed in this study are straightforwardly generalized to three dimensions and (ii) the particles are three dimensional objects with height $`h`$ anyway, so that the use of the word “volume” is justified.
### 2.3 Initial conditions and steady state
The simulations are started in a dilute state with an extended outer ring $`R_o(t=0)>R_o=0.2524`$ m, and the inner ring already rotates counterclockwise with constant angular frequency $`\mathrm{\Omega }=2\pi /T_i=0.1`$ s<sup>-1</sup> and period $`T_i=62.83`$ s. The extended outer ring is used in order to allow for a random, dilute initial configuration. The desired density is then reached by reducing the volume. The radius of the outer ring is reduced within about two seconds to reach its desired value $`R_o`$. Afterwards, the outer ring is kept fixed and the inner ring continues to rotate until at $`t=t_{\mathrm{max}}`$ the simulation ends. If not explicitly mentioned, averages are performed after about one rotation at $`t=60`$ s (to get rid of the arbitrary initial configuration), and during about one rotation, until $`t=119`$ s.
## 3 From the micro- to a macro-description
In the previous section, the microscopic point of view was introduced, as used for the discrete element method. Particles are viewed as independent entities which interact when they come in contact. In this framework, the knowledge of the forces at each contact is sufficient to model the dynamics and the statics of the system. Tensorial quantities like the stress or the deformation gradient are not necessary for a discrete modelling. However, subject of current research is to establish a correspondence to continuum theories by computing tensorial quantities like the stress $`𝝈`$, the strain $`𝜺`$, as well as scalar material properties like, e.g., the bulk and shear moduli goddard86 ; kruyt96 ; liao97b . In the course of this process, we first discuss averaging strategies using the material density as an example.
### 3.1 Averaging strategies
Most of the measurable quantities in granular materials vary strongly on short distances. Thus, computing averages necessitates dealing with or smearing out the fluctuations. In order to suppress the fluctuations, we perform averages in both time and space. This is possible due to the chosen boundary condition. The system can run for long time in a quasi-steady state and, due to the cylindrical symmetry, points at a certain distance $`R`$ from the origin are equivalent to each other. Therefore, averages are taken over many snapshots in time with time steps $`\mathrm{\Delta }t`$ and on rings of material at a distance $`r=RR_i`$ from the inner ring. The width of the averaging rings is $`\mathrm{\Delta }r`$, so that the averaging volume of one ring is $`V_r=2\pi r\mathrm{\Delta }r`$. For the sake of simplicity (and since the procedure is not restricted to cylindrical symmetry), the averaging volume is denoted by $`V=V_r`$ in the following. The rings are numbered from $`s=0`$ to $`B1`$, with $`B=(R_oR_i)/\mathrm{\Delta }r`$, and ring $`s`$ reaches from $`r_s=r\mathrm{\Delta }r/2`$ to $`r_{s+1}=r+\mathrm{\Delta }r/2`$. The averaging over many snapshots is somehow equivalent to an ensemble average. However, we remark that different snapshots are not necessarily independent of each other as discussed in subsection 3.4. Also the duration of the simulation maybe not long enough to explore a representative part of the phase space.
The cylindrical symmetry is accounted for by a rotation of all directed quantities like vectors, depending on the cartesian position $`\stackrel{}{r}_i=(x_i,y_i)`$ of the corresponding particle $`i`$. The orientation of particle $`i`$ is $`\varphi _i=\mathrm{arctan}(y_i/x_i)`$ for $`x_i>0`$ and periodically continued for $`x_i<0`$ so that $`\varphi _i`$ can be found in the interval $`[\pi ,\pi ]`$. The vector $`\stackrel{}{n}^c`$ that corresponds to contact $`c`$ of particle $`i`$ is then rotated about the angle $`\varphi _i`$ from its cartesian orientation before being used for an average. Note that this does not correspond to a transformation into orthonormal cylindrical coordinates.
Finally, we should remark that the most drastic assumption used for our averaging procedure is the fact, that all quantities are smeared out over one particle. Since it cannot be the goal to solve for the stress field inside one particle, we assume that a measured quantity is constant inside the particle. This is almost true for the density, but not e.g. for the stress. However, since we average over all positions with similar distance from the origin, i.e. averages are performed over particles with different positions inside a ring, details of the position dependency inside the particles will be smeared out anyway. An alternative approach was recently proposed by I. Goldhirsch goldhirsch99 who smeared out the averaging quantities along the lines connecting the centers of the particles and weighed the contribution according to the fraction of this line within the averaging volume.
### 3.2 Volume fractions
The first quantity to measure is the local volume fraction
$$\nu =\nu (r)=\frac{1}{V}\underset{pV}{}w_V^pV^p$$
(8)
with the particle volume $`V^p`$. $`w_V^p`$ is the weight of the particle’s contribution to the average.
This formalism can be extended to obtain the mean value of a quantity $`Q`$ in the following way:
$$Q=Q^p=\frac{1}{V}\underset{pV}{}w_V^pV^pQ^p,$$
(9)
with the pre-averaged particle quantity
$$Q^p=\underset{c=1}{\overset{𝒞^p}{}}Q^c.$$
(10)
with the quantity $`Q^c`$ attributed to contact $`c`$ of particle $`p`$. In the following the brackets $`Q`$ denote the average over a quantity $`Q`$ of the particles in an averaging volume $`V`$.
The simplest choice for $`w_V^p`$ is to use $`w_V^p=1`$, if the center of the particle lies inside the ring, and $`w_V^p=0`$ otherwise. This method will be referred to as particle-center averaging in the following.
A more complicated way to account for those parts of the particles which lie partially inside a ring is to use only the fraction of the particle volume that is covered by the averaging volume. Since an exact calculation of the area of a disk that lies in an arbitrary ring is rather complicated, we assume that the boundaries of $`V`$ are locally straight, i.e. we cut the particle in slices, as shown in Fig. 2. This method will be referred to as slicing in the following. The error introduced by using straight cuts is well below one per-cent in all situations considered here.
The volume $`V_s^p=w_V^pV^p`$ of a particle $`p`$ which partially lies between $`r_s`$ and $`r_{s+1}`$ is calculated by subtracting the external volumes $`V_s^{}`$ and $`V_s^+`$ from the particle volume $`V^p=\pi (d/2)^2`$ so that
$`V_s^p`$ $`=`$ $`V^pV_s^{}V_s^+`$ (11)
$`=`$ $`(d/2)^2[\pi \varphi _s+\mathrm{sin}(\varphi _s)\mathrm{cos}(\varphi _s)`$
$`\varphi _{s+1}+\mathrm{sin}(\varphi _{s+1})\mathrm{cos}(\varphi _{s+1})]`$
with $`\varphi _s=\mathrm{arccos}(2(r^pr_s)/d)`$ and $`\varphi _{s+1}=\mathrm{arccos}(2(r_{s+1}r^p)/d)`$. The term $`(d/2)^2\varphi `$ is the area of the segment of the circle with angle $`2\varphi `$, and the term $`(d/2)^2\mathrm{sin}(\varphi )\mathrm{cos}(\varphi )`$ is the area of the triangle belonging to the segment. In Fig. 2 the case $`s=1`$ is highlighted, and the boundaries between $`V_s^{}`$, $`V_s^p`$, and $`V_s^+`$ are indicated as thick solid lines. The two outermost slices $`V_0^p=V_1^{}`$ and $`V_s^p=V_{s1}^+`$ have to be calculated separately.
When the two averaging strategies (particle-center and slicing) are compared for different widths $`\mathrm{\Delta }r`$ of the averaging rings, so that the number of intervals is $`B=20`$, 40, or 60. For $`\mathrm{\Delta }rd_{\mathrm{small}}`$ ($`B=20`$), the results are almost independent of the averaging procedure. For finer binning $`B_s30`$, the slicing procedure converges on a master curve with weak oscillations close to the walls which are due to the wall-induced layering of the particles. The results of the particle-center averaging strongly fluctuate for $`B_c24`$. These oscillations come from ordered layers of the particles close to the walls so that the slicing method reflects the real density distribution for fine enough binning. The particle-center method, on the other hand, leads to peaks, where the centers of the particles in a layer are situated and to much smaller densities where few particle centers are found; the particle-center density is obtained rather than the material density.
### 3.3 Representative elementary volume REV
An important question is, how does the result of an averaging procedure depend on the size of the averaging volume $`V`$. We combine time- and space averaging, i.e. we average over many snapshots and over rings of width $`\mathrm{\Delta }r`$, so that the remaining “size” of the averaging volume is the width of the rings $`\mathrm{\Delta }r`$. In Fig. 3 data for $`\nu `$ at fixed position $`R=0.12`$, 0.13, 0.14, and 0.20 m, but obtained with different width $`\mathrm{\Delta }r`$, are presented. The positions correspond to $`r/d_{\mathrm{small}}2.2`$, 3.6, 4.9, and 13, when made dimensionless with the diameter of the small particles. Both the particle-center method (open symbols) and the slicing method (solid symbols) are almost identical for $`\mathrm{\Delta }r/d_{\mathrm{small}}2`$, for the larger $`\mathrm{\Delta }r`$ the averaging volume can partially lie outside of the system. For very small $`\mathrm{\Delta }r/d_{\mathrm{small}}0.1`$ the different methods lead to strongly differing results, however, the values in the limit $`\mathrm{\Delta }r0`$ are consistent, i.e. independent of $`\mathrm{\Delta }r`$ besides statistical fluctuations. In the intermediate regime $`0.1`$ $`<\mathrm{\Delta }r/d_{\mathrm{small}}<2`$, the particle-center method strongly varies, while the slicing method shows a comparatively smooth variation.
Interestingly, all methods seem to collapse at $`\mathrm{\Delta }r_{\mathrm{REV}}`$ $`d_{\mathrm{small}}`$ (and twice this value), nearly to the size of the majority of the particles. For the examined situations, we observe that the particle-center and the slicing method lead to similar results for $`0.97\mathrm{\Delta }r_{\mathrm{REV}}/d_{\mathrm{small}}1.03`$. This indicates that the systems (and measurements of system quantities) are sensitive to a typical length scale, which is here somewhat smaller than the mean particle size. When using this special $`\mathrm{\Delta }r_{\mathrm{REV}}`$ value, one has $`B=20`$ or $`B=21`$ binning intervals. The open question of this being a typical length scale that also occurs in systems with a broader size spectrum, cannot be answered with our setup, due to the given particle-size ratio.
### 3.4 Time averaging
In order to understand the fluctuations in the system over time, and to test whether subsequent snapshots can be assumed to be independent, the volume fraction $`\nu `$ is displayed for snapshots at different averaging times in Fig. 4.
Changes in density are very weak and mostly occur in the dilated shear zone for small $`r`$. From one snapshot to the next, we frequently find, that the configuration in the outer part of the shear cell has not changed, whereas a new configuration is found in the inner part. Only after rather long times does the density change also in the outer part. Thus, simulation results in the outer part are subject to stronger fluctuations because the average is performed over less independent configurations than in the inner part.
## 4 Macroscopic tensorial quantities
In this section, the averaged, macroscopic tensorial quantities in our model system are presented. The fabric tensor describes the contact network, the stress tensor describes, in this study, the stress due to normal forces, and the deformation gradient is a measure for the corresponding elastic, reversible deformations. A more detailed analysis, where the tangential forces are also taken into account, is in progress latzel99a and some details in that direction can be found in Refs. kruyt96 ; liao97b ; kuhl99 . However, we checked the influence of the tangential forces in our system and found their effect to be always smaller than ten per-cent.
### 4.1 Micro-mechanical fabric tensor
In assemblies of grains, the forces are transmitted from one particle to the next only at the contacts of the particles. In the general case of non-spherical particles, a packing network is characterized by the vectors connecting the centers of the particles and by the particle-contact vectors. Furthermore, the local geometry of each contact is important goddard98 ; cowin88 , see Fig. 5. In our case, with spherical particles, the situation is simpler with respect to both the spherical contact geometry and the fabric.
#### 4.1.1 Fabric tensor for one particle
One quantity that describes the contact configuration of one particle to some extent, is the second order fabric tensor goddard98 ; cowin88
$$𝐅^p=\underset{c=1}{\overset{𝒞^p}{}}\stackrel{}{n}^c\stackrel{}{n}^c,$$
(12)
where $`\stackrel{}{n}^c`$ is the unit normal vector at contact $`c`$ of particle $`p`$ with $`𝒞^p`$ contacts. The symbol $``$ denotes the dyadic product in this study. Other definitions of the fabric use the so-called branch vector $`\stackrel{}{l}^{pc}`$ from the center of particle $`p`$ to its contact $`c`$, however, the unit normal and the unit branch vector are related by $`a_p\stackrel{}{n}=\stackrel{}{l}^{pc}`$ in the case of spheres or disks, so that the definition
$$𝐅^p=\frac{1}{a_p^2}\underset{c=1}{\overset{𝒞^p}{}}\stackrel{}{l}^{pc}\stackrel{}{l}^{pc}.$$
(13)
is identical to Eq. (12).
The fabric tensor in Eq. (12) is symmetric by definition and thus consists of up to three independent scalar quantities in two dimensions. The first of them, the trace (or volumetric part) $`F_V=\mathrm{tr}(𝐅^p)=(F_{\mathrm{max}}+F_{\mathrm{min}})`$, is the number of contacts of particle $`p`$, with the major and the minor eigenvalues $`F_{\mathrm{max}}`$ and $`F_{\mathrm{min}}`$, respectively. One gets from Eqs. (12) or (13) the number of contacts of particle $`p`$
$$\mathrm{tr}(𝐅^p)=\underset{c=1}{\overset{𝒞^p}{}}\mathrm{tr}(\stackrel{}{n}^c\stackrel{}{n}^c)=𝒞^p,$$
(14)
since the scalar product of $`\stackrel{}{n}^c`$ with itself is unity by definition. The second scalar, the deviator $`F_D=F_{\mathrm{max}}F_{\mathrm{min}}`$, accounts for the anisotropy of the contact network to first order, and the third, the angle $`\varphi _F`$, gives the orientation of the “major eigenvector”, i.e. the eigenvector corresponding to $`F_{\mathrm{max}}`$, with respect to the radial direction. In other words, the contact probability distribution is proportional to the function $`[F_V+F_D\mathrm{cos}(2(\varphi \varphi _F))]`$ tsoungui98b ; dubujet98 ; schollmann99 , when averaged over many particles, an approximation which is not always reasonable mehrabadi88 .
#### 4.1.2 Averaged fabric tensor
Assuming that all particles in $`V`$ contribute to the fabric with a weight of their volume $`V^p`$ one has
$$𝐅=𝐅^p=\frac{1}{V}\underset{pV}{}w_V^pV^p\underset{c=1}{\overset{𝒞^p}{}}\stackrel{}{n}^c\stackrel{}{n}^c.$$
(15)
The values of $`F_V`$, $`F_D/F_V`$, and $`\varphi _F`$, as obtained from our simulations, are plotted in Figs. 6(a-c). The trace of the fabric tensor and thus the mean number of contacts increases with increasing distance from the inner ring, and is reduced in the vicinity of the walls. With increasing mean density, the trace of $`𝐅`$ is systematically increasing, while the deviatoric fraction seems to decrease; this means that a denser system is slightly more isotropic concerning the fabric. The major eigendirection is tilted counterclockwise by somewhat more than $`\pi /4`$ from the radial outward direction, except for the innermost layer and for the strongly fluctuating outer region.
In analogy to the trace of the fabric for a single particle, the trace of the averaged fabric is
$$\mathrm{tr}(𝐅)=\frac{1}{V}\underset{pV}{}w_V^pV^p𝒞^p=𝒞^p,$$
(16)
which, in the case of a regular, periodic contact network of almost identical particles with $`a_pa`$, reduces to the sum over all of particles in $`V`$ with the prefactor $`\nu `$ defined in Eq. (8). Now, one has a relation between the coordination number, i.e. the mean number of contacts per particle $`𝒞=𝒞(r)=𝒞^p/\nu `$, the volume fraction $`\nu `$, and the averaged fabric $`𝐅`$, as a combination of Eqs. (8) and (16):
$$\mathrm{tr}(𝐅)\nu 𝒞.$$
(17)
As a test for the averaging procedure, we plot in Fig. 7 $`\mathrm{tr}(𝐅)`$ against $`\nu 𝒞`$ and obtain all data points from all simulations collapsing close to the identity curve. For a theoretical derivation of the small (about 1 per-cent) deviation due to the polydispersity, see Ref. tsoungui99 .
### 4.2 Micro-mechanical stress tensor
The micro-mechanical stress tensor is derived in a way similar to Ref. kruyt96 . For an arbitrary volume $`V`$ with surface $`V`$, the mean stress is defined as
$$𝝈=\frac{1}{V}_V𝑑V^{}𝝈^{},$$
(18)
where $`𝝈^{}=𝝈^{}(\stackrel{}{x})`$ is the position dependent local stress inside $`V`$ and $`\stackrel{}{x}`$ is the coordinate integrated over. Note that $`𝝈^{}`$ can be a strongly fluctuating function of $`\stackrel{}{x}`$.
If the stress in the pore-space vanishes, the integral can be split into a sum over the stresses $`𝝈^p`$, pre-averaged for particles $`p`$, with the respective center of mass vectors $`\stackrel{}{x}^p`$. The mean stress is thus
$$𝝈=𝝈=\frac{1}{V}\underset{pV}{}w_V^pV^p𝝈^p=\frac{1}{V}\underset{pV}{}w_V^p_{V^p}𝑑V^{}𝝈^{},$$
(19)
which corresponds to a smearing out of $`𝝈^{}`$ over the particles, with $`𝝈^p`$, the average stress in particle $`p`$, to be derived in the following.
One can rewrite the transposed stress tensor so that
$$𝝈^\mathrm{T}=\mathrm{grad}\stackrel{}{\mathrm{x}}𝝈^\mathrm{T}=\mathrm{div}(\stackrel{}{\mathrm{x}}𝝈)\stackrel{}{\mathrm{x}}\mathrm{div}𝝈,$$
(20)
by introducing the unit tensor $`𝐈=\mathrm{grad}\stackrel{}{\mathrm{x}}`$, and by applying the series rule in the first term on the right hand side. Note that the tensors $`𝝈`$ on the right hand side are transposed with respect to the left hand side $`𝝈^\mathrm{T}`$. The transposed stress is used for the following operations for the sake of simplicity, however, using the stress directly should lead to the same result. For a more detailed treatment, see Ref. latzel99a . In static equilibrium and in the absence of body forces, the term $`\mathrm{div}𝝈`$ on the right hand side vanishes and $`𝝈`$ is symmetric.
#### 4.2.1 Mean stress for one particle
Using the definition of the transposed stress tensor in Eq. (20) the integral over one particle in Eq. (19) becomes
$`(𝝈^p)^\mathrm{T}`$ $`=`$ $`{\displaystyle \frac{1}{V^p}}{\displaystyle _{V^p}}𝑑V^{}\mathrm{div}(\stackrel{}{\mathrm{x}}𝝈^{})`$ (21)
$`=`$ $`{\displaystyle \frac{1}{V^p}}{\displaystyle _{V^p}}𝑑S(\stackrel{}{x}𝝈^{})\stackrel{}{n}`$
by application of Gauss’ theorem. In Eq. (21) $`dS`$ is the surface element of $`V^p`$ on $`V^p`$ and $`\stackrel{}{n}`$ is the outwards normal unit vector. Using the definition of the stress vector $`\stackrel{}{t}=𝝈^{}\stackrel{}{n}`$ one arrives at
$$(𝝈^p)^\mathrm{T}=\frac{1}{V^p}_{V^p}𝑑S\stackrel{}{x}\stackrel{}{t}.$$
(22)
A force $`\stackrel{}{f}^c`$ acting at a contact $`c`$ with area $`\delta s^c`$ leads to a stress vector $`\stackrel{}{t}^c=\stackrel{}{f}^c/\delta s^c`$ – even in the limit of small contact area $`\delta s^ca_p`$. Here, we apply the simplifying assumption that the force $`\stackrel{}{f}^c`$ is constant on the surface $`\delta s^c`$, i.e. we do not resolve any details at the contact. Explicitly writing the surface element $`dS`$ as a sum,
$$dS=\underset{c=1}{\overset{𝒞^p}{}}\delta (|\stackrel{}{x}\stackrel{}{x}^c|)\delta s^c,$$
(23)
leads, after integration over the delta functions, to the transposed stress
$$(𝝈^p)^\mathrm{T}=\frac{1}{V^p}\underset{c=1}{\overset{𝒞^p}{}}\stackrel{}{x}^c\stackrel{}{f}^c.$$
(24)
Transposing Eq. (24) leads to an exchange of $`\stackrel{}{x}^c`$ and $`\stackrel{}{f}^c`$, and thus to the expression for the mean stress inside particle $`p`$:
$$𝝈^p=\frac{1}{V^p}\underset{c=1}{\overset{𝒞^p}{}}\stackrel{}{f}^c\stackrel{}{x}^c,$$
(25)
which is, at first glance, dependent on the frame of reference associated with the vector $`\stackrel{}{x}^c`$.
Using vector addition, one has $`\stackrel{}{x}^c=\stackrel{}{x}^p+\stackrel{}{l}^{pc}`$, with the position vector $`\stackrel{}{x}^p`$ of particle $`p`$ and the branch vector $`\stackrel{}{l}^{pc}`$, pointing from the center of mass of particle $`p`$ to contact $`c`$. Inserting this relation in Eq. (25) leads to
$`𝝈^p`$ $`=`$ $`{\displaystyle \frac{1}{V^p}}\left({\displaystyle \underset{c=1}{\overset{𝒞^p}{}}}\stackrel{}{f}^c\right)\stackrel{}{x}^p+{\displaystyle \frac{1}{V^p}}{\displaystyle \underset{c=1}{\overset{𝒞^p}{}}}\stackrel{}{f}^c\stackrel{}{l}^{pc}`$ (26)
$`=`$ $`{\displaystyle \frac{1}{V^p}}{\displaystyle \underset{c=1}{\overset{𝒞^p}{}}}\stackrel{}{f}^c\stackrel{}{l}^{pc},`$
since the first sum vanishes in static equilibrium, where
$$\underset{c=1}{\overset{𝒞^p}{}}\stackrel{}{f}^c=\stackrel{}{0}.$$
(27)
#### 4.2.2 Averaged stress tensor
Inserting Eq. (26) in Eq. (19) gives a double sum over all particles with center inside the averaging volume $`V`$ and all their contacts
$$𝝈=𝝈^p=\frac{1}{V}\underset{pV}{}w_V^p\underset{c=1}{\overset{𝒞^p}{}}\stackrel{}{f}^c\stackrel{}{l}^{pc}.$$
(28)
In other words, stress is pre-averaged over all particles and then averaged over $`V`$.
Inserting Eq. (25) in Eq. (19) leads to the mathematically identical sum
$$𝝈=\frac{1}{V}\underset{pV}{}w_V^p\underset{c=1}{\overset{𝒞^p}{}}\stackrel{}{f}^c\stackrel{}{x}^c.$$
(29)
This expression can be transformed into a sum over all contacts $`c(p,q)V`$, where at least one of the participating particles $`p`$ and $`q`$ lies inside the volume $`V`$. The force $`\stackrel{}{f}^c=\stackrel{}{f}^{pq}=\stackrel{}{f}^{qp}`$, acting at contact $`c`$, from particle $`q`$ on $`p`$ is equal, but opposite in direction to the force acting at the same contact from particle $`p`$ on $`q`$, so that
$`𝝈`$ $`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \underset{cV}{}}(w_V^p\stackrel{}{f}^{pq}+w_V^q\stackrel{}{f}^{qp})\stackrel{}{x}^c`$ (30)
$`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \underset{cV}{}}(w_V^pw_V^q)\stackrel{}{f}^c\stackrel{}{x}^c,`$
since $`w_V^pw_V^q=0`$ when both particles are completely inside $`V`$. Note that every contact is visited once in Eq. (30) but twice in Eqs. (28) and (29). The expression $`cV`$ means for the particle-center averaging method that contacts contribute only if one particle is inside $`V`$ while the other is outside. In the framework of the slicing method, contacts contribute if at least one particle is cut by $`V`$ so that $`w_V^pw_V^q0`$.
Given an arbitrary averaging volume $`V`$, and all the information about the forces at all contacts from the discrete element simulations, it is obvious that – from the technical point of view – the sum in Eq. (28) can be done relatively easily, whereas the sum in Eq. (30) requires the identification of the contacts at the “surface” of $`V`$, before the possibly shorter summation can be performed. Note that both expressions Eq. (28) and Eq. (30) have been used by different authors, see e.g. kruyt96 ; chang88 ; cundall82 .
The values of the volumetric stress $`\sigma _V=\mathrm{tr}(𝝈)`$, the deviatoric fraction $`\sigma _D/\sigma _V=\mathrm{dev}(𝝈)/\mathrm{tr}(𝝈)`$, as obtained from our simulations, are plotted in Figs. 8(a-b). The volumetric stress is constant, besides fluctuations, whereas the deviatoric fraction is rather large in the shear-zone and decays with increasing distance from the inner ring, like the fabric. Also, the deviatoric fraction of the stress appears slightly reduced in magnitude for increasing global density. Note that the stress varies over almost two orders of magnitude, while the mean density is changed only weakly from $`\overline{\nu }=0.8084`$ to $`0.8194`$, see table 2. In Fig. 8c $`\varphi _\sigma `$ describes the angle which the principal axis of the stress tensor is tilted from the outward direction. The stress tensor is, on average, rotated counter clockwise by somewhat less than $`\pi /4`$ from the outward direction. It is only in the outermost part that strong fluctuations exist around the mean. We now provide evidence that the fabric and stress are not co-linear.
### 4.3 Mean strain tensor
To achieve the material properties of a granular ensemble one is interested in the stress-strain relationship of the material. One of the simplest techniques used, is the application of “Voigt’s hypothesis” which assumes that the strain is uniform and that every particle displacement conforms to the mean displacement field liao97 . Thus, the expected displacement at contact $`c`$ of particle $`p`$, relative to the force free situation, and due to the mean displacement gradient $`\mathit{ϵ}`$, is $`\mathit{ϵ}\stackrel{}{l}^{pc}`$, with $`\stackrel{}{l}^{pc}=a_p\stackrel{}{n}^c`$, and the mean contact deformation $`\delta `$. Note that the linear, symmetric strain $`𝜺=\frac{1}{2}(\mathit{ϵ}+\mathit{ϵ}^\mathrm{T})`$ is not identical to the displacement gradient, in general. In our study, we follow the approach of Liao et al. liao97b , who used $`\stackrel{}{l}^{pq}`$ instead of $`\stackrel{}{l}^{pc}`$, and assume that the actual displacement field does not coincide with, but fluctuates about the mean displacement field. The difference between the actual displacement $`\stackrel{}{\mathrm{\Delta }}^{pc}`$ and the expected displacement is
$$\stackrel{}{\chi }^{pc}=\mathit{ϵ}\stackrel{}{l}^{pc}\stackrel{}{\mathrm{\Delta }}^{pc}.$$
(31)
The actual displacement is directly related to the simulations via $`\stackrel{}{\mathrm{\Delta }}^{pc}=\delta ^c\stackrel{}{n}^c`$.
If one assumes that the mean displacement field best approximates the actual displacement, one can apply a “least square fit” to the total fluctuations
$$S=\frac{1}{V}\underset{pV}{}w_V^p\underset{c=1}{\overset{𝒞^p}{}}(\stackrel{}{\chi }^{pc})^2,$$
(32)
by minimizing $`S`$ with respect to the mean displacement gradient so that
$`{\displaystyle \frac{S}{\mathit{ϵ}}}`$ $`\stackrel{!}{=}`$ $`\mathrm{𝟎}`$
$`=`$ $`{\displaystyle \frac{2}{V}}{\displaystyle \underset{pV}{}}w_V^p{\displaystyle \underset{c=1}{\overset{𝒞^p}{}}}(\mathit{ϵ}\stackrel{}{l}^{pc}\stackrel{}{\mathrm{\Delta }}^{pc}){\displaystyle \frac{}{\mathit{ϵ}}}(\mathit{ϵ}\stackrel{}{l}^{pc}\stackrel{}{\mathrm{\Delta }}^{pc}).`$
These four equations for the four components of $`\mathit{ϵ}`$ in 2D can be transformed into a relation for the mean displacement tensor as a function of the contact displacements and the branch vectors. By assuming that $`\stackrel{}{\mathrm{\Delta }}^{pc}/\mathit{ϵ}=\stackrel{}{0}`$, the expression “$`\frac{}{\mathit{ϵ}}(\mathit{ϵ}\stackrel{}{l}^{pc})`$” becomes a dyadic product “$`\stackrel{}{l}^{pc}`$”, as can be seen by writing down the equation in index notation. Extracting $`\mathit{ϵ}`$ from the sum (what leaves $`\stackrel{}{l}^{pc}\stackrel{}{l}^{pc}`$, the core of the fabric, in the first term) and multiplying the equation with $`𝐀=𝐅^1`$, the inverse of the fabric, see Eq. (15), we find that
$$\mathit{ϵ}=\frac{\pi }{V}\left(\underset{pV}{}w_V^p\underset{c=1}{\overset{𝒞^p}{}}\stackrel{}{\mathrm{\Delta }}^{pc}\stackrel{}{l}^{pc}\right)𝐀.$$
(34)
The values of $`ϵ_V`$, $`ϵ_D/ϵ_V`$, and $`\varphi _ϵ`$, as obtained from our simulations, are plotted in Figs. 9(a-c). The elastic, volumetric deformation gradient of the granulate is localized in the shear zone and the effect is stronger in the less dense systems. Due to dilation it is easier to compress the dilute material closer to the inner ring compared to the outer part. Like the fabric and stress, the strain also becomes more isotropic with increasing mean density.
## 5 Results
At first, we compute the mean-field expectation values for $`𝝈`$ and $`\mathit{ϵ}`$, in order to get a rough estimate for the orders of magnitude of the following results. Replacing, in Eq. (28), $`f^c`$ by its mean $`\overline{f}=k_n\overline{\delta }`$, $`a_p`$ by $`a=\overline{a}`$, and $`\stackrel{}{l}^{pc}`$ by $`a\stackrel{}{n}^c`$, one gets
$$\overline{𝝈}=(k_n\overline{\delta }/\pi a)𝐅.$$
(35)
Performing some similar replacements in Eq. (34), leads to
$$\overline{\mathit{ϵ}}=(\overline{\delta }/a)𝐈,$$
(36)
or
$$\overline{\mathit{ϵ}}=(\pi /k_n)𝝈𝐀.$$
(37)
The material stiffness, $`\overline{E}`$, can be defined as the ratio of the volumetric parts of stress and strain, so that one obtains from Eq. (35) and (36)
$$\overline{E}=(k_n/2\pi )\mathrm{tr}(𝐅).$$
(38)
In Fig. 10 the rescaled stiffness of the granulate is plotted against the trace of the fabric for all simulations. Note that all data collapse almost on a line, but the mean-field value underestimates the simulation data by a few per-cent. Simulation data for different $`k_n`$ and even data from simulations with neither bottom- nor tangential friction collapse with the data for fixed $`k_n`$ and different volume fractions, shown here.
The deviatoric fraction of $`𝝈`$ is plotted in Fig. 11 against the deviatoric fraction of $`𝐅`$. Here, the data are strongly underestimated by the mean-field result in Eq. (35). A similar plot for the deviatoric fraction of the deformation gradient against $`F_D/F_V`$ shows a gathering of the data close to the identity curve, in disagreement with Eq. (36). Thus, we conclude that the deviatoric parts of stress, deformation gradient and fabric are interconnected in a more complicated way than suggested by the simple mean field estimates kruyt96 ; liao97 ; luding99b ; luding99c .
In Fig. 12 the ratio of the deviatoric parts of stress and strain is plotted against the trace of the fabric. Like the material stiffness, both quantities are proportional, only $`G`$ shows much stronger fluctuations and has a proportionality factor of about $`1/3`$. We did not use the traditional definition of the shear modulus kruyt96 , since our tensors are not co-linear as shown below.
In Fig. 13 the orientations of the tensors are plotted for some simulations and in the inner part of the shear cell. In the outer part, the deviatoric fraction is usually around 10 per-cent, i.e. so small that the orientations become noisier. Simulations A and B are skipped here, because the data of the orientations are rather noisy due to the low density which leads to intermittent behavior with strong fluctuations. We observe that all orientation angles $`\varphi `$ show the same qualitative behavior, however, the fabric is tilted more than the stress which in turn is tilted more than the deformation gradient, where the orientations are measured in counter clockwise direction from the radial outward direction.
As a final cross-check, inserting the measured values for $`𝝈`$ and $`𝐀`$ into Eq. (37) leads to the measured values of $`\mathit{ϵ}`$ within a few percent deviation. However, the orientation of the deformation gradient is not well reproduced – it seems very sensitive to small fluctuations.
## 6 Summary and Conclusion
Discrete element simulations of a 2D Couette shear cell were presented and used as the basis for a micro-macro averaging procedure. In the shear cell a shear band is localized close to the inner, rotating cylindrical wall. The boundary conditions were chosen to allow for averaging over large volumes (rings with width $`\mathrm{\Delta }r`$) and over a steady state and thus over long times. The configurations changed rather rapidly in the shear band, whereas the system is frozen in the outer part, a fact which requires either extremely long simulations or a sampling over different initial conditions in order to allow an ensemble averaging in the outer part.
The simplest averaging strategy involves only the particle-centers as carriers of the quantities to be averaged over, whereas a more advanced method assumes the quantities to be homogeneously smeared out over the whole particle which is cut in slices by the averaging volumes. Both methods agree if the averaging volumes are of the particle size (or multiples), but for other sizes differing results are obtained. The slicing method shows discretization effects in the range of averaging volume widths from one to one fifth of a particle diameter, while the particle-center method shows fluctuations due to the choice of the averaging volume in a much wider range.
The material density, i.e. the volume fraction, the coordination number, the fabric tensor, the stress tensor and the elastic, reversible deformation gradient were obtained by the averaging procedures. The fabric is linearly proportional to the product of volume fraction and coordination number. In the shear band, dilation together with a reduction of the number of contacts is observed. The mean stress is constant in radial direction while the deformation gradient decays with the distance from the inner wall. The ratio of the volumetric parts of stress and strain gives the effective stiffness of the granulate, which is small in the shear band and larger outside, due to dilation.
In the shear band, large deviatoric components of all tensorial quantities are found, however, decreasing with increasing distance from the inner wall. The isotropy of the tensors grows only slightly with increasing density and all tensors are tilted counterclockwise from the radial direction by an angle of the order of $`\pi /4`$. The system organizes itself such that more contacts are created to act against the shear. An essential observation is that the macroscopic tensors are not co-linear, i.e. their orientations are different. The orientation of the fabric is tilted most, that of the deformation gradient is tilted least and thus, the material cannot be described by a simple elastic model involving only two Lamé constants (or bulk modulus and Poisson’s ratio) as the only parameters. Alternatives are to cut the system into pieces with different material properties and thus introducing discontinuities herrmann98 or to use the rank four stiffness tensor for anisotropic materials luding99c . The deviatoric parts of stress and deformation gradient are seemingly interconnected via the fabric tensor.
To conclude, we proposed a consistent averaging formalism to obtain a mean quantity $`Q`$ in average over arbitrary volumes $`V`$. Within this framework, we used for $`Q^c`$ the fabric $`𝐅^c=\stackrel{}{n}^c\stackrel{}{n}^c`$, the stress $`𝝈^c=(1/V^p)\stackrel{}{f}^c\stackrel{}{l}^{pc}`$, and the deformation gradient $`\mathit{ϵ}^c=(\delta ^c/a_p)\stackrel{}{n}^c\stackrel{}{n}^c𝐀`$. Future work will involve the extension of the present analysis to measure also the fluctuations of the above quantities and, in addition, other interesting quantities like, e.g. , the plastic strain and non-symmetric parts of the stress tensor due to the effective moments acting on single particles. More systematic parameter studies are currently in progress.
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# Topological defect system in 𝑂(𝑛) symmetric time-dependent Ginzburg-Landau model
## I Introduction
The importance of the role of time-dependent Ginzburg-Landau (TDGL) model in understanding a variety of problems in physics is clear. The study of the growth kinetics of systems subjected to rapid temperature quenches has recently been extended to include systems with more complex order parameter symmetries. In particular there has been progress on the study of the $`n`$-dimensional vector model with nonconserved Langevin dynamics. An interface description and numerical simulations of a TDGL equation were used to investigate the intra-surface kinetics of phase ordering on toroidal and corrugated surfaces. Two-dimensional $`XY`$ models with resistively-shunted junction (RSJ) dynamics and TDGL dynamics were simulated and it was verified that the vortex responses can be well described by the Minnhagen phenomenology for both types of dynamics.
In recent years, much work has been done on the topological defects in the TDGL model. In certain cosmological and phase ordering problems key questions involve an understanding of the evolution and correlation among defects like vortices, monopoles, domain walls, etc. In studying such objects in field theory questions arise as to how one can define quantities like the density of topological defects and an associated velocity field. Liu and Mazenko have discussed the problem, but unfortunately, for the lack of a powerful method, the topological properties of these systems are not very clear, some important topological information has been lost, also the unified theory of describing the topological properties of all defect system in TDGL model is not established yet.
In our previous work, we discussed how one could use the $`\varphi `$-mapping topological current theory to study the topological structure of the point defects in phase-ordering systems. We deduced a topological current of point-defect system in terms of the order parameter field in the context of a $`d`$-dimensional $`O(n)`$ symmetric TDGL model, where $`d`$ is the spatial dimensionality. Using the expression of the topological current, for point defects ($`n=d`$), the point-defect velocity field was identified, the topological structure and the topological charges of the point-defect system were studied, also the branch process of the point defects was discussed systematically. This analysis will be extended here to the case of arbitrary dimensional defect system in TDGL model where $`n=dk`$, i.e. to study the case for all $`nd`$.
In this paper, in the light of $`\varphi `$–mapping topological current theory, a useful method which plays a important role in studying the topological invariants and the topological structures of physical systems, we will investigate the topological quantization and the branch process of arbitrary dimensional topological defects in $`O(n)`$ symmetric TDGL model. We will show that the topological defects are generated from the zero points of the order parameter field $`\stackrel{}{\varphi }`$, and their topological charges are quantized in terms of the Hopf indices and Brouwer degrees of $`\varphi `$-mapping under the condition that the zero points of field $`\stackrel{}{\varphi }`$ are regular points. While at the critical points of the order parameter field $`\stackrel{}{\varphi }`$, i.e. the limit points and bifurcation points, there exist branch processes, the topological current of defect bifurcates and the topological defects split or merge at such point, this means that the topological defects system is unstable at these points.
This paper is organized as follows. In section 2, we investigate the topological structure of the defect system in TDGL model and give the expression of the defect velocity field. In section 3, we will point out that the topological charges of these defects are the Winding numbers which are determined by the Hopf indices and the Brouwer degrees of the $`\varphi `$–mapping. In section 4, we study the branch process of the defect topological current at the limit points, bifurcation points and higher degenerated points systematically by virtue of the $`\varphi `$–mapping theory and the implicit function theorem.
## II Topological structure of the defect system in TDGl model
We consider a TDGL model for an $`n`$-component order parameter $`\stackrel{}{\varphi }=(\varphi ^1(\stackrel{}{r},t),\mathrm{},\varphi ^n(\stackrel{}{r},t))`$ in $`d`$-spatial dimensions ($`dn=k`$) governed by the Langevin equation
$$\frac{\stackrel{}{\varphi }}{t}=\stackrel{}{K}\mathrm{\Gamma }\frac{\delta F}{\delta \stackrel{}{\varphi }}+\stackrel{}{\eta }$$
(1)
where $`\mathrm{\Gamma }`$ is a kinetic coefficient and $`\stackrel{}{\eta }`$ is a thermal noise which is related to $`\mathrm{\Gamma }`$ by fluctuation-dissipation theorem. $`F`$ is a Ginzburg-Landau effective free energy assumed to be of the form
$$F=d^dr[\frac{1}{2}c(\stackrel{}{\varphi })^2+V(\stackrel{}{\varphi })]$$
(2)
where $`c>0`$ and the potential $`V`$ is assumed to be of the degenerate double-well form. This model is to be supplemented by random, uncorrelated, initial conditions. We assume that there is a rapid temperature quench from a high temperature to zero temperature where the noise $`\stackrel{}{\eta }`$ in (1) can be set to zero. In the scalar case ($`n=1`$) such system order through the growth of domains separated by sharp walls. As time evolves these domains coarsen and order grows to progressively longer length scales. In the case of systems with continuous symmetry ($`n>1`$) the disordering elements will depend on $`n`$ and spatial dimensionality $`d`$. Thus, for example, for $`n=d`$ one has point defects (vortices and monopoles),for $`n=d1`$ one has vortex line or string-like defects while for $`n=dk`$ one has $`k`$-dimensional topological defect objects. For $`n>d`$ there are no stable singular topological objects.
As is well known, the $`n`$-component order parameter field $`\stackrel{}{\varphi }(\stackrel{}{r},t)`$ determines the defect properties of the system, and it can be looked upon as a smooth mapping between the $`(d+1)`$-dimensional space-time $`X`$ and the $`n`$-dimensional Euclidean space $`R^n`$ as $`\varphi :XR^n`$. By analogy with the discussion in our previous work, from this $`\varphi `$-mapping, one can deduce a topological tensor current as
$`j^{\mu _0\mu _1\mathrm{}\mu _k}={\displaystyle \frac{1}{A(S^{n1})(n1)!}}ϵ^{\mu _0\mu _1\mathrm{}\mu _k\mu _{k+1}\mathrm{}\mu _d}ϵ_{a_1\mathrm{}a_n}_{\mu _{k+1}}n^{a_1}_{\mu _{k+2}}n^{a_2}\mathrm{}_{\mu _d}n^{a_n},`$
$$\mu _0,\mathrm{},\mu _d=0,1,\mathrm{},d,a_1,\mathrm{},a_n=1,\mathrm{},n$$
(3)
to describe the $`k`$-dimensional topological defects in TDGL model, and the density tensor of the topological defect system is defined as $`\rho ^{\mu _1\mathrm{}\mu _k}=j^{0\mu _1\mathrm{}\mu _k}`$. In this expression, $`_\mu `$ stands for $`/x^\mu `$, $`A(S^{n1})=2\pi ^{n/2}/\mathrm{\Gamma }(n/2)`$ is the area of $`(n1)`$–dimensional unit sphere $`S^{n1}`$ and $`n^a(x)`$ is the direction field of the $`n`$–component order parameter field $`\stackrel{}{\varphi }`$
$$n^a(x)=\frac{\varphi ^a(x)}{\varphi (x)},\varphi (x)=\sqrt{\varphi ^a(x)\varphi ^a(x)}$$
(4)
with $`n^a(x)n^a(x)=1`$. It is obviously that $`n^a(x)`$ is a section of the sphere bundle $`S(X)`$ and it can be looked upon as a map of $`X`$ onto a $`(n1)`$–dimensional unit sphere $`S^{n1}`$ in order parameter space. Clearly, the zero points of the order parameter field $`\stackrel{}{\varphi }(x)`$ are just the singular points of the unit vector $`n^a(x)`$. It is easy to see that $`j^{\mu _0\mathrm{}\mu _k}`$ are completely antisymmetric, and from the formulas above, we conclude that there exists a conservative equation of the topological tensor current in (3)
$`_{\mu _i}j^{\mu _0\mathrm{}\mu _k}=0,i=0,\mathrm{},k,`$
and from which we have
$`_t\rho ^{\mu _1\mathrm{}\mu _k}+_\mu j^{\mu \mu _1\mathrm{}\mu _k}=0`$
which is just the continuity equation satisfied by $`\rho ^{\mu _1\mathrm{}\mu _k}`$.
In the following, we will investigate the intrinsic structure of the generalized topological current $`j^{\mu _0\mu _1\mathrm{}\mu _k}`$ by making use of the $`\varphi `$–mapping method. From (4), we have
$`_\mu n^a={\displaystyle \frac{1}{\varphi }}_\mu \varphi ^a+\varphi ^a_\mu ({\displaystyle \frac{1}{\varphi }}),{\displaystyle \frac{}{\varphi ^a}}({\displaystyle \frac{1}{\varphi }})={\displaystyle \frac{\varphi ^a}{\varphi ^3}}`$
which should be looked upon as generalized functions. Due to these expressions the generalized topological current (3) can be rewritten as
$`j^{\mu _0\mu _1\mathrm{}\mu _k}`$ $`=`$ $`C_nϵ^{\mu _0\mu _1\mathrm{}\mu _k\mu _{k+1}\mathrm{}\mu _d}ϵ_{a_1\mathrm{}a_n}`$ (5)
$``$ $`_{\mu _{k+1}}\varphi ^a\mathrm{}_{\mu _d}\varphi ^{a_n}{\displaystyle \frac{}{\varphi ^a}}{\displaystyle \frac{}{\varphi ^{a_1}}}(G_n(\varphi )),`$ (6)
where $`C_m`$ is a constant
$`C_n=\{\begin{array}{cc}\frac{1}{A(S^{n1})(n2)(n1)!},& n>2\\ \frac{1}{2\pi },& n=2\end{array},`$
and $`G_n(\varphi )`$ is a Green function
$`G_n(\varphi )=\{\begin{array}{ccc}\frac{1}{\varphi ^{n2}}& ,& n>2\\ \mathrm{ln}\varphi & ,& n=2\end{array}.`$
Defining general Jacobians $`J^{\mu _0\mu _1\mathrm{}\mu _k}(\frac{\varphi }{x})`$ as following
$`ϵ^{a_1\mathrm{}a_n}J^{\mu _0\mu _1\mathrm{}\mu _k}({\displaystyle \frac{\varphi }{x}})=ϵ^{\mu _0\mu _1\mathrm{}\mu _k\mu _{k+1}\mathrm{}\mu _d}_{\mu _{k+1}}\varphi ^{a_1}_{\mu _{k+2}}\varphi ^{a_2}\mathrm{}_{\mu _d}\varphi ^{a_n}`$
and by making use of the $`n`$–dimensional Laplacian Green function relation in $`\varphi `$–space
$`\mathrm{\Delta }_\varphi (G_n(\varphi ))={\displaystyle \frac{4\pi ^{n/2}}{\mathrm{\Gamma }(\frac{n}{2}1)}}\delta (\stackrel{}{\varphi })`$
where $`\mathrm{\Delta }_\varphi =(\frac{^2}{\varphi ^a\varphi ^a})`$ is the $`n`$–dimensional Laplacian operator in $`\varphi `$–space, we do obtain the $`\delta `$–function structure of the defect topological current rigorously
$$j^{\mu _0\mu _1\mathrm{}\mu _k}=\delta (\stackrel{}{\varphi })J^{\mu _0\mu _1\mathrm{}\mu _k}(\frac{\varphi }{x}).$$
(7)
This expression involves the total defect information of the system in TDGL model and it indicates that all the defects are located at the zero points of the order parameter field $`\stackrel{}{\varphi }(x)`$. It must be pointed out that, comparing to similar expressions in other papers, the results in (7) is gotten theoretically in a natural way. From this expression, the density tensor of $`j^{\mu _0\mu _1\mathrm{}\mu _k}`$ is also changed into a compact form
$`\rho ^{\mu _1\mathrm{}\mu _k}=j^{0\mu _1\mathrm{}\mu _k}=\delta (\stackrel{}{\varphi })D^{\mu _1\mathrm{}\mu _k}({\displaystyle \frac{\varphi }{x}}),`$
where $`D^{\mu _1\mathrm{}\mu _k}(\frac{\varphi }{x})=J^{0\mu _1\mathrm{}\mu _k}(\frac{\varphi }{x})`$. We find that $`j^{\mu _0\mu _1\mathrm{}\mu _k}0`$, $`\rho ^{\mu _1\mathrm{}\mu _k}0`$ only when $`\stackrel{}{\varphi }=0`$, which is just the singularity of $`j^{\mu _0\mu _1\mathrm{}\mu _k}`$. In detail, the Kernel of the $`\varphi `$–mapping is the singularities of the topological tensor current $`j^{\mu _0\mu _1\mathrm{}\mu _k}`$ in $`X`$, i.e. the inner structure of the topological tensor current is labeled by the zeroes of $`\varphi `$-mapping. We think that this is the essential of the topological tensor current theory and $`\varphi `$–mapping is the key to study this theory.
From the above discussions, we see that the kernel of $`\varphi `$–mapping plays an important role in the topological tensor current theory, so we are focused on the zero points of $`\stackrel{}{\varphi }`$ and will search for the solutions of the equations $`\varphi ^a(x)=0`$ $`(a=1,\mathrm{},n)`$ by means of the implicit function theorem. These points are topological singularities in the orientation of the order parameter field $`\stackrel{}{\varphi }(x)`$.
Suppose that the vector field $`\stackrel{}{\varphi }(x)`$ possesses $`l`$ zeroes, according to the implicit function theorem, when the zeroes are regular points of $`\varphi `$–mapping at which the rank of the Jacobian matrix $`[_\mu \varphi ^a]`$ is $`n`$, the solutions of $`\stackrel{}{\varphi }=0`$ can be expressed parameterizedly by
$$x^\mu =z_i^\mu (t,u^1,\mathrm{},u^k),i=1,\mathrm{},l,$$
(8)
where the subscript $`i`$ represents the $`i`$–th solution and the parameters $`u^I`$ ($`I=1,\mathrm{},k`$), combining with the time parameter $`t`$, span a $`(k+1)`$–dimensional subspace which is called the $`i`$–th singular subspace $`N_i`$ in the space-time $`X`$ corresponding to the $`\varphi `$–mapping. These singular subspaces $`N_i`$ are just the world volumes of the topological defects, the parameters $`u^I`$ play the role of the spatial parameters of the topological defects. For each singular subspace $`N_i`$, we can define a normal subspace $`M_i`$ in $`X`$ which is spanned by the parameters $`v^A`$ ($`A=1,\mathrm{},n`$), and the intersection point of $`M_i`$ and $`N_i`$ is denoted by $`p_i`$ which can be expressed parameterizedly by $`v^A=p_i^A`$. In fact, in the words of differential topology, $`M_i`$ is transversal to $`N_i`$ at the point $`p_i`$. By virtue of the implicit function theorem at the regular point $`p_i`$, it should be held true that the Jacobian matrices $`J(\frac{\varphi }{v})`$ satisfies $`J(\frac{\varphi }{v})=\frac{D(\varphi ^1,\mathrm{},\varphi ^n)}{D(v^1,\mathrm{},v^n)}0`$.
On the other hand, putting the solutions (8) into $`\stackrel{}{\varphi }(x)`$, we have $`\stackrel{}{\varphi }(z_i)0`$, from this expression, and since we expect the instantaneous velocity to be orthogonal to the local orientation of the topological defects, we can define the velocity via
$$J^{\mu \mu _1\mathrm{}\mu _k}=v^{[\mu }D^{\mu _1\mathrm{}\mu _k]},$$
(9)
that is, the topological current and the density tensor should satisfy
$$j^{\mu \mu _1\mathrm{}\mu _k}=v^{[\mu }\rho ^{\mu _1\mathrm{}\mu _k]}.$$
(10)
This expression for the velocity can be used to find the defect velocity distribution in the case of phase-ordering kinetics for a non-conserved order parameter. It is very useful because it avoid the problem of having to specify the positions of the topological defects explicitly. The positions are implicitly determined by the zeros of the order parameter field. The general expression with $`J^{\mu \mu _1\mathrm{}\mu _k}(\varphi /x)`$ should be useful in looking at the motion of the topological defects in TDGL model in the presence of external fields beyond a growth kinetics context.
If we restrict the discussion to the simplest case of point defects ($`n=d`$), the corresponding topological current is reduced to $`j^\mu =\delta (\varphi )J^\mu (\varphi /x)`$, and Eq.(10) can be put into a conventional form $`j^\mu =v^\mu \rho `$ with $`v^\mu `$ taking the expression of $`v^\mu =J^\mu (\varphi /x)/D(\varphi /x)`$. Also, the string in TDGl model, which was discussed by Mazenko, is a special case for $`n=d1`$ in our theory.
## III Topological quantization of the defect charges in TDGL model
In the following, we will investigate the topological charges of the topological defects and their quantization. Let $`\mathrm{\Sigma }_i`$ be a neighborhood of $`p_i`$ on $`M_i`$ with boundary $`\mathrm{\Sigma }_i`$ satisfying $`p_i\mathrm{\Sigma }_i`$, $`\mathrm{\Sigma }_i\mathrm{\Sigma }_j=\mathrm{}`$. Then the generalized winding number $`W_i`$ of $`n^a(x)`$ at $`p_i`$ can be defined by the Gauss map $`n:\mathrm{\Sigma }_iS^{n1}`$
$$W_i=\frac{1}{A(S^{n1})(n1)!}_{\mathrm{\Sigma }_i}n^{}(ϵ_{a_1\mathrm{}a_n}n^{a_1}dn^{a_2}\mathrm{}dn^{a_n})$$
(11)
where $`n^{}`$ denotes the pull back of map $`n`$. The generalized winding numbers is a topological invariant and is also called the degree of Gauss map. It means that, when the point $`v^A`$ covers $`\mathrm{\Sigma }_i`$ once, the unit vector $`n^a`$ will cover a region $`n[\mathrm{\Sigma }_i]`$ whose area is $`W_i`$ times of $`A(S^{n1})`$, i.e. the unit vector $`n^a`$ will cover the unit sphere $`S^{n1}`$ for $`W_i`$ times. Using the Stokes’ theorem in exterior differential form and duplicating the above process, we get the compact form of $`W_i`$
$$W_i=_{\mathrm{\Sigma }_i}\delta (\stackrel{}{\varphi })J(\frac{\varphi }{v})d^nv.$$
(12)
By analogy with the procedure of deducing $`\delta (f(x))`$, since
$$\delta (\stackrel{}{\varphi })=\{\begin{array}{cc}+\mathrm{},& for\stackrel{}{\varphi }(x)=0\\ 0,& for\stackrel{}{\varphi }(x)0\end{array}=\{\begin{array}{cc}+\mathrm{},& forxN_i\\ 0,& forxN_i\end{array},$$
(13)
we can expand the $`\delta `$–function $`\delta (\stackrel{}{\varphi })`$ as
$$\delta (\stackrel{}{\varphi })=\underset{i=1}{\overset{l}{}}c_i\delta (N_i),$$
(14)
where the coefficients $`c_i`$ must be positive, i.e. $`c_i=c_i`$. $`\delta (N_i)`$ is the $`\delta `$–function in space-time $`X`$ on a submanifold $`N_i`$
$$\delta (N_i)=_{N_i}\delta ^n(\stackrel{}{x}\stackrel{}{z}_i(t,u^1,\mathrm{},u^k))𝑑td^ku.$$
(15)
Substituting (14) into (12), and calculating the integral, we get the expression of $`c_i`$
$$c_i=\frac{\beta _i}{J(\frac{\varphi }{v})_{p_i}}=\frac{\beta _i\eta _i}{J(\frac{\varphi }{v})_{p_i}},$$
(16)
where $`\beta _i=|W_i|`$ is a positive integer called the Hopf index of $`\varphi `$-mapping on $`M_i,`$ it means that when the point $`v`$ covers the neighborhood of the zero point $`p_i`$ once, the function $`\stackrel{}{\varphi }`$ covers the corresponding region in $`\stackrel{}{\varphi }`$-space $`\beta _i`$ times, and $`\eta _i=signJ(\frac{\varphi }{v})_{p_i}=\pm 1`$ is the Brouwer degree of $`\varphi `$-mapping. Substituting this expression of $`c_i`$ and (14) into (7), we gain the total expansion of the topological current
$`j^{\mu _0\mu _1\mathrm{}\mu _k}={\displaystyle \underset{i=1}{\overset{l}{}}}{\displaystyle \frac{\beta _i\eta _i}{J(\frac{\varphi }{v})|_{p_i}}}\delta (N_i)J^{\mu _0\mu _1\mathrm{}\mu _k}({\displaystyle \frac{\varphi }{x}}).`$
or in terms of parameters $`y^A^{^{}}=(t,v^1,\mathrm{},v^n,u^1,\mathrm{},u^k)`$
$$j^{A_0^{^{}}A_1^{^{}}\mathrm{}A_k^{^{}}}=\underset{i=1}{\overset{l}{}}\frac{\beta _i\eta _i}{J(\frac{\varphi }{v})|_{p_i}}\delta (N_i)J^{A_0^{^{}}A_1^{^{}}\mathrm{}A_k^{^{}}}(\frac{\varphi }{y}).$$
(17)
From the above equation, we conclude that the inner structure of $`j^{\mu _0\mu _1\mathrm{}\mu _k}`$ or $`j^{A_0^{^{}}A_1^{^{}}\mathrm{}A_k^{^{}}}`$ is labeled by the total expansion of $`\delta (\stackrel{}{\varphi })`$, and it just represents $`l`$ $`(k)`$-dimensional topological defects with topological charges $`g_i=\beta _i\eta _i`$ moving in the $`(d+1)`$–dimensional space-time $`X`$. The $`(k+1)`$-dimensional singular subspaces $`N_i(i=1,\mathrm{}l)`$ are their world sheets in the space-time. Mazenko and Halperin also got similar results for the case of point-like defects and line defects, but unfortunately, they did not consider the case $`\beta _i1`$. In fact, what they lost sight of is just the most important topological information for the charge of topological defects. In detail, the Hopf indices $`\beta _i`$ characterize the absolute values of the topological charges of these defects and the Brouwer degrees $`\eta _i=+1`$ correspond to defects while $`\eta _i=1`$ to antidefects. Furthermore, they did not discuss what will happen when $`\eta _i`$ is indefinite, which we will study in detail in section 4.
## IV The branch process of the topological current
With the discussion mentioned above, we know that the results in the above section are obtained straightly from the topological view point under the condition $`J(\varphi /v)|_{p_i}0`$, i.e. at the regular points of the order parameter field $`\stackrel{}{\varphi }`$. When the condition fails, i.e. the Brouwer degree $`\eta _i`$ are indefinite, there should exist some kind of branch processes in the topological current of the topological defect system in TDGL model. In what follows, we will study the case when $`J(\varphi /v)|_{p_i}=0`$. It often happens when the zero points of field $`\stackrel{}{\varphi }`$ include some branch points, which lead to the bifurcation of the topological current.
In this section, we will discuss the branch processes of these topological defects. In order to simplify our study, let the spatial parameters $`u^I`$ be fixed, i.e. to choose a fixed point on the topological defect. In this case, the Jacobian matrices $`J^{A_0^{^{}}A_1^{^{}}\mathrm{}A_k^{^{}}}(\frac{\varphi }{y})`$ are reduced to
$`J^{AI_1\mathrm{}I_k}({\displaystyle \frac{\varphi }{y}})J^A({\displaystyle \frac{\varphi }{y}}),J^{ABI_1\mathrm{}I_{k1}}({\displaystyle \frac{\varphi }{y}})=0,J^{(n+1)\mathrm{}d}({\displaystyle \frac{\varphi }{y}})=J({\displaystyle \frac{\varphi }{v}}),`$
$$A,B=0,1,\mathrm{},n,I_j=n+1,\mathrm{},d,$$
(18)
for $`y^A=v^A(An),y^0=t,y^{n+I}=u^I(I1)`$. The branch points are determined by the $`n+1`$ equations
$$\varphi ^a(t,v^1,\mathrm{},v^n,\stackrel{}{u})=0,a=1,\mathrm{},n$$
(19)
and
$$\varphi ^{n+1}(t,v^1,\mathrm{},v^n,\stackrel{}{u})J(\frac{\varphi }{v})=0$$
(20)
for the fixed $`\stackrel{}{u}`$. and they are denoted as $`(t^{},p_i)`$. In $`\varphi `$–mapping theory usually there are two kinds of branch points, namely the limit points and bifurcation points, satisfying
$$J^1(\frac{\varphi }{y})|_{(t^{},p_i)}0$$
(21)
and
$$J^1(\frac{\varphi }{y})|_{(t^{},p_i)}=0,$$
(22)
respectively. In the following, we assume that the branch points $`(t^{},p_i)`$ of $`\varphi `$–mapping have been found.
### A The branch process at the limit point
We first discuss the branch process at the limit point satisfying the condition (21). In order to use the theorem of implicit function to study the branch process of topological defects at the limit point, we use the Jacobian $`J^1(\frac{\varphi }{y})`$ instead of $`J(\frac{\varphi }{v})`$ to discuss the problem. In fact, this means that we have replaced the parameter $`t`$ by $`v^1`$. Then, taking account of the condition(21) and using the implicit function theorem, we have an unique solution of the equations (19) in the neighborhood of the limit point $`(t^{},p_i)`$
$$t=t(v^1,\stackrel{}{u}),v^i=v^i(v^1,\stackrel{}{u}),i=2,3,\mathrm{},n$$
(23)
with $`t^{}=t(p_i^1,\stackrel{}{u})`$. In order to show the behavior of the defects at the limit points, we will investigate the Taylor expansion of (23) in the neighborhood of $`(t^{},p_i)`$. In the present case, from (21) and (20), we get
$`{\displaystyle \frac{dv^1}{dt}}|_{(t^{},p_i)}={\displaystyle \frac{J^1(\frac{\varphi }{y})}{J(\frac{\varphi }{v})}}|_{(t^{},p_i)}=\mathrm{},`$
i.e.
$`{\displaystyle \frac{dt}{dv^1}}|_{(t^{},p_i)}=0.`$
Therefore, the Taylor expansion of (23) at the point $`(t^{},p_i)`$ gives
$$tt^{}=\frac{1}{2}\frac{d^2t}{(dv^1)^2}|_{(t^{},p_i)}(v^1p_i^1)^2$$
(24)
which is a parabola in the $`v^1`$$`t`$ plane. From (24), we can obtain the two solutions $`v_{(1)}^1(t,\stackrel{}{u})`$ and $`v_{(2)}^1(t,\stackrel{}{u})`$, which give the branch solutions of the system (19) at the limit point. If $`\frac{d^2t}{(dv^1)^2}|_{(t^{},p_i)}>0`$, we have the branch solutions for $`t>t^{}`$ (Fig 1(a)), otherwise, we have the branch solutions for $`t<t^{}`$ (Fig 1(b)). The former is related to the creation of defect and antidefect in pair at the limit points, and the latter to the annihilation of the topological defects, since the topological current of topological defects is identically conserved, the topological quantum numbers of these two generated topological defects must be opposite at the limit point, i.e. $`\beta _1\eta _1+\beta _2\eta _2=0`$.
### B The branch process at the bifurcation point
In the following, let us consider the case (22), in which the restrictions of the system (19) at the bifurcation point $`(t^{},p_i)`$ are
$$J(\frac{\varphi }{v})|_{(t^{},p_i)}=0,J^1(\frac{\varphi }{y})|_{(t^{},p_i)}=0.$$
(25)
These two restrictive conditions will lead to an important fact that the dependency relationship between $`t`$ and $`v^1`$ is not unique in the neighborhood of the bifurcation point $`(t^{},p_i).`$ In fact, we have
$$\frac{dv^1}{dt}|_{(t^{},p_i)}=\frac{J^1(\frac{\varphi }{y})}{J(\frac{\varphi }{v})}|_{(t^{},p_i)}$$
(26)
which under the restraint (25) directly shows that the tangential direction of the integral curve of equation (26) is indefinite at the point $`(t^{},p_i)`$. Hence, (26) does not satisfy the conditions of the existence and uniqueness theorem of the solution of a differential equation. This is why the very point $`(t^{},\stackrel{}{z}_i)`$ is called the bifurcation point of the system (19).
As we have mentioned above, at the bifurcation point $`(t^{},p_i)`$, the rank of the Jacobian matrix $`[\frac{\varphi }{v}]`$ is smaller than $`n`$. For the aim of searching for the different directions of all branch curves at the bifurcation point, we firstly consider the rank of the Jacobian matrix $`[\frac{\varphi }{v}]`$ is $`n1`$. The case of a more smaller rank will be discussed in next subsection. Let $`J_1(\frac{\varphi }{v})=[\varphi _A^a](a=1,\mathrm{},n1;A=2,\mathrm{},n)`$ be one of the $`(n1)\times (n1)`$ submatrix of the Jacobian matrix $`[\frac{\varphi }{v}]`$ with $`detJ_1(\frac{\varphi }{v})0`$ at the point $`(t^{},p_i)`$ (otherwise, we have to rearrange the equations of (19)), where $`\varphi _A^a`$ stands for $`(\varphi ^a/v^A)`$. By means of the implicit function theorem we obtain one and only one functional relationship in the neighborhood of the bifurcation point $`(t^{},p_i)`$
$$v^A=f^A(v^1,t,\stackrel{}{u}),A=2,3,\mathrm{},n.$$
(27)
We denote the partial derivatives as $`f_1^A=\frac{v^A}{v^1}`$, $`f_t^A=\frac{v^A}{t}`$, $`f_{11}^A=\frac{^2v^A}{(v^1)^2}`$, $`f_{1t}^A=\frac{^2v^A}{v^1t}`$, $`f_{tt}^A=\frac{^2v^A}{t^2}`$. From (19) and (27), we have for $`a=1,\mathrm{},n1`$
$`\varphi ^a=\varphi ^a(v^1,f^2(v^1,t,\stackrel{}{u}),\mathrm{},f^n(v^1,t,\stackrel{}{u}),t,\stackrel{}{u})0`$
which gives
$$\underset{A=2}{\overset{n}{}}\frac{\varphi ^a}{v^A}f_1^A=\frac{\varphi ^a}{v^1},a=1,\mathrm{},n1$$
(28)
$$\underset{A=2}{\overset{n}{}}\frac{\varphi ^a}{v^A}f_t^A=\frac{\varphi ^a}{t},a=1,\mathrm{},n1.$$
(29)
By differentiating (28) and (29) with respect to $`v^1`$ and $`t`$, and applying the Gaussian elimination method, we can find the second order derivatives $`f_{11}^A`$, $`f_{1t}^A`$ and $`f_{tt}^A`$. The above discussions do not matter to the last component $`\varphi ^n(v^1,\mathrm{},v^n,t,\stackrel{}{u})`$. In order to find the different values of $`dv^1/dt`$ at the bifurcation point, let us investigate the Taylor expansion of $`\varphi ^n(v^1,\mathrm{},v^n,t,\stackrel{}{u})`$ in the neighborhood of $`(t^{},p_i)`$. Substituting (27) into $`\varphi ^n(v^1,\mathrm{},v^n,t,\stackrel{}{u})`$, we get the function of two variables $`v^1`$ and $`t`$
$$F(t,v^1,\stackrel{}{u})=\varphi ^m(v^1,f^2(v^1,t,\stackrel{}{u}),\mathrm{},f^m(v^1,t,\stackrel{}{u}),t,\stackrel{}{u})$$
(30)
which according to (19) must vanish at the bifurcation point
$$F(t^{},p_i)=0.$$
(31)
From (30), we can calculate the first order partial derivatives of $`F(t,v^1,\stackrel{}{u})`$ with respect to $`v^1`$ and $`t`$ respectively at the bifurcation point $`(t^{},p_i)`$
$$\frac{F}{v^1}=\varphi _1^n+\underset{A=2}{\overset{n}{}}\varphi _A^nf_1^A,\frac{F}{t}=\varphi _t^n+\underset{A=2}{\overset{n}{}}\varphi _A^nf_t^A.$$
(32)
By making use of (28) and (29), with the Cramer’s rule, the first equation of (25) is expressed as
$`{\displaystyle \frac{F}{v^1}}detJ_1({\displaystyle \frac{\varphi }{v}})|_{(t^{},p_i)}=0.`$
Since $`detJ_1(\frac{\varphi }{v})|_{(t^{},p_i)}0`$, the above equation leads to
$$\frac{F}{v^1}|_{(t^{},p_i)}=0.$$
(33)
With the same reasons, we can prove that
$$\frac{F}{t}|_{(t^{},p_i)}=0.$$
(34)
The second order partial derivatives of the function $`F(t,v^1,\stackrel{}{u})`$ are easily to find out from (32) which at $`(t^{},p_i)`$ are denoted by
$$A=\frac{^2F}{(v^1)^2}_{(t^{},p_i)},B=\frac{^2F}{v^1t}_{(t^{},p_i)},C=\frac{^2F}{t^2}_{(t^{},p_i)}.$$
(35)
Then, by virtue of (31), (33), (34) and (35), the Taylor expansion of $`F(t,v^1,\stackrel{}{u})`$ in the neighborhood of the bifurcation point $`(t^{},p_i)`$ gives
$$A(v^1p_i^1)^2+2B(v^1p_i^1)(tt^{})+C(tt^{})^2=0.$$
(36)
Dividing (36) by $`(v^1p_i^1)^2`$ or $`(tt^{})^2`$, and taking the limit $`tt^{}`$ as well as $`v^1p_i^1`$ respectively, we get two equations
$$A(\frac{dv^1}{dt})^2+2B\frac{dv^1}{dt}+C=0.$$
(37)
and
$$C(\frac{dt}{dv^1})^2+2B\frac{dt}{dv^1}+A=0.$$
(38)
So we get the different directions of the branch curves at the bifurcation point from the solutions of (37) or (38). There are four possible cases:
Firstly, $`A0,`$ $`\mathrm{\Delta }=4(B^2AC)>0`$, from Eq. (37) we get two different solutions: $`dv^1/dt_{1,2}=(B\pm \sqrt{B^2AC})/A`$, which is shown in Fig. 2, where two topological defects meet and then depart at the bifurcation point. Secondly, $`A0,\mathrm{\Delta }=4(B^2AC)=0`$, there is only one solution: $`dv^1/dt=B/A`$, which includes three important cases: (a) two topological defects tangentially collide at the bifurcation point (Fig 3(a)); (b) two topological defects merge into one topological defect at the bifurcation point (Fig 3(b)); (c) one topological defect splits into two topological defects at the bifurcation point (Fig 3(c)). Thirdly, $`A=0,C0,`$ $`\mathrm{\Delta }=4(B^2AC)>0`$, from Eq. (38) we have $`dt/dv^1=0`$ and $`2B/C`$. There are two important cases: (i) One topological defect splits into three topological defects at the bifurcation point (Fig 4(a)); (ii) Three topological defects merge into one at the bifurcation point (Fig 4(b)). Finally, $`A=C=0`$, Eqs. (37) and (38) give respectively $`dv^1/dt=0`$ and $`dt/dv^1=0`$. This case is obvious as in Fig. 5, which is similar to the third situation.
In order to determine the branches directions of the remainder variables, we will use the relations simply
$`dv^A=f_1^Adv^1+f_t^Adt,A=2,3,\mathrm{},n`$
where the partial derivative coefficients $`f_1^A`$ and $`f_t^A`$ have given in (28) and (29). Then, respectively
$`{\displaystyle \frac{dv^A}{dv^1}}=f_1^A+f_t^A{\displaystyle \frac{dt}{dv^1}}`$
or
$$\frac{dv^A}{dt}=f_1^A\frac{dv^1}{dt}+f_t^A.$$
(39)
where partial derivative coefficients $`f_1^A`$ and $`f_t^A`$ are given by (28) and (29). From this relations we find that the values of $`dv^A/dt`$ at the bifurcation point $`(t^{},z_i)`$ are also possibly different because (38) may give different values of $`dv^1/dt`$.
### C The branch process at the higher degenerated point
In the following, let us discuss the branch process at a higher degenerated point. In the above subsection, we have analyzed the case that the rank of the Jacobian matrix $`[\varphi /v]`$ of the equation (20) is $`n1`$. In this section, we consider the case that the rank of the Jacobian matrix is $`n2`$ (for the case that the rank of the matrix $`[\varphi /v]`$ is lower than $`n2`$, the discussion is in the same way). Let the $`(n2)\times (n2)`$ submatrix $`J_2(\frac{\varphi }{v})`$ of the Jacobian matrix $`[\varphi /v]`$ be
$`J_2({\displaystyle \frac{\varphi }{v}})=\left(\begin{array}{cccc}\varphi _3^1& \varphi _4^1& \mathrm{}& \varphi _n^1\\ \varphi _3^2& \varphi _4^2& \mathrm{}& \varphi _n^2\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \varphi _3^{n2}& \varphi _4^{n2}& \mathrm{}& \varphi _n^{n2}\end{array}\right)`$
and suppose that $`detJ_2(\frac{\varphi }{v})|_{(t^{},p_i)}0.`$ With the same reasons of obtaining (27), we can have the function relations
$$v^A=f^A(v^1,v^2,t,\stackrel{}{u}),A=3,4,\mathrm{},n.$$
(40)
For the partial derivatives $`f_1^A`$, $`f_2^A`$ and $`f_t^A`$, we can easily derive the system similar to the equations (28) and (29), in which the three terms at the right hand of can be figured out at the same time. In order to determine the 2–order partial derivatives $`f_{11}^A`$, $`f_{12}^A`$, $`f_{1t}^A`$, $`f_{22}^A`$, $`f_{2t}^A`$ and $`f_{tt}^A`$, we can use the method similar to the above mentioned. Substituting the relations (40) into the last two equations of the system (19), we have the following two equations with respect to the arguments $`v^1,v^2,t,\stackrel{}{u}`$
$$\{\begin{array}{c}F_1(v^1,v^2,t,\stackrel{}{u})=\varphi ^{n1}(v^1,v^2,f^3(v^1,v^2,t,\stackrel{}{u}),\mathrm{},f^n(v^1,v^2,t,\stackrel{}{u}),t,\stackrel{}{u})=0\hfill \\ F_2(v^1,v^2,t,\stackrel{}{u})=\varphi ^n(v^1,v^2,f^3(v^1,v^2,t,\stackrel{}{u}),\mathrm{},f^n(v^1,v^2,t,\stackrel{}{u}),t,\stackrel{}{u})=0.\hfill \end{array}$$
(41)
Calculating the partial derivatives of the function $`F_1`$ and $`F_2`$ with respect to $`v^1`$, $`v^2`$ and $`t`$, taking notice of (40) and using six similar expressions to (33) and (34), i.e.
$$\frac{F_j}{v^1}_{(t^{},p_i)}=0,\frac{F_j}{v^2}_{(t^{},p_i)}=0,\frac{F_j}{t}_{(t^{},p_i)}=0,j=1,2,$$
(42)
we have the following forms of Taylor expressions of $`F_1`$ and $`F_2`$ in the neighborhood of $`(t^{},p_i)`$
$`F_j(v^1,v^2,t,\stackrel{}{u})A_{j1}(v^1p_i^1)^2+A_{j2}(v^1p_i^1)(v^2p_i^2)+A_{j3}(v^1p_i^1)`$
$`(tt^{})+A_{j4}(v^2p_i^2)^2+A_{j5}(v^2p_i^2)(tt^{})+A_{j6}(tt^{})^2=0`$
$$j=1,2.$$
(43)
In the case of $`A_{j1}0,A_{j4}0`$, by dividing (43) by $`(tt^{})^2`$ and taking the limit $`tt^{}`$, we obtain two quadratic equations of $`\frac{dv^1}{dt}`$ and $`\frac{dv^2}{dt}`$
$$A_{j1}(\frac{dv^1}{dt})^2+A_{j2}\frac{dv^1}{dt}\frac{dv^2}{dt}+A_{j3}\frac{dv^1}{dt}+A_{j4}(\frac{dv^2}{dt})^2+A_{j5}\frac{dv^2}{dt}+A_{j6}=0$$
(44)
$`j=1,2.`$
Eliminating the variable $`dv^1/dt`$, we obtain a equation of $`dv^2/dt`$ in the form of a determinant
$$\left|\begin{array}{cccc}A_{11}& A_{12}Q+A_{23}& A_{14}Q^2+A_{15}Q+A_{16}& 0\\ 0& A_{11}& A_{12}Q+A_{13}& A_{14}Q^2+A_{15}Q+A_{16}\\ A_{21}& A_{22}Q+A_{23}& A_{24}Q^2+A_{25}Q+A_{26}& 0\\ 0& A_{21}& A_{22}Q+A_{23}& A_{24}Q^2+A_{25}Q+A_{26}\end{array}\right|=0$$
(45)
where $`Q=dv^2/dt`$, which is a $`4th`$ order equation of $`dv^2/dt`$
$$a_0(\frac{dv^2}{dt})^4+a_1(\frac{dv^2}{dt})^3+a_2(\frac{dv^2}{dt})^2+a_3(\frac{dv^2}{dt})+a_4=0.$$
(46)
Therefore we get different directions at the higher degenerated point corresponding to different branch curves. The number of different branch curves is four at most. If the degree of degeneracy of the matrix $`[\frac{\varphi }{v}]`$ is more higher, i.e. the rank of the matrix $`[\frac{\varphi }{v}]`$ is more lower than the present $`(n2)`$ case, the procedure of deduction will be more complicate. In general supposing the rank of the matrix $`[\frac{\varphi }{x}]`$ be $`(ns)`$, the number of the possible different directions of the branch curves is $`2^s`$ at most.
At the end of this section, we conclude that there exist crucial cases of branch processes in our theory of topological defect system in TDGL model. This means that a topological defect, at the bifurcation point, may split into several (for instance $`m`$) topological defects along different branch curves with different charges. Since the topological current is a conserved current, the total quantum number of the splitting topological defects must precisely equal to the topological charge of the original defect i.e.
$`{\displaystyle \underset{j=1}{\overset{m}{}}}\beta _{i_j}\eta _{i_j}=\beta _i\eta _i`$
for fixed $`i`$. This can be looked upon as the topological reason of the defect splitting. Here we should point out that such splitting is a stochastic process, the sole restriction of this process is just the conservation of the topological charge of the topological defects during this splitting process. Of course, the topological charge of each splitting defects is an integer.
## Acknowledgment
The author gratefully acknowledges the support of K. C. Wong Education Foundation, Hong Kong. The author is also gratefully indebted to Prof. Y. S. Duan for his warm-hearted helps and useful discussion.
## Figures’ Captions
Fig. 1. (a) The creation of two topological defects. (b) Two topological defects annihilate in collision at the limit point.
Fig. 2. Two topological defects collide with different directions of motion at the bifurcation point.
Fig. 3. Topological defects have the same direction of motion. (a) Two topological defects tangentially collide at the bifurcation point. (b) Two topological defects merge into one topological defect at the bifurcation point. (c) One topological defect splits into two topological defects at the bifurcation point.
Fig. 4. (a) One topological defect splits into three topological defects at the bifurcation point. (b) Three topological defects merge into one topological defect at the bifurcation point.
Fig. 5. This case is similar to Fig. 4. (a) Three topological defects merge into one topological defect at the bifurcation point. (b) One topological defect splits into three topological defects at the bifurcation point.
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# Transport with three-particle interaction
## I Introduction
The equation of state of nuclear matter is known to saturate at the density of $`n_0=0.16`$fm<sup>-3</sup> with a binding energy of $`16`$MeV. Refined two-particle calculations using Bruckner theory and beyond are not crossing the “Coester line” in that they lead to binding energies and/or densities above the saturation ones. Only density dependent Skyrme parameterizations originally introduced by Skyrme , three-particle interactions or relativistic treatments can reproduce the correct binding energy and saturation density of nuclear matter. The Skyrme parameterizations are derived from an effective three-particle interaction . This leads to a nonlinear density dependence in the effective three - particle part which is responsible for saturation. The importance of three-particle collisions in nuclear matter transport has been demonstrated, e.g. in . The density dependence deviating from that arising by three-body contact interaction has been introduced and compared with experiments in .
The relativistic approach on the other hand yields immediately the correct saturation with two - particle exchange interactions. This has been traced down to the nonlinear density dependence of scalar density and consequently nonlinear mean fields which leads to a density contribution to the binding energy of $`(n/n_0)^{8/3}`$. Higher order effects have been shown to lead to a $`(n/n_0)^{3.4}`$ dependence, see and citations therein. Therefore the physics of the saturation mechanism is certainly a nontrivial density dependence of the mean field. Though the Skyrme interaction has been overwhelmingly successful and the evidence of saturation by three-particle interaction , it is puzzling that a transport theory with three-particle contact interaction has not been formulated. In this paper we will derive the corresponding kinetic equation using nonequilibrium Green’s function and we will show that the three-body interaction term leads to a $`(n/n_0)^{10/3}`$ contribution to the binding energy.
The transport theory including three - particle interactions seems to be of wider interest e.g. in it has been found that the three-body interactions have a measurable influence on thermodynamic properties of fluids in equilibrium. The description of nonequilibrium in nuclear matter is mostly based on the Boltzmann (BUU) equation which uses the Skyrme parameterization to determine the mean field and drift side of the kinetic equation. The collisional integral is then calculated with a cross section which arises from different theoretical impact. In addition one usually neglects the three-particle collision integral. These hybrid models have worked quite successfully despite their weak rigour in microscopic foundation.
In this paper the transport equation will be derived for two- and three - particle contact interactions. We will see that a natural density dependent mean field appears, in agreement with variational methods. Further we will obtain, on the same ground, density dependent cross sections. This has the advantage that we can derive a BUU equation which has the same microscopic impact on the drift and collisional side. Finally the three -particle collision integral appears naturally from this treatment. From balance equations we will derive the density dependent energy which gives the equation of state.
## II Three - particle Kadanoff and Baym equation
We consider a Hamilton system with the Hamiltonian
$`H`$ $`=`$ $`{\displaystyle \underset{i}{}}a_i^+{\displaystyle \frac{\mathrm{}^2^2}{2m}}a_i`$ (3)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}a_i^+a_j^+a_ia_jV_2(i,j)`$
$`+{\displaystyle \frac{1}{6}}{\displaystyle \underset{ijk}{}}a_i^+a_j^+a_k^+a_ia_ja_kV_3(i,j,k)`$
where $`a_i^+,a_i`$ are creation and annihilation operators with cumulative variables $`i=(x_i,t_i,\mathrm{})`$. We assume the potential is contact - like as
$`V_2(i,j)`$ $`=`$ $`t_0\delta (ij)`$ (4)
$`V_3(i,j,k)`$ $`=`$ $`t_3\delta (ij)\delta (ik).`$ (5)
The Heisenberg equations of motion for the annihilation operators read
$`i\mathrm{}_{t_1}a_1`$ $`=`$ $`[a_1,H]`$ (6)
$`=`$ $`{\displaystyle \frac{\mathrm{}^2_1^2}{2m}}a_1+{\displaystyle \underset{2}{}}V(1,2)a_2^+a_1a_2`$ (8)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{23}{}}V(1,2,3)a_2^+a_3^+a_1a_2a_3.`$
From this we get the equation of motion for the causal Green function $`G(1,2)=\frac{1}{i}<Ta_1^+a_2>`$ with the time ordering $`T`$ and the average taken about the nonequilibrium density operator
$`(i\mathrm{}_{t_1}{\displaystyle \frac{\mathrm{}^2_1^2}{2m}})G(1,1^{})=\delta (11^{})`$ (9)
$`+{\displaystyle d2V(1,2)G_2(1,2,1^{},2^+)}`$ (10)
$`+{\displaystyle \frac{1}{2}}{\displaystyle d2d3V(1,2,3)G_3(1,2,3,1^{},2^+,3^+)}`$ (11)
$`\delta (11^{}){\displaystyle 𝑑x_2\mathrm{\Sigma }_{HF}(x_1,x_2,t)G(x_2,t,x_1,t)}`$ (12)
$`+{\displaystyle d2(\mathrm{\Sigma }(1,2)G(2,1^{})\mathrm{\Sigma }^>(1,2)G^<(2,1^{}))}`$ (13)
where the $`1^+`$ denotes the space - time point $`1=(x_1,t_1+ϵ)`$ with a time infinitesimally larger than $`1`$. This equation is enclosed in the standard manner introducing the self energy $`\mathrm{\Sigma }`$. We have split the mean field parts $`\mathrm{\Sigma }_{HF}`$ including exchange and have introduced for the rest the self energy on the Keldysh contour. This is equivalent to the weakening of initial correlations. The correlation functions are $`G^<(1,2)=<a_2^+a_1>`$ and $`G^>(1,2)=<a_1a_2^+>`$ respectively and their equation of motion follows from (13) in the form of the known Kadanoff and Baym equation
$`i(G_0^1G^<G^<G_0^1)=\mathrm{\Sigma }^RG^<G^<\mathrm{\Sigma }^A`$ (14)
where one takes advantage of the retarded and advanced functions $`G^{R/A}=i\theta [\pm (t_1t_2)](G^>G^<)`$ and understands products as integration over inner variables.
### A Mean field parts
We first calculate the mean field parts. These parts come from the first order interaction diagrams of figure 1. To this end we use the two- and three- particle Green function in lowest order approximation seen in figure 2.
Introducing these diagrams into figure 1, one obtains for contact interaction (5)
$`\mathrm{\Sigma }^{\mathrm{HF}}(1,1^{})=\left[{\displaystyle \frac{g\pm 1}{g}}t_0n(1)+{\displaystyle \frac{(g\pm 1)(g\pm 2)}{2g^2}}t_3n(1)^2\right]\delta (11^{})`$ (15)
(16)
where we have used the fact that the density is $`n(1)=G^<(1,1)/g`$. Here and in the following we write the upper sign for Bosons and the lower sign for Fermions. The result (16) is the known Skyrme mean field expression in nuclear matter for $`g=4`$. It resembles an effective density dependent two-particle interaction arising from three-particle contact interaction. As a consistency check we see that for $`g=2`$ degenerated Fermionic system, like spin -1/2 Fermions, no three-body term appears. Pauli-blocking forbids two particles to meet at the same point with equal quantum numbers which one would have for three - particle contact interaction and degeneracy of $`g=2`$.
### B Kinetic equation
For the kinetic equation we need as the lowest order approximation the next diagrams of figure 2 including one interaction line. While there are different expansion techniques at hand we will perform here a scheme as close as possible to the causal Green’s function including all exchanges which are presented in figure 2. This will give us the advantage that we can consider all diagrams which differ by exchange of outgoing lines as equal. To this aim we introduce the abbreviated symmetrized Green’s function with respect to incoming lines in figure 3.
The expansion of the three -body Green’s function is given in figure 5. Besides the three -particle interaction potential we have to consider all possible single line interactions between two incoming lines. This is shortened by the introduction of the symmetrized Green function in figure 3 and the vertices of figure 4.
Since we can also have nontrivial combinations of two incoming lines and one free line for the three - body Green function we have introduced an auxiliary two - body Green function $`G_2^\mathrm{s}`$ which is defined in figure 5. Please note that the exchange of outgoing lines does not lead to distinguished diagrams. The two - particle Green function is then given by figure 5. Introducing now this expansion into the definition of self energy in figure 1 we obtain the 4 diagrams of figure 6. All other combinations lead either to disconnected diagrams or to equivalent ones by interchanging outgoing lines. Here the seemingly disconnected diagrams for $`G_3^1`$ in figure 5 lead together with the corresponding diagrams from $`G_2^1`$ to the vertex in front of $`G_2^s`$ of figure 6.
The enclosing diagram about the auxiliary Green function $`G_2^s`$ is given in figure 7. Please denote that here the seemingly disconnected diagrams of figure 5 vanish. We have used again the fact that the interchange of outgoing lines does not lead to distinguished diagrams.
Inserting figure 7 in figure 6 we obtain the final result for the self energy in the Born approximation
$`\mathrm{\Sigma }^<(1,1^{})=(g\pm 1)[t_0^2`$ (17)
$`+(g\pm 2)(4t_0^2+6t_3^2{\displaystyle \frac{n(1)n(1^{})}{g^2}}+5t_0t_3({\displaystyle \frac{n(1)}{g}}+{\displaystyle \frac{n(1^{})}{g}}))]`$ (18)
$`\times G^<(1,1^{})^2G^>(1^{}1)`$ (19)
$`+{\displaystyle \frac{1}{2}}(g\pm 1)(g\pm 2)t_3^2G^<(1,1^{})^3G^>(1^{}1)^2.`$ (20)
Note that all interactions vanish for the fermionic case $`g=1`$ since the Pauli-principle does not allow contact interaction in this case. For $`g=2`$ only the two - particle interactions survive for fermions as one expects from the Pauli principle. That the result is symmetric in the densities at the two time-space points needs no further comment.
From the self energy we can write down the Kadanoff and Baym equation (14) in closed form. If we employ standard gradient expansion and convert the equation (14) into an equation for the pole part of the Green’s function we obtain the final kinetic equation
$`{\displaystyle \frac{}{t}}f_1(t)+{\displaystyle \frac{p_1}{m}}{\displaystyle \frac{}{r}}f_1(t){\displaystyle \frac{}{r}}\sigma ^{\mathrm{HF}}(t){\displaystyle \frac{}{p_1}}f_1(t)`$ (21)
$`=`$ $`{\displaystyle \frac{2}{\mathrm{}^2}}{\displaystyle \frac{dp_2dp_1^{}dp_2^{}}{(2\pi \mathrm{})^6}T_2^2\delta (p_1+p_2p_1^{}p_2^{})}`$ (22)
$`\times `$ $`{\displaystyle _{t_0}^t}𝑑\tau \mathrm{cos}\left[{\displaystyle \frac{1}{\mathrm{}}}(E_1+E_2E_1^{}E_2^{})(t\tau )\right]`$ (23)
$`\times `$ $`\left\{f_1^{}(\tau )f_2^{}(\tau )\overline{f}_1(\tau )\overline{f}_2(\tau )f_1(\tau )f_2(\tau )\overline{f}_1^{}(\tau )\overline{f}_2^{}(\tau )\right\}`$ (24)
$`+`$ $`{\displaystyle \frac{2}{\mathrm{}^2}}{\displaystyle \frac{dp_2dp_3dp_1^{}dp_2^{}dp_3^{}}{(2\pi \mathrm{})^{12}}T_3^2\delta (p_1+p_2+p_3p_1^{}p_2^{}p_3^{})}`$ (25)
$`\times `$ $`{\displaystyle _{t_0}^t}𝑑\tau \mathrm{cos}\left[(E_1+E_2+E_3E_1^{}E_2^{}E_3^{}){\displaystyle \frac{t\tau }{\mathrm{}}}\right]`$ (27)
$`\times \{f_1^{}(\tau )f_2^{}(\tau )f_3^{}(\tau )\overline{f}_1(\tau )\overline{f}_2(\tau )\overline{f}_3(\tau )`$
$``$ $`f_1(\tau )f_2(\tau )f_3(\tau )\overline{f}_1^{}(\tau )\overline{f}_2^{}(\tau )\overline{f}_3^{}(\tau )\}`$ (28)
with $`\overline{f}=1f`$, and the particle dispersion $`E=p^2/2m+\sigma ^{\mathrm{HF}}`$. The quasiparticle distribution functions are normalized to the density as $`g\frac{dp}{(2\pi \mathrm{})^3}f(p)=n`$. For the sake of simplicity we have suppressed the notation of centre of mass space dependence. Neglecting the retardation in the distribution functions and taking $`t_0\mathrm{}`$ gives the standard Boltzmann two- and three- particle collision integrals.
The introduced two - and three- particle T-matrices are read off from (20) as first Born approximation
$`T_2^2`$ $`=`$ $`{\displaystyle \frac{g\pm 1}{g^2}}[t_0^2`$ (30)
$`+`$ $`(g\pm 2)(4t_0^2+6t_3^2{\displaystyle \frac{n(1)n(1^{})}{g^2}}+5t_0t_3({\displaystyle \frac{n(1)}{g}}+{\displaystyle \frac{n(1^{})}{g}}))]`$ (31)
$`T_3^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{g\pm 1}{g^3}}(g\pm 2)t_3^2.`$ (32)
The reader is remind that we have generally to sum over collision partners and to average over the outgoing products. While the summation is already taken into account explicitly in the diagrammatic approach above we had to divide by $`g^2`$ and $`g^3`$ for the two- and three - particle outgoing collision partners respectively. An effective cross section can be defined in the usual way and read up to constants
$`{\displaystyle \frac{d\sigma _2}{d\mathrm{\Omega }}}`$ $``$ $`T_2^2={\displaystyle \frac{3}{16}}\left[9t_0^2+{\displaystyle \frac{3}{4}}t_3^2n^2+5t_0t_3n\right]`$ (35)
$`={\displaystyle \frac{27}{16}}\left[(t_0+{\displaystyle \frac{1}{6}}t_3n)^2+{\displaystyle \frac{2}{9}}t_3n(t_0+{\displaystyle \frac{1}{4}}t_3n)\right]`$
$`{\displaystyle \frac{d\sigma _3}{d\mathrm{\Omega }}}`$ $``$ $`T_3^2={\displaystyle \frac{3}{64}}t_3^2.`$ (36)
If one compares this with the effective two- particle cross section derived in one sees that it differs by the last two terms of (36) proportional to $`t_3`$. This is due to the fact that the underlying three - particle process is taken into account explicitly here while in an effective two - particle kinetic theory has been developed. Obviously we have to face a partial cancellation between three - and two - particle collisions. This will become explicit in the discussion of correlation energies. We will see that the two - and three - particle correlation energies indeed cancel partially concerning the $`t_3`$ terms.
The ratio of the effective cross sections are for the case of nuclear matter ($`g=4`$)
$`{\displaystyle \frac{d\sigma _2}{d\sigma _3}}={\displaystyle \frac{T_2^2}{T_3^2}}`$ $`=`$ $`4\left[9{\displaystyle \frac{t_0^2}{t_3^2}}+{\displaystyle \frac{3}{4}}n^2+5{\displaystyle \frac{t_0}{t_3}}n\right].`$ (37)
This could serve as the measure of relative importance of the corresponding collision processes, but we prefer to discuss this later in terms of dispersed energy by the two- and three- particle collision integrals since this includes also the Pauli blocking.
We would like to point out that we have fulfilled our task and have derived a kinetic equation where the drift as well as the collision integral follows from the same microscopic impact. The drift represents a Skyrme mean field and the collision side shows a two - and three - particle collision integral where the cross sections are calculated from the same parameters as the mean field.
## III Balance equation
By multiplying the kinetic equation with $`1,p`$ and $`E`$, one obtains the balance for density $`n`$, momentum $`u`$ and energy $``$. Since the collision integrals vanish for density and momentum balance we get the usual balance equations
$`{\displaystyle \frac{n}{t}}+{\displaystyle \frac{}{R}}{\displaystyle \frac{dp}{(2\pi \mathrm{})^3}\frac{E}{p}f}=0`$ (38)
$`{\displaystyle \frac{u_i}{t}}+{\displaystyle \frac{}{R_j}}{\displaystyle \frac{dp}{(2\pi \mathrm{})^3}(p_i\frac{E}{p_j}f+\delta _{ij})}=0`$ (39)
where the mean field energy of the system varies as
$`\delta `$ $`=`$ $`{\displaystyle \frac{dp}{(2\pi \mathrm{})^3}\frac{\delta }{\delta f(p)}\delta f(p)}`$ (40)
$`=`$ $`{\displaystyle \frac{dp}{(2\pi \mathrm{})^3}E\delta f(p)}`$ (41)
such that from (16) follows
$`=<{\displaystyle \frac{k^2}{2m}}>+{\displaystyle \frac{g\pm 1}{2g}}t_0n^2+{\displaystyle \frac{(g\pm 1)(g\pm 2)}{6g^2}}t_3n^3.`$ (42)
With the help of this quantity the balance of energy density reads
$`{\displaystyle \frac{}{t}}+{\displaystyle \frac{}{R}}{\displaystyle \frac{dp}{(2\pi \mathrm{})^3}E\frac{E}{p}f}={\displaystyle \frac{}{t}}E_{\mathrm{corr}_2}(t){\displaystyle \frac{}{t}}E_{\mathrm{corr}_3}(t)`$ (43)
with the two-particle correlation energy
$`E_{\mathrm{corr}_2}(t)`$ (45)
$`={\displaystyle \frac{g}{\mathrm{}}}{\displaystyle \underset{t_0}{\overset{t}{}}}𝑑\tau {\displaystyle \frac{dp_1dp_2dp_1^{}dp_2^{}}{(2\pi \mathrm{})^9}T_2^2\delta (p_1+p_2p_1^{}p_2^{})}`$ (46)
$`\times \mathrm{sin}\left[(E_1+E_2E_1^{}E_2^{}){\displaystyle \frac{t\tau }{\mathrm{}}}\right]f_1^{}(\tau )f_2^{}(\tau )\overline{f}_1(\tau )\overline{f}_2(\tau )`$ (47)
(48)
and the complete analogous expression for the three particle energy.
Please note that in the balance equations for the density and momentum no correlated density or correlated flux appears. This is due to the Born approximation. For more nontrivial approximations as e.g. the ladder summation these correlated observables appear .
### A Fit to nuclear matter ground state
As one can see, it is not enough to use the mean field parameterization to fit the equation of state as done in most approaches so far. Since the collision integral induces a two - particle and three - particle correlation energy we have to take this into account and have to refit the parameter $`t_0,t_3`$ to the binding properties of nuclear matter. To illustrate this fact let us evaluate the fit without two - and three - particle correlation energy and therefore without collision integrals.
#### 1 Hartree-Fock parameterization
Taking into account only Hartree-Fock mean field correlations (42) the ground state energy for nuclear matter reads
$`{\displaystyle \frac{E}{A}}={\displaystyle \frac{}{n}}={\displaystyle \frac{3}{5}}ϵ_f+{\displaystyle \frac{3}{8}}t_0n+{\displaystyle \frac{1}{16}}t_3n^2.`$ (49)
Then the nuclear binding at $`n_0=0.16`$ fm<sup>-3</sup> with an energy of $`E_B=15.68`$MeV is reproduced by the set
$`t_0`$ $`=`$ $`{\displaystyle \frac{16}{15n_0}}(2ϵ_f5E_B)=1026.67\mathrm{MeV}\mathrm{fm}^3`$ (50)
$`t_3`$ $`=`$ $`{\displaystyle \frac{16}{5n_0^2}}(ϵ_f5E_B)=14625\mathrm{M}\mathrm{e}\mathrm{v}\mathrm{fm}^6`$ (51)
which leads to a compressibility of
$`K=9n^2{\displaystyle \frac{^2E/A}{n^2}}={\displaystyle \frac{9}{8}}t_3n_0^2{\displaystyle \frac{6}{5}}ϵ_f=377\mathrm{M}\mathrm{e}\mathrm{V}.`$ (52)
#### 2 Parameterization with two - and three - particle correlation energy
Now we consider the two - particle correlation energy. From (21) we obtain the total energy
$`+E_{\mathrm{corr}_2}+E_{\mathrm{corr}_3}.`$ (53)
We want to calculate the long time limit which represents the equilibrium value. From the identity $`\underset{0}{\overset{\mathrm{}}{}}\mathrm{sin}xtdt=\frac{𝒫}{x}`$ a principle value integration replaces the $`sin`$ term in (48)
$`E_{\mathrm{corr}_2}(\mathrm{})=g{\displaystyle \frac{dp_1dp_2dp_1^{}dp_2^{}}{(2\pi \mathrm{})^9}T_2^2\frac{𝒫}{E_1+E_2E_1^{}E_2^{}}}`$ (54)
$`\times \delta (p_1+p_2p_1^{}p_2^{})f_1^{}f_2^{}\overline{f}_1\overline{f}_2.`$ (55)
For ground state Fermi distributions this expression can be integrated analytically and we find from (55) the known Galitskii result for the two - particle ground state correlation energy
$`{\displaystyle \frac{E_{\mathrm{corr}_2}}{n}}`$ $`=`$ $`4ϵ_f{\displaystyle \frac{2\mathrm{log}211}{35}}({\displaystyle \frac{p_fmT_2}{4\pi ^2\mathrm{}^3}})^2`$ (56)
$`=`$ $`{\displaystyle \frac{5.7910^5}{\mathrm{MeV}\mathrm{fm}^2}}n^{4/3}\left[9t_0^2+{\displaystyle \frac{3}{4}}t_3^3n^2+5t_3t_0n\right].`$ (57)
As pointed out in we had to subtract here an infinite value, i.e. the term proportional to $`f_1f_2`$ in (55). This can be understood as renormalization of the contact potential and is formally hidden in $`E_{\mathrm{corr}}(0)`$ when time integrating equation (43). For finite range potentials we have an intrinsic cut-off due to range of interaction and such problems do not occur.
The correlational two - particle energy (57) is always negative for fermionic degeneracies $`2<g<8`$ and for bosonic degeneracy $`g<4`$. Since the leading density goes with $`n^{10/3}`$ it will dominate over the positive kinetic part which goes like $`n^{2/3}`$. We find that for densities around $`3n_0`$ the total energy has a maximum and starts to decrease continuously for higher densities. The maximal energy is changed towards higher values if we now include three \- particle correlations since the latter remain positive and have the same leading density term as the two - particle correlational energy.
The three - particle part reads
$`E_{\mathrm{corr}_3}(\mathrm{})=g{\displaystyle \frac{dp_1dp_2dp_3dp_1^{}dp_2^{}dp_3^{}}{(2\pi \mathrm{})^{15}}T_3^2}`$ (58)
$`\times `$ $`\delta (p_1+p_2+p_3p_1^{}p_2^{}p_3^{})f_1^{}f_2^{}f_3^{}\overline{f}_1\overline{f}_2\overline{f}_3`$ (59)
$`\times `$ $`{\displaystyle \frac{𝒫}{E_1+E_2+E_3E_1^{}E_2^{}E_3^{}}}.`$ (60)
The analytic result for the contact potential has not been given in the literature to our knowledge and reads
$`{\displaystyle \frac{E_{\mathrm{corr}_3}}{n}}`$ $`=`$ $`4ϵ_f{\displaystyle \frac{9013}{29257713}}({\displaystyle \frac{p_f^4mT_3}{4\pi ^4\mathrm{}^6}})^2`$ (61)
$`=`$ $`{\displaystyle \frac{2.3710^6}{\mathrm{MeV}\mathrm{fm}^2}}n^{10/3}t_3^2.`$ (62)
This three - particle correlation energy remains positive and has the same leading density behaviour of $`n^{10/3}`$ as the two - particle correlational energy. Since the pre-factor is smaller than the one for the two - particle correlational energy we obtain a maximum at 3 times nuclear density beyond which the two - particle part dominates and the total energy diverges negatively which would mean a collapse of the system. This clearly marks the limit of the Born approximation.
Taking now the correlational energy into account via (53) instead of only the Hartree Fock energy (49) we obtain a fit to the nuclear binding of
$`\stackrel{~}{t}_0`$ $`=`$ $`745.71\mathrm{MeV}\mathrm{fm}^3`$ (63)
$`\stackrel{~}{t}_3`$ $`=`$ $`8272.8.\mathrm{MeV}\mathrm{fm}^6`$ (64)
which leads to a compressibility of
$`\stackrel{~}{K}=351\mathrm{M}\mathrm{e}\mathrm{V}`$ (65)
which is somewhat lower than the mean field compressibility of (52). The comparison of the two equation of states with and without two - and three-particle correlations can be seen in figure 8. The inclusion of two - particle correlation energy leads to a maximum at 3 times nuclear density above which the system collapses. The complete result including three particle correlational energy leads to a higher reachable maximum.
### B Importance of three particle collisions
We have seen that the three - particle contact interaction induces in a natural way a density dependence of the two - particle cross section. Additionally we have obtained an explicit three - particle collision integral on the same microscopic footing. We want now to answer the question how important the three - particle collisions are. Therefore we use as a measure the ratio of the two- and three - particle ground state correlation energies, (57) and (62), since this gives the measure of how much energy is maximally dispersed by the corresponding integrals. We obtain
$`\left|{\displaystyle \frac{E_2}{E_3}}\right|`$ $`=`$ $`{\displaystyle \frac{2^551113}{9013}}(112\mathrm{log}2)`$ (66)
$`\times `$ $`\left[9{\displaystyle \frac{t_0^2}{n^2t_3^2}}+{\displaystyle \frac{3}{4}}+5{\displaystyle \frac{t_0}{t_3n}}\right].`$ (67)
This ratio is decreasing until it reaches the density
$`n_{\mathrm{min}}={\displaystyle \frac{18t_0}{5t_3}}=0.32\mathrm{fm}^3`$ (68)
where the ratio has the minimum
$`\left|{\displaystyle \frac{E_2}{E_3}}\right|_{\mathrm{min}}=1.4.`$ (69)
For higher densities the constant value
$`\left|{\displaystyle \frac{E_2}{E_3}}\right|_{\mathrm{}}=18.3`$ (70)
is approached.
I other words this means that the importance of three particle collisions increase with increasing density up to twice the nuclear density where the correlational energies are nearly equal. For higher densities we have 18 times larger two -particle correlational energies than three -particle ones. At nuclear saturation density this ratio is
$`\left|{\displaystyle \frac{E_2}{E_3}}\right|_{n_0}=19.2`$ (71)
indicating that the three -particle collisions become important between nuclear density and twice the nuclear density while it can be neglected in the other cases.
## IV Conclusions
For a microscopic two - and three - particle contact interaction consisting of two parameters the Kadanoff and Baym equation of motion is derived. From this a Boltzmann kinetic equation is obtained with drift which turns out to be the known Skyrme Hartree Fock expression while the collision side consists of two - and three - particle collision integrals. The two - particle collision integral contains an explicit density - dependent cross section arising from the three - particle contact interaction. By this way both the drift side as well as collision side is derived from the same microscopic footing and no hybrid assumptions about separate density dependent mean field and cross section is needed.
The correlational energy for the three- and the two - particle part are calculated analytically. Due to the density dependent two - particle collision integral both correlational energies have the same leading power of $`(n/n_0)^{10/3}4`$ in density which shows that both contributions are of the same importance if three - particle interaction has to be considered. This is clearly motivated by the saturation point where non - relativistic two - particle approaches fail to overcome the “Coester line”.
We find that there is a maximum in the energy at 3 times nuclear density if two - and three particle correlational energies are included. Beyond this density the Born approximation fails at least in that the system collapses towards diverging negative energy.
While the two - particle collision cross section has a natural density dependence due to three - particle contact interaction it turns out that the explicit three particle collision integral can be neglected as long as one is below nuclear density. Around twice nuclear density the three particle collision integral has the same importance as the two - particle one since it disperses the same amount of energy. For higher densities the three - particle collision integral is again negligible.
I would like to thank H.S. Köhler and P. Lipavsky for numerous discussions and the LPC for a friendly and hospitable atmosphere. P. Chocian is thanked for reading the manuscript.
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# Topological Symmetries
## 1 Introduction
Since Witten’s pioneering work on supersymmetric quantum mechanics , there has been a growing interest in supersymmetry and its applications . The interest in supersymmetry has also motivated the development of various geralizations of supersymmetry. Most notable of these are parasupersymmetries , the $`q`$-deformed supersymmetries and the fractional supersymmetries .
One of the most intriguing aspects of supersymmetric quantum mechanics is its relationship with the Atiyah-Singer index theorem . This has already been noticed by Witten in early 1980’s and subsequently led to new proofs of this theorem . Supersymmetric proofs of the index theorem together with Witten’s supersymmetric derivation of the Morse inequalities and its impact on Floer’s theory are among the greatest mathematical achievments of supersymmetric quantum mechanics.
The effectiveness of supersymmetry in providing insight into some of the most fundamental results of differential geometry and topology suggests the study of the topological content of various generalizations of supersymmetry. The only known generalization of supersymmetry which exhibits similar topological properties is a certain type of $`p=2`$ parasupersymmetries. The topological properties of $`p=2`$ parasupersymmetry (PSUSY) has been extensively studied in Refs. .
The purpose of the present article is to introduce a generalization of supersymmetry (SUSY) which shares its topological properties. We shall term such a symmetry a topological quantum mechanical symmetry or simply topological symmetry (TS).
The first step towards such a generalization is to recall the basic properties of the $`N=1`$ supersymmetric quantum mechanics.<sup>1</sup><sup>1</sup>1We use the terminology in which $`N`$ is half of the number of Hermitian generators. What we call $`N=1`$ SUSY sometimes is referred to as the $`𝒩=2`$ SUSY where $`𝒩`$ is the number of Hermitian SUSY generators.
$`N=1`$ supersymmetric quantum mechanics is specified by a $`ZZ_2`$-graded Hilbert space $`=_+_{}`$ and the superalgebra
$`[H,𝒬]`$ $`=`$ $`0,`$ (1)
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{𝒬,𝒬^{}\},`$ (2)
$`𝒬^2`$ $`=`$ $`0,`$ (3)
where $`H`$ and $`𝒬`$ are the Hamiltonian and the generator of the supersymmetry, respectively. The $`ZZ_2`$-grading of the Hilbert space is implemented using a chirality or parity operator $`\tau :`$ satisfying
$`\tau ^2`$ $`=`$ $`1,\tau ^{}=\tau ,`$ (4)
$`[\tau ,H]`$ $`=`$ $`0,`$ (5)
$`\{\tau ,𝒬\}`$ $`=`$ $`0.`$ (6)
The subspaces $`_+`$ and $`_{}`$ are identified with the eigenspaces of $`\tau `$,
$$_\pm :=\{|\psi |\tau |\psi =\pm |\psi \}.$$
(7)
The elements of $`_\pm `$ are said to have definite parity $`\pm `$. An operator $`O`$ acting on $``$ is said to have definite parity $`+`$ or $``$, if it commutes or anticommutes with $`\tau `$, respectively.
It is well-known that using Eqs. (1) – (3) one can obtain the degeneracy structure of a general $`N=1`$ supersymmetric system . These systems have a nonnegative spectrum and the eigenspaces corresponding to positive eigenvalues are spanned by pairs of eigenvectors of opposite parity. This particular degeneracy structure of supersymmetric systems is sufficient to show the topological invariance of the Witten index:
$`\mathrm{index}_\mathrm{W}:=n_0^+n_0^{},`$ (8)
$`n_0^\pm :=\mathrm{number}\mathrm{of}\mathrm{zero}\mathrm{energy}\mathrm{states}\mathrm{with}\mathrm{parity}\pm .`$ (9)
The degeneracy structure is obtained from the algebraic structure. Therefore, the situation may be described by the following diagram.
$$\mathrm{algebraic}\mathrm{structure}\mathrm{degeneracy}\mathrm{structure}\mathrm{topological}\mathrm{invariants}$$
(10)
The same analysis is valid for the case of the $`p=2`$ PSUSY studied in Refs. .
The idea pursued in this article is to reverse the first arrow in (10). More specifically, we wish to
* find and postulate the type of degeneracy structures which lead to topological invariants such as the Witten index, and
* obtain the algebraic structure of symmetries which support this type of degeneracy structures.
## 2 Topological Symmetries and Their Invariants
* Definition 1: Let $`m_+`$ and $`m_{}`$ be two positive integers. Then a quantum mechanical symmetry is called a $`ZZ_2`$-graded topological symmetry (TS) of type $`(m_+,m_{})`$, if the following conditions are fulfilled.
+ The Hilbert space is $`ZZ_2`$-graded. The $`ZZ_2`$-grading is achieved via a parity operator $`\tau `$ satisfying Eqs. (4), i.e., $`=_+_{}`$ where $`_\pm `$ are given by Eq. (7).
+ The energy spectrum of the system is nonnegative.
+ There is an energy eigenbasis consisting of definite parity state vectors, i.e., Eq. (5) holds.
+ For every positive energy eigenvalue E, there exists a positive integer $`\lambda _E`$ such that $`E`$ is $`m_E:=\lambda _E(m_{}+m_+)`$ fold degenerate. Furthermore, the corresponding eigenspace is spanned by $`\lambda _Em_{}`$ negative parity eigenvectors and $`\lambda _Em_+`$ positive parity eigenvectors.
Clearly, SUSY is an example of TS of type $`(1,1)`$. Therefore, TSs are generalizations of the SUSY.
* Definition 2: A $`ZZ_2`$-graded topological symmetry is said to be uniform, if for all $`E>0`$, $`\lambda _E=1`$.
Given the above definition of a $`ZZ_2`$-graded uniform topological symmetry (UTS), we can easily prove the following theorem.
* Theorem: Consider a $`ZZ_2`$-graded UTS of type $`(m_{},m_+)`$ and let $`n_0^\pm `$ denote the number of zero energy eigenstates of the system with parity $`\pm `$. Then the quantity
$$\mathrm{\Delta }_{(m_+,m_{})}:=m_{}n_0^+m_+n_0^{}$$
(11)
is a topological invariant, i.e., it is invariant under the continuous changes of the quantum system<sup>2</sup><sup>2</sup>2These include continuous changes of the Hamiltonian and the boundary conditions. that do not destroy the UTS.
* Proof: Suppose that under a continuous change of the system a zero energy state vector with positive parity is elevated to a positive energy level. This positive energy level must have $`m_+`$ positive parity states and $`m_{}`$ negative parity states. Hence, the initial zero energy state must be accompanied by $`(m_+1)`$ positive parity zero energy eigenstates and $`m_{}`$ negative parity zero energy eigenstates. This implies that after the change $`\mathrm{\Delta }_{(m_+,m_{})}`$ is given by
$$\mathrm{\Delta }_{(m_+,m_{})}^{\mathrm{after}}=m_{}(n_0^+m_+)m_+(n_0^{}m_{})=m_{}n_0^+m_+n_0^{}=\mathrm{\Delta }_{(m_+,m_{})}^{\mathrm{before}}.$$
Every possible change of the zero energy states is a combination of this particular change and its converse. Therefore, in general $`\mathrm{\Delta }_{(m_+,m_{})}`$ remains invariant. $`\mathrm{}`$
The same proof is valid for the case of nonuniform $`ZZ_2`$-graded TSs. Therefore, $`\mathrm{\Delta }_{(m_+,m_{})}`$ is a topological invariant for any $`ZZ_2`$-graded TS.
We shall next provide the basic framework for addressing the characterization problem for TSs. In this paper we shall only consider the uniform TSs. But our method applies to nonuniform TSs as well.
We shall demand TSs to have symmetry generators $`𝒬_a`$ (with $`a=1,2,\mathrm{},N`$) which have negative parity. In particular, we shall only consider the $`N=1`$ UTSs where the label $`a=1`$ can be dropped. In this case, Eq. (6) is valid.
Next we introduce the Hermitian symmetry generators,
$$Q_1:=\frac{1}{\sqrt{2}}(𝒬+𝒬^{}),\mathrm{and}Q_2=\frac{i}{\sqrt{2}}(𝒬𝒬^{}).$$
(12)
In view of Eqs. (1), (4), (6) and (12), we have
$`[H,Q_i]`$ $`=`$ $`0,`$ (13)
$`\{\tau ,Q_i\}`$ $`=`$ $`0,`$ (14)
where $`i\{1,2\}`$.
We can use Eq. (5) to construct an orthonormal basis of the Hilbert space in which $`H`$ and $`\tau `$ are diagonal. Our strategy will be to use the information on the degeneracy structure of the energy eigenspaces and Eqs. (4), (13), and (14) to obtain matrix representations of $`Q_1`$ and $`Q_2`$ in the energy eigenspaces $`_E`$ with eigenvalue $`E>0`$. We shall denote the representation of an operator $`O`$ in the eigenspace $`_E`$ by $`O^E`$. Clearly for the Hamiltonian $`H`$, we have $`H^E=EI_m`$, where $`I_m`$ is the $`m\times m`$ identity matrix and $`m:=m_++m_{}`$. Finally, we shall try to use these representations to obtain the most general algebraic relations satisfied by $`𝒬`$ and $`H`$.
## 3 Uniform Topological Symmetries of Type $`(1,1)`$
For the $`ZZ_2`$-graded TS of type $`(1,1)`$, the positive energy levels are doubly degenerate ($`m=2`$). In a basis that diagonalizes $`H`$ and $`\tau `$, we have (up to a permutation of the basis vectors)
$$\tau ^E=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),$$
(15)
where $`E>0`$ and we have used Eqs. (4). Next let us note that $`Q_i^E`$ are $`2\times 2`$ Hermitian matrices satisfying (14). These conditions are sufficient to conclude that
$$Q_1^E=\left(\begin{array}{cc}0& \mu ^{}\\ \mu & 0\end{array}\right)\mathrm{and}Q_2^E=\left(\begin{array}{cc}0& \nu ^{}\\ \nu & 0\end{array}\right),$$
(16)
where $`\mu `$ and $`\nu `$ are complex numbers.
Using the matrix representations (15) and (16), we can easily compute
$`𝒬^E`$ $`:=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(Q_1^E+iQ_2^E)={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}0& \mu ^{}+i\nu ^{}\\ \mu +i\nu & 0\end{array}\right),`$ (19)
$`(Q_1^E)^2`$ $`=`$ $`|\mu |^2I_2,`$ (20)
$`(Q_2^E)^2`$ $`=`$ $`|\nu |^2I_2,`$ (21)
$`(𝒬^E)^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[|\mu |^2|\nu |^2+i(\mu \nu ^{}+\mu ^{}\nu )\right]I_2=det(𝒬^E)I_2,`$ (22)
where ‘$`det`$’ stands for ‘determinant’. Next we introduce the Hermitian operators $`M`$, $`K_1`$ and $`K_2`$ which commute with $`H`$ and have the following representations in $`_E`$ with $`E>0`$.
$$M^E=|\mu |^2I_2,K_1^E=(|\mu |^2|\nu |^2)I_2,K_2^E=(\mu \nu ^{}+\mu ^{}\nu )I_2.$$
(23)
In view of Eqs. (20) – (23), $`M^E`$, $`K^E`$, and $`K_i^E`$ commute with $`Q_i^E,𝒬^E`$ and $`\tau ^E`$. Generalizaing these equations to operator identities, we find
$`[M,Q_i]=[K_j,Q_i]=[M,K_j]=0,`$ (24)
$`Q_1^2=M,`$ (25)
$`Q_2^2=MK_1,`$ (26)
$`\{Q_1,Q_2\}=K_2.`$ (27)
We can also express Eqs. (24) – (27) in terms of $`𝒬`$. This yields
$`[M,𝒬]=[K,𝒬]=[M,K]=0,`$ (28)
$`{\displaystyle \frac{1}{2}}\{𝒬,𝒬^{}\}=M{\displaystyle \frac{1}{2}}(K+K^{}),`$ (29)
$`𝒬^2=K,`$ (30)
where $`K:=(K_1+iK_2)/2`$.
Next let us note that under a linear transformations of the form:
$$\begin{array}{ccc}Q_1& & \stackrel{~}{Q}_1:=aQ_1+bQ_2,\\ Q_2& & \stackrel{~}{Q}_2:=cQ_1+dQ_2,\end{array}$$
(31)
the algebra (24) – (27) is left form-invariant.<sup>3</sup><sup>3</sup>3Here the coefficients $`a,b,c`$ and $`d`$ are assumed to be Hermitian operators commuting with $`H`$, $`Q_i`$, $`K_i`$ and $`M`$. More specifically, the transformed generators $`\stackrel{~}{Q}_i`$ satisfy the same algebra provided that the operators $`M`$ and $`K_i`$ are transformed according to
$`M`$ $``$ $`\stackrel{~}{M}:=(a^2+b^2)Mb^2K_1+abK_2,`$ (32)
$`K_1`$ $``$ $`\stackrel{~}{K}_1:=(a^2+b^2c^2d^2)M+(d^2b^2)K_1+(abcd)K_2,`$ (33)
$`K_2`$ $``$ $`\stackrel{~}{K}_2:=2(ac+bd)M2bdK_1+(ad+bc)K_2.`$ (34)
This observation may be used to find a new set of negative parity symmetry generators $`\stackrel{~}{Q}_i`$ which would satisfy the algebra (24) – (27) with $`K_i`$ set to zero. The most general linear transformations (31) for which $`\stackrel{~}{K}_1=\stackrel{~}{K}_2=0`$ are the ones satisfying (either of)
$$\frac{a+ic}{b+id}=\frac{K_2}{2M}\pm i\sqrt{1\frac{K_1}{M}\frac{K_2^2}{4M^2}}.$$
(35)
Using the energy eigenbasis in which Eqs. (23) hold, we can show that indeed the terms inside the square root in (35) add up to a positive operator.<sup>4</sup><sup>4</sup>4Strictly speaking most of the conclusions drawn in this article are based on the information obtained from the restriction of the relevant operators to $`_0`$. We will however suppose that the same results are generally valid in $``$. This is consistent with the fact that we have not given the representations of $`M`$ and $`K_i`$ in $`_0`$. This is a necessary condition for the corresponding linear transformation to be invertible.
The above analysis shows that without loss of generality we can set $`K=0`$ in Eqs. (28) – (30). This yields
$`[M,𝒬]=0,`$ (36)
$`{\displaystyle \frac{1}{2}}\{𝒬,𝒬^{}\}=M,`$ (37)
$`𝒬^2=0.`$ (38)
Since $`M`$ has the same degeneracy structure as the Hamiltonian (at least in $`_0`$), we can write $`H`$ as a function of $`M`$. In particular, we can identify $`H`$ with $`M`$, in which case the algebra (36) – (38) reduces to the superalgebra (1) – (3). Therefore, the algebra of the $`ZZ_2`$-graded UTS of type $`(1,1)`$ is the same as the algebra of SUSY. The above analysis may be viewed as a derivation of the superalgebra (1) – (3) from a set of basic principles, i.e., the definition of the $`ZZ_2`$-graded UTS of type $`(1,1)`$.
## 4 Uniform Topological Symmetry of Type $`(2,1)`$
For the $`ZZ_2`$-graded UTS of type $`(2,1)`$, the positive energy levels are triply degenerate ($`m=3`$). We will again work in a basis in which $`H`$ and $`\tau `$ are diagonal. Restricting to an eigenspace $`_E`$ of $`H`$ with $`E>0`$ and enforcing Eqs. (4) and (14), we can easily show that (up to permutations of the basis vectors of $`_E`$)
$`\tau ^E=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right),Q_1^E=\left(\begin{array}{ccc}0& 0& \mu _1^{}\\ 0& 0& \mu _2^{}\\ \mu _1& \mu _2& 0\end{array}\right),`$ (45)
$`Q_2^E=\left(\begin{array}{ccc}0& 0& \nu _1^{}\\ 0& 0& \nu _2^{}\\ \nu _1& \nu _2& 0\end{array}\right),𝒬^E={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{ccc}0& 0& \mu _1^{}+i\nu _1^{}\\ 0& 0& \mu _2^{}+i\nu _2^{}\\ \mu _1+i\nu _1& \mu _2+i\nu _2& 0\end{array}\right).`$ (52)
Now in order to obtain the most general algebraic identities satisfied by these matrices we appeal to the Cayley-Hamilton theorem . This theorem states that any $`m\times m`$ matrix $`A`$ satisfies its characteristic equation, $`P_A(x)=0`$, where $`P_A(x)`$ is the characteristic polynomial for $`A`$. It is not difficult to show that the characteristic polynomial for a $`3\times 3`$ matrix $`A`$ of the form
$$A=\left(\begin{array}{ccc}0& 0& \alpha \\ 0& 0& \beta \\ \gamma & \delta & 0\end{array}\right)$$
is given by $`P_A(x)=x^3(\alpha \gamma +\beta \delta )x`$. Using this equation and the identity $`P_A(A)=0`$ for $`Q_1^E`$, $`Q_2^E`$ and $`𝒬^E`$, we find
$`(Q_1^E)^3`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}|\mu _j|^2Q_1^E,`$ (53)
$`(Q_2^E)^3`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}|\nu _j|^2Q_2^E,`$ (54)
$`(𝒬^E)^3`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{2}{}}}[|\mu _j|^2|\nu _j|^2+i(\mu _j\nu _j^{}+\mu _j^{}\nu _j)]𝒬^E.`$ (55)
Next we introduce the Hermitian operators $`M`$, $`K_1`$, and $`K_2`$ which commute with $`H`$ and have the following matrix representations in $`_E`$ with $`E>0`$.
$`M^E`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}|\mu _j|^2I_3,`$ (56)
$`K_1^E`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}(|\mu _j|^2|\nu _j|^2)I_3,`$ (57)
$`K_2^E`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}(\mu _j\nu _j^{}+\mu _j^{}\nu _j)I_3.`$ (58)
Now setting $`K:=(K_1+iK_2)/2`$ and making use of Eqs. (56) – (58), we can generalize Eqs. (53) – (55) to the operator relations
$`Q_1^3`$ $`=`$ $`MQ_1,`$ (59)
$`Q_2^3`$ $`=`$ $`(MK_1)Q_2,`$ (60)
$`𝒬^3`$ $`=`$ $`K𝒬.`$ (61)
Eqs. (56) – (58) also suggest that the operators $`M`$ and $`K_i`$ commute among themselves and with $`Q_i`$. Furthermore, we can rewrite Eq. (61) in terms of $`Q_1`$ and $`Q_2`$. Simplifying the resulting equation using Eqs. (59) and (60) we obtain the following general algebra
$`[M,Q_i]=[M,K_i]=[K_i,Q_j]=0,`$ (62)
$`Q_1^3=MQ_1,`$ (63)
$`Q_2^3=(MK_1)Q_2,`$ (64)
$`Q_2Q_1Q_2+\{Q_1,Q_2^2\}=(MK_1)Q_1+K_2Q_2,`$ (65)
$`Q_1Q_2Q_1+\{Q_2,Q_1^2\}=MQ_2+K_2Q_1.`$ (66)
We can show by direct computation that this algebra remains form-invariant under linear transformations (31) of $`Q_i`$. More remarkable is the fact that under such a transformation the operators $`M`$ and $`K_i`$ transform according to the same relations as in the case of UTS of type $`(1,1)`$, namely Eqs. (32) – (34). Therefore, we can always transform to a new set of symmetry generators for which $`K_i=0`$. Setting $`K_i=0`$ in Eqs. (62) – (66) and rewriting these equations in terms of $`𝒬`$, we find
$`[M,𝒬]=0,`$ (67)
$`\{𝒬^2,𝒬^{}\}+𝒬𝒬^{}𝒬=2M𝒬,`$ (68)
$`𝒬^3=0.`$ (69)
The operator $`M`$ has the same degeneracy structure as the Hamiltonian. Therefore, we can identify $`H`$ with a function of $`M`$. In particular, we can set $`H=M/2`$. In this case the algebra (67) – (69) becomes identical to the algebra of $`p=2`$ parasupersymmetry . In other words, the algebra of the $`ZZ_2`$-graded UTS of type $`(2,1)`$ is precisely the algebra of the $`p=2`$ parasupersymmetry (of Rubakov and Spiridonov ). Again, the above analysis may be viewed as a derivation of the algebra of $`p=2`$ parasupersymmetry from a set of basic principles, i.e., the definition of the $`ZZ_2`$-graded UTS of type $`(2,1)`$.
As shown in Ref. , the algebra (67) – (69) does not imply the degeneracy structure of type $`(2,1)`$ UTS. Therefore, the type $`(2,1)`$ $`ZZ_2`$-graded UTSs belong to a special class of symmetries whose generator $`𝒬`$ satisfies Eqs. (67) – (69). A method for constructing the corresponding moduli spaces is given in Ref. .
## 5 Conclusion
In this article we have introduced a general notion of a topological symmetry. We have provided a simple framework for the study of the $`ZZ_2`$-graded topological symmetries. We showed that the algebras of the $`ZZ_2`$-graded UTS of order $`(1,1)`$ and $`(2,1)`$ are essentially the algebras of supersymmetric quantum mechanics and $`p=2`$ parasupersymmetric quantum mechanics, respectively.
By construction, topological symmetries involve a class of integer-valued topological invariants $`\mathrm{\Delta }_{(m_+,m_{})}`$. These are the analogues of the Witten index of supersymmetry and the parasupersymmetric topological invariant . The physical interpretation of these invariants is that they are a measure of the existence of zero-energy ground states. For the known cases, the latter is an indication of the exactness of symmetry. The mathematical interpretation of $`\mathrm{\Delta }_{(m_+,m_{})}`$ is not quite clear. For the known cases they are related to the analytic indices of Fredholm operators . The general case requires a detailed study of the algebraic structure of general topological symmetries. The algebras of $`ZZ_2`$-graded TS of arbitrary type $`(m_+,m_{})`$ are currently under investigation.
Finally, one can easily generalize the definition of the $`ZZ_2`$-graded topological symmetry of type $`(m_+,m_{})`$ to a $`ZZ_n`$-graded topological symmetry of type $`(m_1,m_2,\mathrm{},m_n)`$. Such a system will have states with definite ‘color’ taking values in $`\{1,2,\mathrm{},n\}`$. The spectrum will be nonnegative. The positive energy eigenvalues $`E`$ will be $`m_E=\lambda _E_{\mathrm{}=1}^nm_{\mathrm{}}`$ fold degenerate. The energy eigenspaces with positive eigenvalue will all have $`\lambda _Em_1`$ states of color $`1`$, $`\lambda _Em_2`$ states of color $`2`$, $`\mathrm{}`$, and $`\lambda _Em_n`$ states of color $`n`$. One may try to define topological invariants for these more general topological symmetries. For example for a $`ZZ_3`$-graded TS of type $`(1,1,1)`$, one can introduce the invariant $`\mathrm{\Delta }_{(1,1,1)}=(n_0^{(1)}n_0^{(2)})^2+(n_0^{(2)}n_0^{(3)})^2+(n_0^{(3)}n_0^{(1)})^2`$, where $`n_0^{(\mathrm{})}`$ is the number of zero energy states of color $`\mathrm{}`$. The basic properties of $`ZZ_n`$-graded topological symmetries will be explored in .
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# 1 Introduction
## 1 Introduction
Since the discovery of asymptotic freedom in non-Abelian gauge field theories, like Quantum Chromo Dynamics (QCD), many perturbative calculations have been performed to hadron-hadron and semi-leptonic processes. The most reactions are computed up to next-to-leading order (NLO) only, except for some semi-leptonic processes. The latter is due to the simplicity of the Born approximation where the basic process is given by the interaction of the intermediate vector boson with a quark. Higher order corrections to the other reactions are still missing because of the very complicated Feynman and phase space integrals which arise in the calculations. Here we will discuss some new methods which maybe will enable us to compute the QCD and also QED corrections beyond the NLO level.
## 2 Asymptotic Expansions
In this section we concentrate on the computation of the coefficient functions and anomalous dimensions which show up in deep inelastic lepton-hadron scattering. The relevant quantity in this process is the structure function defined by
$`F(x,Q^2)={\displaystyle \frac{1}{4\pi }}{\displaystyle d^4ze^{iq.z}p[J(z),J(0)]p}q^2=Q^2<0x={\displaystyle \frac{Q^2}{2\nu }}.`$ (1)
In the commutator above $`J(z)`$ represents the electro-weak current where we have suppressed Lorentz indices for convenience. In the Bjorken limit i.e. $`Q^2\mathrm{}`$ and $`x=fixed`$ the light cone region dominates the integrand so that one can make an Operator Product Expansion (OPE)
$`[J(z),J(0)]\stackrel{=}{z^20}{\displaystyle \underset{N}{}}C^N(z^2\mu ^2)O^N(\mu ^2,0).`$ (2)
This expansion also holds for the time ordered product corresponding to forward Compton scattering
$`T(x,Q^2)`$ $`=`$ $`{\displaystyle \frac{i}{4\pi ^2}}{\displaystyle d^4ze^{iq.z}p[J(z),J(0)]p},`$
$`\mathrm{Im}T(x,Q^2)`$ $`=`$ $`\pi F(x,Q^2).`$ (3)
Let us assume that we can write an unsubtracted dispersion relation in the variable $`\nu =pq`$
$`T(\nu ,Q^2)={\displaystyle _{Q^2/2}^{\mathrm{}}}𝑑\nu ^{}{\displaystyle \frac{F(\nu ^{},Q^2)}{\nu ^{}\nu }},`$ (4)
so that after substitution of the variables $`\nu =\frac{Q^2}{2x}`$, $`\nu ^{}=\frac{Q^2}{2x^{}}`$ we can derive the following relation
$`T(\nu ,Q^2)`$ $`=`$ $`x{\displaystyle _0^1}{\displaystyle \frac{dx^{}}{x^{}}}{\displaystyle \frac{F(x^{},Q^2)}{xx^{}}}`$ (5)
$`=`$ $`{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{2pq}{Q^2}}\right)^N{\displaystyle _0^1}𝑑x^{}x^{N1}F(x^{},Q^2)`$
$`=`$ $`{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{2pq}{Q^2}}\right)^NA^{(N)}(\mu ^2)𝒞^{(N)}\left({\displaystyle \frac{Q^2}{\mu ^2}}\right),`$
where the operator matrix element and the coefficient function are given by
$`A^{(N)}(\mu ^2)=pO^N(\mu ^2,0)p,`$ (6)
and
$`𝒞^{(N)}\left({\displaystyle \frac{Q^2}{\mu ^2}}\right)={\displaystyle d^4ze^{iq.z}C^N(z^2\mu ^2)},`$ (7)
respectively. Both quantities satisfy a renormalization group equation
$`\left[\mu {\displaystyle \frac{}{\mu }}+\beta (g){\displaystyle \frac{}{g}}+\gamma ^{(N)}(g)\right]A^{(N)}(\mu ^2)=0,`$
$`\left[\mu {\displaystyle \frac{}{\mu }}+\beta (g){\displaystyle \frac{}{g}}\gamma ^{(N)}(g)\right]𝒞^{(N)}\left({\displaystyle \frac{Q^2}{\mu ^2}}\right)=0.`$ (8)
In the equations above the quantities $`\beta (g)`$ and $`\gamma ^{(N)}(g)`$ represent the beta-function and the anomalous dimension respectively. The latter determine the $`Q^2`$-evolution of the structure function $`F(x,Q^2)`$ which can be measured in experiment and provides us with one of the tests of perturbative QCD. If one computes the coefficient functions in the conventional way one encounters phase space integrals and loop integrals (see e.g. ). However if one wants to compute the coefficient function beyond order $`g^4`$ it is more convenient to try another method which is explained in . Taking $`\varphi _6^3`$-theory as an example one can expand the propagator in $`T(x,Q^2)`$ (2) given by $`1/(kp)^2`$ as follows
$`T(x,Q^2)T`$ $`=`$ $`(ig)^2{\displaystyle \frac{d^nk}{(2\pi )^n}\frac{i^3}{(k^2)^2(kp)^2(k+q)^2}}`$ (9)
$`=`$ $`ig^2{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d^nk}{(2\pi )^n}\frac{(2kp)^N}{(k^2)^{3+N}(k+q)^2}},`$
where we have used n-dimensional regularization to regularize the collinear (C) and ultraviolet (UV) singularities. The asymptotic expansion of the propagator converts the Compton amplitude into a self energy type of integral. The latter is much easier to compute than expressions for box- or triangle graphs. Another feature is that this asymptotic expansion transforms the C-divergence in $`T`$ at $`n=6`$ into an UV singularity appearing in the expression for the self energy integral. Therefore the operator matrix element in Eq. (5), containing the collinear divergence denoted by $`1/\epsilon _C`$, will be replaced by the corresponding operator renormalization constant $`Z_{O^N}`$
$`Q^2T`$ $`=`$ $`{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{2pq}{Q^2}}\right)^NA^{(N)}({\displaystyle \frac{1}{\epsilon _C}},\mu ^2)𝒞^{(N)}\left({\displaystyle \frac{Q^2}{\mu ^2}}\right)`$ (10)
$`{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{2pq}{Q^2}}\right)^NZ_{O^N}({\displaystyle \frac{1}{\epsilon _{UV}}},\mu ^2)𝒞^{(N)}\left({\displaystyle \frac{Q^2}{\mu ^2}}\right),`$
where $`Z_{O^N}`$ depends on the UV singularity denoted by $`1/\epsilon _{UV}`$. This transformation leaves the coefficient function unaltered. A straightforward calculation yields
$`Q^2T`$ $`=`$ $`g^2{\displaystyle \frac{\pi ^{n/2}}{(2\pi )^n}}{\displaystyle \frac{\mathrm{\Gamma }(n/21)\mathrm{\Gamma }(n/23)}{\mathrm{\Gamma }(n4)}}{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{2pq}{Q^2}}\right)^N{\displaystyle \frac{\mathrm{\Gamma }(4+Nn/2)}{\mathrm{\Gamma }(3+N)}}\left({\displaystyle \frac{Q^2}{\mu ^2}}\right)^{n/23}.`$ (11)
One can expand the gamma-functions above around $`\epsilon =n6`$ and obtain
$`Q^2T`$ $`=`$ $`{\displaystyle \frac{g^2}{64\pi ^3}}{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{2pq}{Q^2}}\right)^N[{\displaystyle \frac{1}{N+1}}{\displaystyle \frac{1}{N+2}}][{\displaystyle \frac{2}{\epsilon }}+\gamma _E\mathrm{ln}4\pi `$ (12)
$`+\mathrm{ln}{\displaystyle \frac{Q^2}{\mu ^2}}1{\displaystyle \underset{k=1}{\overset{N}{}}}{\displaystyle \frac{1}{k}}].`$
In order $`g^2`$ Eq. (10) can be expressed as
$`Q^2T={\displaystyle \frac{g^2}{64\pi ^3}}{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}\left[{\displaystyle \frac{\gamma _O^{(N)}}{2}}\left\{{\displaystyle \frac{2}{\epsilon }}+\gamma _E\mathrm{ln}4\pi +\mathrm{ln}{\displaystyle \frac{Q^2}{\mu ^2}}\right\}+c_q^{(N)}\right],`$ (13)
so that one can read $`\gamma _O^{(N)}`$ and $`c_q^{(N)}`$. If the method is extended up to order $`g^4`$ one encounters graphs where at least three propagators, carrying the momentum $`p`$ (see the solid line in Fig. (1)), have to be expanded
This procedure leads to triple sums which are very difficult to unravel
$`{\displaystyle \underset{i}{}}{\displaystyle \underset{j}{}}{\displaystyle \underset{l}{}}\left({\displaystyle \frac{2pk_1}{k_1^2}}\right)^i\left({\displaystyle \frac{2pk_2}{k_2^2}}\right)^j\left({\displaystyle \frac{2p(k_1k_2q)}{(k_1k_2q)^2}}\right)^j.`$ (14)
To avoid this one has to remove one or more propagators, carrying the momentum $`p`$, before doing the expansion. First one can adopt the method
1. Integration by Parts
This trick has been successfully applied to compute self energy graphs . An example is
$`I={\displaystyle \frac{d^nk_1}{(2\pi )^n}\frac{d^nk_2}{(2\pi )^n}\frac{1}{k_1^2k_2^2(k_1+q)^2(k_2+q)^2(k_1k_2)^2}}.`$ (15)
The integral above appears in the partial integration w.r.t. the momentum $`k_1`$ of the expression below (see Fig. (2))
$`0={\displaystyle \frac{d^nk_1}{(2\pi )^n}\frac{d^nk_2}{(2\pi )^n}\frac{}{k_{1,\mu }}\frac{(k_1k_2)_\mu }{k_1^2k_2^2(k_1+q)^2(k_2+q)^2(k_1k_2)^2}}.`$ (16)
Hence $`I`$ in Eq. (15) can be reduced to a set of integrals each containing one propagator less than in the original expression.
2. Recursion Relations (Difference Equations)
Another method is developed in . Using the method of ”Difference Equations” one is able to knock out one propagator (see the example in Fig. (3)). One can show that all graphs can be reduced to two building blocks only (see Fig. (4)). The latter contain one propagator carrying the momentum $`p`$ so that one only obtains a single sum which is easy to perform. At this moment it is not clear whether the method can be extended beyond order $`g^4`$ (two-loop).
In higher order it is only possible to obtain the integral over the structure function in Eq. (5) for finite moments $`N`$ (see ). This is achieved by
$`{\displaystyle \frac{^NT}{p_{\mu _1}\mathrm{}p_{\mu _N}}}_{p_{\mu _i}=0}=2^NN!{\displaystyle \frac{q_{\mu _1}\mathrm{}q_{\mu _N}}{(Q^2)^N}}Z_{O^N}({\displaystyle \frac{1}{\epsilon _{UV}}},\mu ^2)𝒞^{(N)}\left({\displaystyle \frac{Q^2}{\mu ^2}}\right)`$ (17)
In our example Eq. (9) this procedure leads to the result
$`{\displaystyle \frac{^NT}{p_{\mu _1}\mathrm{}p_{\mu _N}}}_{p_{\mu _i}=0}=2^NN!(ig)^2{\displaystyle \frac{d^nk}{(2\pi )^n}\frac{k_{\mu _1}\mathrm{}k_{\mu _N}}{(k^2)^{N+2}(k+q)^2}}`$ (18)
which is again of the self energy type.
## 3 Mellin Barnes Techniques
A method which turns out to be very useful to compute three-point and four-point functions even up to two-loop level is given by the Mellin-Barnes technique . A simple Mellin-Barnes transformation takes the form
$`(A_1+A_2)^\nu ={\displaystyle \frac{1}{2\pi i}}{\displaystyle _i\mathrm{}^i\mathrm{}}A_1^\sigma A_2^{\nu \sigma }{\displaystyle \frac{\mathrm{\Gamma }(\sigma )\mathrm{\Gamma }(\nu +\sigma )}{\mathrm{\Gamma }(\nu )}},`$ (19)
with $`|arg(A_1)arg(A_2)|<\pi `$ and the poles in $`\sigma `$ are located on the real axis. To illustrate this technique we take as an example the three-point function where the propagators are raised to an arbitrary power
$`J(n;\nu _1,\nu _2,\nu _3)`$ $`=`$ $`{\displaystyle d^nk\frac{1}{[(q_1+k)^2m_1^2]^{\nu _1}[(q_2+k)^2m_2^2]^{\nu _2}[(q_3+k)^2m_3^2]^{\nu _3}}}`$ (20)
$`=`$ $`i^{1n}\pi ^{n/2}{\displaystyle \frac{\mathrm{\Gamma }(_i\nu _in/2)}{_i\mathrm{\Gamma }(\nu _i)}}{\displaystyle \underset{i}{}d\alpha _i\alpha _i^{\nu _i1}\delta (1\underset{j}{}\alpha _j)}`$
$`\times {\displaystyle \frac{1}{[_j\alpha _jm_j^2+\alpha _2\alpha _3p_1^2+\alpha _1\alpha _3p_2^2+\alpha _1\alpha _2p_3^2]^{_k\nu _kn/2}}},`$
where $`p_1=q_3q_2`$, $`p_2=q_1q_3`$ and $`p_3=q_2q_1`$. Next we apply the double Mellin Barnes integral transformation given by
$`{\displaystyle \frac{1}{(X+Y+Z)^a}}`$ $`=`$ $`{\displaystyle \frac{1}{Z^a}}{\displaystyle _i\mathrm{}^i\mathrm{}}{\displaystyle \frac{du}{2\pi i}}{\displaystyle _i\mathrm{}^i\mathrm{}}{\displaystyle \frac{dv}{2\pi i}}{\displaystyle \frac{\mathrm{\Gamma }(a+u+v)\mathrm{\Gamma }(u)\mathrm{\Gamma }(v)}{\mathrm{\Gamma }(a)}}`$ (21)
$`\left({\displaystyle \frac{X}{Z}}\right)^u\left({\displaystyle \frac{Y}{Z}}\right)^v,`$
so that the integral in Eq. (20) can be written as
$`J(n;\nu _1,\nu _2,\nu _3)={\displaystyle \frac{\pi ^{n/2}i^{1n}}{_j\mathrm{\Gamma }(\nu _j)}}{\displaystyle \underset{i}{}d\alpha _i\alpha _i^{\nu _i1}\delta (1\underset{j}{}\alpha _j)\left(\alpha _1\alpha _2p_3^2\right)^{n/2_k\nu _k}}`$
$`{\displaystyle _i\mathrm{}^i\mathrm{}}{\displaystyle \frac{du}{2\pi i}}{\displaystyle _i\mathrm{}^i\mathrm{}}{\displaystyle \frac{dv}{2\pi i}}\mathrm{\Gamma }({\displaystyle \underset{i}{}}\nu _in/2+u+v)\mathrm{\Gamma }(u)\mathrm{\Gamma }(v)\left({\displaystyle \frac{\alpha _3p_1^2}{\alpha _1p_3^2}}\right)^u\left({\displaystyle \frac{\alpha _3p_2^2}{\alpha _2p_3^2}}\right)^v,`$ (22)
where we have put for simplicity $`m_i^2=0`$. Using
$`{\displaystyle \underset{i}{}}{\displaystyle _0^1}𝑑\alpha _i\delta (1{\displaystyle \underset{j}{}}\alpha _j)\alpha _1^a\alpha _2^b\alpha _3^c={\displaystyle \frac{\mathrm{\Gamma }(1+a)\mathrm{\Gamma }(1+b)\mathrm{\Gamma }(1+c)}{\mathrm{\Gamma }(1+a+b+c)}},`$ (23)
the three-point function reads
$`J(n;\nu _1,\nu _2,\nu _3)=`$
$`\pi ^{n/2}i^{1n}{\displaystyle _i\mathrm{}^i\mathrm{}}{\displaystyle \frac{du}{2\pi i}}{\displaystyle _i\mathrm{}^i\mathrm{}}{\displaystyle \frac{dv}{2\pi i}}{\displaystyle \frac{\mathrm{\Gamma }(n/2\nu _1\nu _3)\mathrm{\Gamma }(n/2\nu _2\nu _3)\mathrm{\Gamma }(\nu _3+u+v)}{_i\mathrm{\Gamma }(\nu _i)\mathrm{\Gamma }(n_j\nu _j)}}`$
$`\mathrm{\Gamma }({\displaystyle \underset{k}{}}\nu _k+u+vn/2)\mathrm{\Gamma }(u)\mathrm{\Gamma }(v)\left({\displaystyle \frac{p_1^2}{p_3}}\right)^u\left({\displaystyle \frac{p_2^2}{p_3}}\right)^u\left(p_3\right)^{n/2_j\nu _j}.`$ (24)
A special case is $`\nu _1+\nu _2+\nu _3=n`$ (Uniqueness Condition). After application of the residue theorem to expression (3) one obtains
$`J(n;\nu _1,\nu _2,\nu _3)=\pi ^{n/2}i^{1n}{\displaystyle \underset{j=1}{\overset{3}{}}}{\displaystyle \frac{\mathrm{\Gamma }(n/2\nu _j)}{\mathrm{\Gamma }(\nu _j)}}\left(p_j^2\right)^{\nu _jn/2}.`$ (25)
The Mellin Barnes techniques has been successfully applied to compute two-loop graphs. Examples are the three-point functions in Fig. 5 and the four point functions in Fig. 6 which are computed in for $`p_i^20`$. Recent progress has been made for the two-loop box graphs in Fig. 6 which has been computed for all external momenta on-shell ($`p_i^2=0`$) , . This is very useful for applications to the second order corrections in Bhabha scattering and in di-jet production.
## 4 Negative Dimension Approach
This method has been recently developed in to express one-loop integrals containing different masses into special functions. To illustrate the technique we take the two-point function as an example. The latter is given by the integral
$`I_2^n(\nu _1,\nu _2,q^2,M_1^2,M_2^2)={\displaystyle \frac{d^nk}{i\pi ^{n/2}}\frac{1}{A_1^{\nu _1}A_2^{\nu _2}}},`$ (26)
Using the Schwinger representation for the propagator
$`{\displaystyle \frac{1}{A_i^{\nu _i}}}={\displaystyle \frac{(1)^{\nu _i}}{\mathrm{\Gamma }(\nu _i)}}{\displaystyle _0^{\mathrm{}}}𝑑x_ix_i^{\nu _i1}exp\left(x_iA_i\right),`$ (27)
expression (26) becomes equal to
$`I_2^n={\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{(1)^{\nu _i}}{\mathrm{\Gamma }(\nu _i)}}{\displaystyle _0^{\mathrm{}}}𝑑x_ix_i^{\nu _i1}{\displaystyle \frac{d^nk}{i\pi ^{n/2}}exp\left(\underset{j=1}{\overset{2}{}}x_jA_j\right)}.`$ (28)
The expression above can be evaluated in two different ways. The first way proceeds by expanding the exponents in a power series. This yields
$`I_2^n`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{(1)^{\nu _i}}{\mathrm{\Gamma }(\nu _i)}}{\displaystyle _0^{\mathrm{}}}𝑑x_ix_i^{\nu _i1}{\displaystyle \frac{d^nk}{i\pi ^{n/2}}\underset{n_1=0}{\overset{\mathrm{}}{}}\frac{(x_1A_1)^{n_1}}{n_1!}\underset{n_2=0}{\overset{\mathrm{}}{}}\frac{(x_2A_2)^{n_2}}{n_2!}}=`$ (29)
$`{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{(1)^{\nu _i}}{\mathrm{\Gamma }(\nu _i)}}{\displaystyle _0^{\mathrm{}}}𝑑x_ix_i^{\nu _i1}{\displaystyle \underset{n_1=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n_2=0}{\overset{\mathrm{}}{}}}I_2^n(n_1,n_2,q^2,M_1^2,M_2^2){\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle \frac{x_j^{n_j}}{n_j!}}.`$
Using the identity
$`{\displaystyle \frac{(1)^{\nu _i}}{\mathrm{\Gamma }(\nu _i)}}{\displaystyle _0^{\mathrm{}}}𝑑x_ix_i^{\nu _i+n_j1}=\delta _{\nu _i+n_j,0},`$ (30)
we obtain
$`I_2^n=I_2^n(\nu _1,\nu _2,q^2,M_1^2,M_2^2){\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{1}{\mathrm{\Gamma }(1\nu _i)}}.`$ (31)
There is a second way to compute Eq. (28). First one shifts the momentum $`k`$ which becomes $`k=k^{}qx_2/(x_1+x_2)`$. Subsequently one uses the following identities
$`{\displaystyle \frac{d^nk^{}}{i\pi ^{n/2}}\left(k^2\right)^l}=l!\delta _{l+n/2,0}{\displaystyle \frac{d^nk^{}}{i\pi ^{n/2}}exp\left(\alpha k^2\right)}={\displaystyle \frac{1}{\alpha ^{n/2}}}.`$ (32)
Notice that since $`l>0`$ one obtains $`n<0`$. This the reason for the name ”Negative Dimension Approach”. The expression in Eq. (28) can now be written as
$`I_2^n={\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{(1)^{\nu _i}}{\mathrm{\Gamma }(\nu _i)}}{\displaystyle _0^{\mathrm{}}}𝑑x_ix_i^{\nu _i1}{\displaystyle \frac{1}{(x_1+x_2)^{n/2}}}exp\left({\displaystyle \frac{x_1x_2}{x_1+x_2}}q^2\right)exp\left({\displaystyle \underset{j=1}{\overset{2}{}}}x_jM_j^2\right).`$ (33)
The integrand can be written as
$`{\displaystyle \frac{1}{(x_1+x_2)^{n/2}}}exp\left({\displaystyle \frac{x_1x_2}{x_1+x_2}}q^2\right)exp\left({\displaystyle \underset{j=1}{\overset{2}{}}}x_jM_j^2\right)`$
$`={\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{l!}}(x_1+x_2)^{ln/2}(x_1x_2q^2)^l{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(x_1M_1^2x_2M_2^2)^m}{m!}}.`$ (34)
Using the identities
$`({\displaystyle \underset{i=0}{\overset{M}{}}}x_iA_i)^N=\mathrm{\Gamma }(N+1){\displaystyle \underset{k=1}{\overset{M}{}}}{\displaystyle \underset{i_k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{x_{i_k}A_{i_k}}{i_k!}}\text{with}{\displaystyle \underset{k=1}{\overset{M}{}}}i_k=N,`$ (35)
one can write
$`I_2^n`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{(1)^{\nu _i}}{\mathrm{\Gamma }(\nu _i)}}{\displaystyle _0^{\mathrm{}}}𝑑x_ix_i^{\nu _i1}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(x_1x_2q^2)^l}{l!}}\mathrm{\Gamma }(1ln/2){\displaystyle \underset{p_1=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{p_2=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{x_1^{p_1}}{p_1!}}{\displaystyle \frac{x_2^{p_2}}{p_2!}}`$ (36)
$`{\displaystyle \underset{m_1=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m_2=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(x_1M_1^2)^{m_1}}{m_1!}}{\displaystyle \frac{(x_2M_2^2)^{m_2}}{m_2!}}\text{with}p_1+p_2=ln/2.`$
Equating Eqs. (31) and (36) one obtains the result
$`I_2^n(\nu _1,\nu _2,q^2,M_1^2,M_2^2)`$ $`=`$ $`{\displaystyle \underset{l,p_i,m_i}{}}{\displaystyle \frac{\mathrm{\Gamma }(1ln/2)\mathrm{\Gamma }(1\nu _1)\mathrm{\Gamma }(1\nu _2)}{\mathrm{\Gamma }(1+l)\mathrm{\Gamma }(1+p_1)\mathrm{\Gamma }(1+p_2)\mathrm{\Gamma }(1+m_1)\mathrm{\Gamma }(1+m_2)}}`$ (37)
$`(q^2)^{l+m_1+m_2}\times \left({\displaystyle \frac{M_1^2}{q^2}}\right)^{m_1}\left({\displaystyle \frac{M_2^2}{q^2}}\right)^{m_2},`$
with
$`p_1+p_2=ln/2p_1=l\nu _1m_1p_2=l\nu _2m_2.`$ (38)
Choosing two independent summation indices out of five i.e. $`l,p_1,p_2,m_1,m_2`$ one can express $`I_2^n`$ into the following special functions . given by the Hyper-geometric Functions $`{}_{2}{}^{}F_{1}^{}`$, $`{}_{3}{}^{}F_{2}^{}`$, the Appell Functions $`F_i`$ ($`i=14`$), the Horn Function $`H_2`$ and the Kampé de Fériet Functions $`S_i`$ ($`i=1,2`$). In the case of the two-point function in Eq. (26) there are eight possibilties to express the integral into the special functions above depending on the convergence of the sums in a specific kinematic region. If we choose $`|M_i^2/q^2|<1`$ one has to take $`m_1,m_2`$. Using the following identity
$`{\displaystyle \frac{\mathrm{\Gamma }(zn)}{\mathrm{\Gamma }(z)}}=(1)^n{\displaystyle \frac{\mathrm{\Gamma }(1z)}{\mathrm{\Gamma }(1z+n)}},`$ (39)
the expression in Eq. 37 can be written as
$`I_2^n(\nu _1,\nu _2,q^2,M_1^2,M_2^2)={\displaystyle \underset{m_1=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m_2=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(1n+\nu _1+\nu _2+m_1+m_2)}{\mathrm{\Gamma }(n/2+1\nu _1\nu _2m_1m_2)}}`$
$`\times {\displaystyle \frac{\mathrm{\Gamma }(1\nu _1)\mathrm{\Gamma }(1\nu _2)}{\mathrm{\Gamma }(1+m_1)\mathrm{\Gamma }(1+m_2)}}(q^2)^{n/2\nu _1\nu _2}\times \left({\displaystyle \frac{M_1^2}{q^2}}\right)^{m_1}\left({\displaystyle \frac{M_2^2}{q^2}}\right)^{m_2}`$
$`=(1)^{n/2}(q^2)^{n/2\nu _1\nu _2}{\displaystyle \frac{\mathrm{\Gamma }(n/2\nu _1)\mathrm{\Gamma }(n/2\nu _2)\mathrm{\Gamma }(\nu _1+\nu _2n/2)}{\mathrm{\Gamma }(\nu _1)\mathrm{\Gamma }(\nu _2)\mathrm{\Gamma }(n\nu _1\nu _2)}}`$
$`F_4(1+\nu _1+\nu _2n,\nu _1+\nu _2n/2,1+\nu _1n/2,1+\nu _2n/2,{\displaystyle \frac{M_1^2}{q^2}},{\displaystyle \frac{M_2^2}{q^2}}).`$ (40)
Here the fourth Appell Function is given by
$`F_4(\alpha ,\beta ,\gamma ,\gamma ^{},x,y)`$ $`=`$ $`{\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(\alpha +m+n)}{\mathrm{\Gamma }(\alpha )}}{\displaystyle \frac{\mathrm{\Gamma }(\beta +m+n)}{\mathrm{\Gamma }(\beta )}}{\displaystyle \frac{\mathrm{\Gamma }(\gamma )}{\mathrm{\Gamma }(\gamma +m)}}{\displaystyle \frac{\mathrm{\Gamma }(\gamma ^{})}{\mathrm{\Gamma }(\gamma ^{}+n)}}`$ (41)
$`\times {\displaystyle \frac{x^m}{\mathrm{\Gamma }(m+1)}}{\displaystyle \frac{y^n}{\mathrm{\Gamma }(n+1)}}.`$
We conclude that using the method above one can find different representations for an one-loop integral in different kinematic regions depending on the radius of convergence of the ratios between the scales $`q^2,M_1^2,M_2^2`$. It is unclear whether this method, which seems to us very complicate, can be used to compute two-loop integrals.
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# 1 Introduction
## 1 Introduction
This note is intended to clarify the realization and interpretation of the Randall-Sundrum compactification scenario within string theory. In the model of , our 4-d world is extended with an extra direction $`r`$ to a 5-d space-time with the warped metric
$$ds^2=e^{2\sigma (r)}\eta _{\mu \nu }dx^\mu dx^\nu +dr^2$$
(1)
with $`\sigma (r)=k|r|`$. Even while the range of $`r`$ is infinite, the warped form of the metric ensures that the effective volume of the extra direction is finite. As a result, matter particles sufficiently close to the domain wall region near $`r=0`$ will experience ordinary 4-d gravity at long distances .
At first sight, this proposal seems like a rather drastic departure from the more conventional Kaluza-Klein framework. Indeed, in most works on string compactifications thus far, the four uncompactified directions and the compact manifold are assumed to form a simple direct product. Although it was realized for a long time that this basic KK set-up can be generalized to include the possibility of warped products, the physics of these more general scenarios is still largely unexplored.
A second important ingredient of the RS-scenario is that part or all of the observable matter may be thought of as confined to a 4-d sub-manifold of the higher dimensional space-time. A concrete theoretical realization of such world-branes are the D3-branes of IIB string theory, which confine open strings to their world-volume. D3-branes, however, do not bind 4-d gravity. Possible supersymmetric realizations of the Planck-brane, located around $`r=0`$ in (1), are therefore rather expected to be found in the form of domain wall type configurations, or possible stringy generalizations thereof. Various attempts have been made to find smooth domain wall solutions of this type within 5-d gauged supergravity, but thus far without real success .
There are several reasons for why this is indeed a hard problem. Even for a given compactification from 10 dimensions, it is an elaborate task to derive the dimensionally reduced theory. Thus far this has been done only for reductions over rather special symmetric 5-manifolds $`K_5`$ such as $`S^5`$ or $`S^5/Z_2`$, etc, and/or for special theories with extended supersymmetry. However, while it seems feasible to classify the possible types of supersymmetric solutions for each of these special dimensional reductions, there is no guarantee that they provide a general enough framework.
Instead of following the above procedure of $`(1)`$ performing some special dimensional reduction to 5 dimensions and $`(2)`$ looking for RS domain wall type solutions, it seems more practical to reverse the two steps. Since the scalar fields $`\varphi ^a`$ arise as moduli of some internal 5-d compact space $`K_5`$, any domain wall solution in 5-d gauged supergravity describes (upon lifting it back up to 10-dimensions) some specific warped compactification of the 10-d theory. It will therefore be much more general – and also easier – to first (a) identify a general class of warped compactifications of the 10-dimensional theory, and then (b) perform the same type of dimensional reduction from 10 to 5 dimensions. In the end, one can then hope to identify a class of 10-d geometries for which the resulting dimensionally reduced solution has all the required properties.
As will be described below, such a class of warped IIB geometries indeed exists in the form of quite generic F-theory compactifications on Calabi-Yau four-folds.<sup>1</sup><sup>1</sup>1Another special realisation of an RS geometry in terms of a toroidal type IIB orientifold compactification has been described in . These have been studied in some detail in the recent literature – a list of references include and – and indeed none of our equations will be new. Given the current interest in the subject, however, it seems useful to collect some of the known facts about these compactifications, since it has not been generally appreciated that supersymmetric RS-geometries indeed exist in string theory, and furthermore that they are in fact quite generic.
Since all derivations are contained in existing papers, we will here only present the general form of the compactification geometry without any proof that it is really a supersymmetric solution to the 10-d equations of motion. This proof can however be quite directly extracted from the literature, in particular from the very clear discussion by Becker and Becker . Their analysis was done in the context of M-theory compactifications on C-Y four-folds. It can however be straightforwardly translated to the F-theory context by performing the T-duality transformation outlined in . An explicit example of this T-duality transformation is discussed in .
Although the 10-d perspective will allow us to identify a large class of RS-type compactification geometries, their geometrical structure is rather involved. It is therefore not easy to explicitly perform the dimensional reduction of these solutions to 5-dimensions. We will nonetheless attempt to make this 5-d perspective as transparent as possible. In particular we will show that they indeed give rise to a 5-d metric of the generic form (1).
## 2 Warped Compactification in F-theory
F-theory is a geometric language for describing compactifications of type IIB string theory, in which the expectation values of the dilaton and axion fields are allowed to vary non-trivially along the compactification manifold . Compactifications of F-theory down to four-dimensions are specified by means of a Calabi-Yau four-fold that admit an elliptic fibration with a section. In other words, these are 8 dimensional compact manifolds $`K_8`$ that locally look like a product of a complex three-fold $`K_6`$ times a two-torus $`T^2`$. The two-torus will be taken to shrink to zero size. It can however be taken to change its shape when moving along the base $`K_6`$. In particular, it can have non-trivial monodromies around singular co-dimension 2 loci inside the $`K_6`$, where the elliptic fiber degenerates.
The four-fold $`K_8`$ is not the actual compactification geometry; rather it gives an economical way to characterize the compactification geometry as well as the expectation values of other fields. Moreover, due to the special geometric properties of the Calabi-Yau four-fold $`K_8`$ – vanishing first Chern class and $`SU(4)`$ holonomy – the associated IIB background by construction will preserve 4-d supersymmetry, at least at the classical and perturbative level.
The warped geometry of this type of F-theory compactifications has been derived in , by direct translation of the M-theory analysis of . The full solution for the 10-dimensional IIB string metric, in the string frame, takes the form
$$ds_{_{IIB}}^2=e^{2\alpha (y)}\eta _{\mu \nu }dx^\mu dx^\nu +e^{2\alpha (y)}g_{^_\mathrm{M}^_\mathrm{N}}(y)dy^^_\mathrm{M}dy^^_\mathrm{N}$$
(2)
where $`g_{^_\mathrm{M}^_\mathrm{N}}`$ denotes the metric on the 6-dimensional base-manifold $`K_6`$. The shape of the warp-factor $`e^{2\alpha }`$ will depend on the detailed geometry of the CY four-fold $`K_8`$, as well as on other data such as the possible non-zero expectation values of other fields and the locations of the possible D-branes.
Besides the ten-dimensional space-time metric, the fields that can take non-trivial expectation values are the following:
(i) the dilaton field $`\varphi `$
(ii) the RR-scalar or axion field $`\stackrel{~}{\varphi }`$
(iii) the NSNS 3-form field strength $`H^{NS}`$
(iv) the RR 3-form field strength $`H^R`$
(v) the RR 5-form field strength $`F^{(5)}`$
The expectation values of all these fields can all be conveniently characterized in terms of the geometry of $`K_8`$.
In the type IIB theory, the modulus $`\tau `$ of the elliptic fibration, the shape of the two-torus inside the $`K_8`$, describes the variation along the 6-d base manifold $`K_6`$ of the dilaton and axion fields, $`\varphi `$ and $`\stackrel{~}{\varphi }`$, via the identification
$$\tau =\stackrel{~}{\varphi }+ie^\varphi .$$
(3)
As mentioned above, a key feature of F-theory is that this modulus in general has non-trivial monodromies around 4-d submanifolds inside $`K_6`$. These 4-d sub-manifolds are associated with the locations of D7-branes, of which the remaining 3+1-dimensions span the uncompactified space-time directions. In going around one of the D7-branes, the modulus field $`\tau `$ can pick up an $`SL(2,𝐙)`$ monodromy
$$\tau \frac{a\tau +b}{c\tau +d},$$
(4)
which leaves the geometric shape of the two-torus fibre inside $`K_8`$ invariant, but nonetheless via (3) amounts to a non-trivial duality transformation of the IIB string theory. We thus notice that the dilaton and axion field are not smooth single-valued functions, but instead are multi-valued with branch cut singularities at the locations of the D7-branes. The full non-perturbative string theory, however, is expected to be well-behaved everywhere.
For the following, it will be convenient to combine the NSNS and RR three-from field strengths, $`H^{NS}`$ and $`H^R`$, of the IIB supergravity into a single four-form field-strength $`G`$ on $`K_8`$ as follows . Let $`z`$ and $`\overline{z}`$ denote the coordinates along the $`T^2`$ fiber. Then we can write
$$G=\frac{\pi }{i\tau _2}(Hd\overline{z}\overline{H}dz)$$
(5)
$$H=H^R\tau H^{NS};\overline{H}=H^R\overline{\tau }H^{NS}$$
(6)
For supersymmetric configurations, $`H`$ defines an integral harmonic (1,2)-form on $`K_6`$ satisfying $`Hk=0`$ with $`k`$ the Kähler class of $`K_6`$ . It transforms under the $`SL(2,Z)`$ monodromy transformations (4) around the seven-branes as $`HH/(c\tau +d)`$. The field-strength $`G`$ is invariant under these transformations.
An important aspect of F-theory compactifications is that they typically carry, via their non-trivial topology, an effective total D3-brane charge. The value of this charge is proportional to the Euler characteristic $`\chi (K_8)`$ of the original Calabi-Yau four-fold $`K_8`$. Here $`\chi (K_8)`$ is defined via
$$\frac{1}{24}\chi (K_8)=_{K_8}I_8(R)$$
(7)
where
$$I_8(R)=\frac{1}{192}\left(\mathrm{tr}R^4\frac{1}{4}(\mathrm{tr}R^2)^2\right)$$
(8)
with $`R`$ the curvature two-form on $`K_8`$. Global tadpole cancellation, or conservation of the RR 5-form flux, requires that this charge must be canceled by other sources. These other sources come from possible non-zero fluxes of the NSNS or RR two-form fields, or from the explicit insertion of $`N`$ D3-brane world branes, that is, D3-branes that span the 3+1-d uncompactified world but are localized as point-like objects inside the $`K_6`$. The number of such D3-branes is therefore not free, but completely determined via charge conservation. This global tadpole cancellation relation reads
$$N=\frac{1}{24}\chi (K_8)\frac{1}{8\pi ^2}_{K_8}GG$$
(9)
Depending on the topology of $`K_8`$, $`N`$ can reach values of up to $`10^3`$ or larger. An example with $`N=972`$, mentioned in , is provided by an elliptically fibered CY four-fold over $`𝐏^3`$. The Euler number $`\chi (K_8)`$ can be non-zero only if $`K_6`$ has a non-vanishing first Chern class, that is, provided the F-theory compactification makes use of a non-zero number of D7-branes.
The equation of motion for the warp factor $`e^{2\alpha }`$ obtained in and reads as follows
$$\mathrm{\Delta }^{(8)}e^{4\alpha }=\mathrm{\hspace{0.33em}4}\pi ^2\{I_8(R)\frac{1}{8\pi ^2}GG\underset{i=1}{\overset{N}{}}\delta ^{(8)}(yy_i)\}$$
(10)
where $`\mathrm{\Delta }^{(8)}`$ denotes the Laplacian and $``$ the Hodge star on $`K_8`$. The points $`y=y_i`$ correspond to the location of the $`N`$ D3-branes. Here, following , we have written the equation on the full 8-d manifold $`K_8`$, even though in F-theory the elliptic fiber $`T^2`$ inside $`K_8`$ has been shrunk to zero size. In this limit, the solution for $`\alpha `$ obtained via (10) only depends on the 6 coordinates $`y^^_\mathrm{M}`$ on $`K_6`$. Alternatively, using the analysis of , one may also first reduce the right-hand side to $`K_6`$, via integration over the $`T^2`$ fiber, and then solve the reduced equation to obtain the function $`e^{4\alpha }`$ directly on $`K_6`$.
Finally, there is also an non-trivial expectation value for the self-dual RR five-form field strength, equal to
$$F_{\mu \nu \lambda \sigma ^_\mathrm{M}}=ϵ_{\mu \nu \lambda \sigma }_^_\mathrm{M}e^{4\alpha }.$$
(11)
We note that via (10) the D3-branes indeed form a source for this field strength, but that via (9) the total charge adds up to zero.
## 3 Shape of the warp factor
Let us summarize. Starting from the elliptically fibered CY four-fold $`K_8`$ we can extract a complete characterization of the warped compactification. First, since $`K_8K_6\times T^2`$, we obtain the metric $`g_{^_\mathrm{M}^_\mathrm{N}}`$ on the base $`K_6`$, as well as the dilaton and axion via (3). We then deduce the form of the warp factor $`e^{2\alpha }`$ from (10), which incorporates the complete backreaction due to the $`G`$-flux and D3-branes. Finally from (2), we obtain the actual compactification geometry. Note that, as indicated in fig 2, the rescaling by $`e^{2\alpha }`$ of $`g_{^_\mathrm{M}^_\mathrm{N}}`$ in (2) may have a drastic effect on the shape of the compactification manifold, which indeed may look quite different from that of the original $`K_6`$. In particular, it is possible that near the locations of the D3-branes one of the internal directions may become non-compact.
We may formally solve the equation (10) via
$$e^{4\alpha (y)}=e^{4\alpha _0}+4\pi ^2d^8y^{}\sqrt{g}𝒢(y,y^{})\left[I_8(R(y^{}))\frac{1}{8\pi ^2}GG\underset{i=1}{\overset{N}{}}\delta ^{(8)}(y^{}y_i)\right]$$
(12)
where $`𝒢(y,y^{})`$ denotes the Green function for $`\mathrm{\Delta }^{(8)}`$. The term $`e^{4\alpha _0}`$ parametrizes the constant zero mode of $`e^{4\alpha }`$, which is not fixed by eqn (10). Note that for $`e^{4\alpha }`$ to be everywhere positive, this constant $`e^{4\alpha _0}`$ can not be arbitrarily small, since the second term on the r.h.s. of (12) can become negative. This implies that the warped 6-geometry automatically has a minimal volume<sup>2</sup><sup>2</sup>2We thank Sav Sethi for bringing this feature to our attention..
An interesting limiting case is when all $`D3`$ branes are concentrated in one point, say $`y=y_0`$. Close to this point, the warp function $`\alpha (y)`$ reduces to
$$\alpha (y)\mathrm{log}|yy_0|+\mathrm{const}.$$
(13)
Via (2) this describes the familiar semi-infinite near-horizon geometry of $`N`$ D3-branes: $`AdS_5\times S^5`$ with radius $`R={}_{}{}^{4}\sqrt{4\pi Ng_s}`$. (See fig 2.) Although the radial $`AdS_5`$ coordinate $`rR\mathrm{log}|yy_0|`$ runs over semi-infinite range, the compactification geometry (2) still gives rise to a 4-d Einstein action with a finite 4-d Newton constant $`1/\mathrm{}_4^2`$ equal to
$$\frac{1}{(\mathrm{}_4)^2}=\frac{1}{(\mathrm{}_{_{10}})^8}_{K_6}\sqrt{g}e^{2\varphi 4\alpha }$$
(14)
with $`\mathrm{}_{10}`$ the 10-d Planck length.
## 4 Reduction to 5 dimensions.
We would now like to show that these F-theory compactifications, upon performing a suitable dimensional reduction to five dimensions, reduce to supersymmetric RS domain wall solutions. To this end we will look for a specific coordinate system
$$y^^_\mathrm{M}=(y^m,r)$$
(15)
where $`m`$ now runs over 5 values, such that the 10-d metric takes the following form
$$ds_{RS}^2=e^{2\sigma (r)}\eta _{\mu \nu }dx^\mu dx^\nu +dr^2+h_{mn}(y,r)dy^mdy^n$$
(16)
where $`ds_{RS}^2`$ related to the original IIB string metric (2) via the rescaling
$$ds_{RS}^2=e^{\varphi /2}(V_5)^{1/4}ds_{_{IIB}}^2V_5=\frac{1}{(\mathrm{}_{10})^5}_{K_5}\sqrt{h}$$
(17)
Here the prefactor in (17) is chosen such that $`ds_{RS}^2`$ is the metric in the 5-d Einstein frame (where we have taken the 5-d Planck length $`\mathrm{}_5`$ equal to the 10-d one). The 5-d slices of constant $`r`$ define 5-d submanifolds $`K_5`$, on which $`2\alpha (y)+\varphi (y)/2=`$ constant; this correspondence guarantees that the warp-factor $`e^{2\sigma }`$ in $`ds_{RS}^2`$ just depends on $`r`$ and not on the remaining $`y^m`$’s.
In this way we indeed obtain a solution that from the 5-d perspective looks just like an RS-type warped geometry. For large negative $`r`$, close to the D3-branes, the warp factor behaves like $`e^{2\sigma }e^{2|r|/R}`$ with $`R`$ the AdS-radius of the $`N`$ D3-brane solution. On the other end, somewhere outside the throat region of the AdS-tube, near the ‘equator’ of the $`K_6`$, the warp factor $`e^{2\sigma (r)}`$ reaches some maximal value. Eventually, there is a boundary value for the coordinate $`r`$, which we can take to be $`r=0`$, at which the transverse 5-manifold $`K_5`$ shrinks to zero size. However, as is clear from fig 2, this does not correspond to any singularity, but just to a smooth cap closing off the 6-d manifold $`K_6`$. This fact that $`r`$ takes a maximal value $`r=0`$ implies that the AdS-space is indeed compactified, in the sense that, relative to the 4-d space-time, it has a finite volume. It therefore produces to a finite 4-d Newton constant equal to (cf )
$$\frac{1}{(\mathrm{}_4)^2}=\frac{1}{(\mathrm{}_5)^3}\underset{\mathrm{}}{\overset{0}{}}𝑑re^{2\sigma (r)}.$$
(18)
It is easily verified that via (17) and (16), this result coincides with the value (14) found via direct reduction from 10 dimensions.
Upon dimensional reduction, the metric $`h_{mn}`$ of the internal $`K_5`$, as well as other fields such as the dilaton, all reduce to scalar fields that provide the matter multiplets of the 5-d supergravity. All these fields vary with the radial coordinate $`r`$, and since this radial flow is supersymmetric, it should in principle be described as some gradient flow driven by some appropriate superpotential. However, due to the rather complex geometrical structure of the typical F-theory compactification, it unfortunately seems impossibly hard to find the explicit form of this potential.
## 5 Discussion
In this note we have summarized the geometrical description of warped F-theory compactifications, and shown they can be used to obtain geometries very analogous to the RS-scenario. Due to the presence of the D7-branes in the F-theory geometry, the solutions are not completely smooth: the dilaton and axion field have isolated branch cut singularities at the D7-brane loci around which $`\tau `$ is multi-valued. The string theory, however, is well-defined in this background.
Our description of the solutions makes clear that the internal structure of the RS Planck brane is not that of an actual brane, but rather that of a compactification geometry (which however also contains the D7-branes). Consequently, the localized graviton zero mode of is just the standard KK zero mode of the 10-d metric; because its wave-function is sharply peaked near the wall region, where the warp factor $`e^{2\sigma }`$ is maximal, the 4-d graviton indeed looks like some bound state. From the higher dimensional viewpoint, however, this RS-localization of gravity is not a new phenomenon.
The most interesting aspect of these warped string compactifications, is that there is no clear distinction between the low energy and extra dimensional physics. Kaluza-Klein excitations, when localized far inside the AdS region can describe particles with masses much smaller than the inverse size of the original $`K_6`$. While these particles naively look like new degrees of freedom arising from the presence of extra dimensions, the holographic AdS/CFT correspondence tells us that they are in fact localized excitations of the low energy gauge theory. The same holds for string excited modes in this region. Hence via the holographic identification of the RG scale with a real extra dimension, the two usually separate stages of dimensional and low energy reduction should now be combined into one single procedure.
Finally, as a closely related point, we need to emphasize that the solutions as given here are generically unstable against small perturbations. The best way to understand this instability is via the RG language: in general there will exist relevant operators whose couplings, once turned on in the UV region, will quickly grow and typically produce some singularity that effectively closes off the AdS-tube . In our way of obtaining the solutions, we did not immediately notice this instability, because we required that the original $`K_6`$ geometry and all other fields, except for the warp-factor and the 5-form flux, are smooth at the locations of the D3-branes. This requirement is special, however, and we should allow for deformations that may spoil this property.
In order to ensure that the perturbed geometry contains a substantial intermediate AdS-like region, we either need to fine-tune the UV initial conditions or introduce a symmetry that eliminates these unstable modes. In RG language, this means that the dual 4-d field theory should be made approximately conformal invariant over a large range of scales, separating the Planck scale from the scale set by the non-trivial gauge dynamics. Via the AdS/CFT dictionary, the problem of realizing an RS-geometry in string theory therefore is reduced back to the original problem it was designed to solve, namely how to generate a large gauge hierarchy. Or stated in more positive terms, in searching for realistic string compactification scenarios, the observed gauge hierarchy can be viewed as an indication that warped geometries of this type deserve serious attention.
Acknowledgements
This work is supported by NSF-grant 98-02484. C.S.C is supported in part by an NSF Graduate Research Fellowship. We would like to thank M. Berkooz, M. Cvetic, R. Dijkgraaf, M. Gremm, R. Kallosh, L. Randall, S. Sethi and E. Verlinde for helpful discussions.
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# REFERENCES
hep-th/0003045
Gravity and antigravity in a brane world with metastable gravitons.
Comment on: G. Dvali, G. Gabadadze and M. Porrati, “Metastable gravitons and infinite volume extra dimensions,” hep-th/0002190
and
C. Csaki, J. Erlich, T. J. Hollowood, “Graviton propagators, brane bending and bending of light in theories with quasi-localized gravity,” hep-th/0003020
Ruth Gregory<sup>1</sup>, Valery A. Rubakov<sup>2</sup> and Sergei M. Sibiryakov<sup>2</sup>
<sup>1</sup> Centre for Particle Theory, Durham University, South Road, Durham, DH1 3LE, U.K.
<sup>2</sup> Institute for Nuclear Research of the Russian Academy of Sciences,
60th October Anniversary prospect, 7a, Moscow 117312, Russia.
## Abstract
In the framework of a five-dimensional three-brane model with quasi-localized gravitons we evaluate metric perturbations induced on the positive tension brane by matter residing thereon. We find that at intermediate distances, the effective four-dimensional theory coincides, up to small corrections, with General Relativity. This is in accord with Csaki, Erlich and Hollowood and in contrast to Dvali, Gabadadze and Porrati. We show, however, that at ultra-large distances this effective four-dimensional theory becomes dramatically different: conventional tensor gravity changes into scalar anti-gravity.
The papers by Dvali, Gabadadze and Porrati, , and Csaki, Erlich and Hollowood, , address the issue of whether four-dimensional gravity is phenomenologically acceptable in a class of brane models with infinite extra dimensions in which the five-dimensional gravitons have a metastable “bound state”, rather than a genuine zero mode. A model of this sort has been proposed in Refs. and is a variation of the Randall–Sundrum (RS) scenario for a non-compact fifth dimension . The construction with metastable gravitons has been put in a more general setting in Refs. . It has been argued in Ref. that models with metastable gravitons are not viable: from the four-dimensional point of view, gravitons are effectively massive and hence appear to suffer from a van Dam–Veltman–Zakharov discontinuity in the propagator in the massless limit. In particular, it has been claimed that the prediction for the deflection of light by massive bodies would be considerably different from that of General Relativity. The issue has recently been analyzed in more detail in Ref. , where explicit calculations of four dimensional gravity have been performed along the lines of Garriga and Tanaka , and Giddings, Katz and Randall . The outcome of that analysis is that four-dimensional gravity has been claimed to be in fact Einsteinian, despite the peculiarity of apparently massive gravitons.
In this comment we also apply the Garriga–Tanaka (GT) technique to obtain effective four-dimensional gravity at the linearized level, considering as an example the model of Refs. . We find that at intermediate distances (which should extend from microscopic to very large scales in a phenomenologically acceptable model) four dimensional gravity is indeed Einsteinian, in accord with Ref. and in contrast to Ref. . However, at ultra-large scales we find a new phenomenon: four-dimensional gravity changes dramatically, becoming scalar anti-gravity rather than tensor gravity. This may or may not signal an internal inconsistency of the models under discussion.
To recapitulate, the set up of Refs. is as follows. The model has five dimensions and contains one brane with tension $`\sigma >0`$ and two branes with equal tensions $`\sigma /2`$ placed at equal distances to the right and to the left of the positive tension brane in the fifth direction. There is a reflection symmetry, $`zz`$, which enables one to consider explicitly only the region to the right of the positive tension brane (hereafter $`z`$ denotes the fifth coordinate). Conventional matter resides on the central positive tension brane. The bulk cosmological constant between the branes, $`\mathrm{\Lambda }`$, is negative, whereas it is equal to zero to the right of the negative tension brane. With appropriately tuned $`\mathrm{\Lambda }`$, there exists a solution to the five-dimensional Einstein equations for which both positive and negative tension branes are at rest at $`z=0`$ and $`z=z_c`$ respectively, $`z_c`$ being an arbitrary constant. The metric of this solution is
$$ds^2=a^2(z)\eta _{\mu \nu }dx^\mu dx^\nu dz^2$$
(1)
where
$$a(z)=\{\begin{array}{cc}e^{kz}\hfill & 0<z<z_c\hfill \\ e^{kz_c}a_{}\hfill & z>z_c\hfill \end{array}$$
(2)
The constant $`k`$ is related to $`\sigma `$ and $`\mathrm{\Lambda }`$. The four-dimensional hypersurfaces $`z=const.`$ are flat, the five-dimensional space-time is flat to the right of the negative-tension brane and anti-de Sitter between the branes. The spacetime to the left of the positive tension brane is a mirror image of this set-up.
This background has two different length scales, $`k^1`$ and
$$r_c=k^1e^{3kz_c}$$
(3)
These are assumed to be well separated, $`r_ck^1`$. It has been argued in Ref. that the extra dimension “opens up” both at short distances, $`rk^1`$ and ultra-long ones, $`rr_c`$.
To find the four-dimensional gravity experienced by matter residing on the positive tension brane, we follow GT and consider a Gaussian-Normal (GN) gauge
$$g_{zz}=1g_{z\mu }=0$$
(4)
In the bulk, one can further restrict the gauge to be transverse-tracefree (TTF)
$$h_\mu ^\mu =h_{\nu ,\mu }^\mu =0$$
(5)
Hereafter $`h_{\mu \nu }`$ are metric perturbarions; indices are raised and lowered by the four-dimensional Minkowski metric. The linearized Einstein equations in the bulk take one and the same simple form for all components of $`h_{\mu \nu }`$,
$$\{\begin{array}{cc}h^{\prime \prime }4k^2h\frac{1}{a^2}\mathrm{}^{(4)}h=0\hfill & 0<z<z_c\hfill \\ h^{\prime \prime }\frac{1}{a_{}^2}\mathrm{}^{(4)}h=0\hfill & z>z_c\hfill \end{array}$$
(6)
It is convenient, however, to formulate the junction conditions on the positive tension brane in the local GN frame. In this frame, metric perturbations $`\overline{h}_{\mu \nu }`$ are not transverse-tracefree, so the two sets of perturbations are related in the bulk between the two branes by a five-dimensional gauge transformation preserving (4),
$$\overline{h}_{\mu \nu }=h_{\mu \nu }+\frac{1}{k}\widehat{\xi }_{,\mu \nu }^52ka^2\eta _{\mu \nu }\widehat{\xi }^5+a^2(\xi _{\mu ,\nu }+\xi _{\nu ,\mu })$$
(7)
where $`\xi _\mu (x)`$ and $`\widehat{\xi }^5(x)`$ are the gauge parameters. Notice that if $`\xi ^5`$ is not zero, there is a ‘shift’ in the location of the wall relative to an observer at infinity, i.e. the wall appears bent to such an observer (as discussed in ). Physically, this simply represents the fact that the wall GN frame is constructed by integrating normal geodesics from the wall, and in the presence of matter these geodesics will be distorted, thereby altering the proper distance between the wall and infinity. In fact, one finds a similar “bending” of the equatorial plane in the Schwarzschild spacetime if one tries to impose a local GN frame away from the horizon.
In the presence of additional matter on the positive tension brane with energy momentum $`T_{\mu \nu }`$, the junction conditions on this brane read
$$\overline{h}_{\mu \nu }^{}+2k\overline{h}_{\mu \nu }=8\pi G_5\left(T_{\mu \nu }\frac{1}{3}\eta _{\mu \nu }T_\lambda ^\lambda \right)$$
(8)
where $`G_5`$ is the five-dimensional gravitational constant. The solution to equations (5) – (8) has been obtained by Garriga and Tanaka. They found that $`\widehat{\xi }^5`$ obeys
$$\mathrm{}^{(4)}\widehat{\xi }^5=\frac{4\pi }{3}G_5T_\lambda ^\lambda $$
(9)
We will need the expression for the induced metric on the positive tension brane. Up to terms that can be gauged away on this brane, the induced metric is
$$\overline{h}_{\mu \nu }(z=0)=h_{\mu \nu }^{(m)}2k\eta _{\mu \nu }\widehat{\xi }^5$$
(10)
where
$$h_{\mu \nu }^{(m)}=16\pi G_5𝑑x^{}G_R^{(5)}(x,x^{};z=z^{}=0)\left(T_{\mu \nu }\frac{1}{3}\eta _{\mu \nu }T_\lambda ^\lambda \right)(x^{})$$
(11)
Here $`G_R^{(5)}`$ is the retarded Green’s function of eq. (6) with appropriate (source-free) junction conditions on the two branes. This Green’s function is mirror-symmetric and obeys
$$\left[_z^24k^2\theta (z_cz)\frac{1}{a^2}\mathrm{}^{(4)}+4k\delta (z)2k\delta (zz_c)\right]G_R^{(5)}(x,x^{};z,z^{})=\delta (xx^{})\delta (zz^{})$$
(12)
Let us consider the case of the static source first. It has been found in Ref. that for $`k^1rr_c`$, the leading behavior of the static Green’s function (given by $`𝑑tG_R^{(5)}(z=z^{}=0)`$) is the same as in the RS model (up to small corrections), and corresponds to a $`1/r`$ potential. Hence, at intermediate distances the analysis is identical to GT, and the induced metric is the same as in the linearized four-dimensional General Relativity. This is in accord with Ref. .
On the other hand, it follows from Ref. that at ultra-large distances, $`rr_c`$, the contribution (11) behaves like $`1/r^2`$ (the fifth dimension “opens up”). There remains, however, the second term in eq.(10). Since eq.(9) has a four-dimensional form, this term gives rise to a $`1/r`$ potential (missed in Ref. ) even at ultra-large distances. For a point-like static source of unit mass, the corresponding gravitational potential is
$$V(r)\frac{1}{2}\overline{h}_{00}(r)=+\frac{1}{3}G_4\frac{1}{r}$$
(13)
where $`G_4=kG_5`$ is the four-dimensional Newton’s constant entering also into the conventional Newton’s law at intermediate distances. We see that at $`rr_c`$, four-dimensional gravity is induced by the trace of energy-momentum tensor and has a repulsive $`1/r`$ potential. At ultra-large distances tensor gravity changes to scalar anti-gravity.
Likewise, the four-dimensional gravitational waves emitted by non-static sources are conventional tensor ones at intermediate distances and transform into scalar waves at ultra-large distances (the relevant distance scale being different from $`r_c`$ due to relativistic effects, see Ref. ). Indeed, the first term in eq. (10) dissipates , whereas the second term survives, again due to the four-dimensional structure of eq. (9).
These two cases illustrate the general property of eq. (10): the first term becomes irrelevant at ultra-large distances (the physical reason being the metastability of the five-dimensional graviton bound state), so the four-dimensional gravity (in effect, anti-gravity) is entirely due to the second, scalar term.
This bizarre feature of models with a metastable graviton bound state obviously deserves further investigation. In particular, it will be interesting to identify the four-dimensional massless scalar mode, which is present at ultra-large distances, among the free sourceless perturbations. This mode is unlikely to be the radion , studied in this model in Ref. : the radion would show up at intermediate distances, as well as at ultra-large ones<sup>*</sup><sup>*</sup>*The radion presumably couples exponentially weakly to the matter on the positive tension brane, as it does in the two-brane model of Randall and Sundrum . This may be the reason why the radion effects have not been revealed by the analyses made in Ref. and this note.; furthermore, the experience with models where the distance between the branes is stabilized suggests that the massless four-dimensional mode parametrized by $`\widehat{\xi }^5`$ exists even if the radion is made massive.
More importantly, one would like to understand whether anti-gravity at ultra-large distances is a signal of an intrinsic inconsistency of this class of models, or simply a signal that physics is intrinsically five-dimensional at these scales. In four dimensions, scalar antigravity requires either negative kinetic and gradient energy or a ghost. Whether or not a similar feature is inherent in models with extra dimensions remains an open question. If it is, there still would remain a possibility that fields with negative energy might be acceptable, as their effect might show up at ultra-large distances only.
We note finally, that anti-gravity may not be a special feature of models with quasi-localized gravitons. It is also possible that this phenomenon may be present in models of the type suggested by Kogan et. al. , where some Kaluza–Klein graviton excitations are extremely light. The same question about internal consistency then would apply to these models as well.
We would like to thank Sergei Dubovsky, Dmitry Gorbunov, Maxim Libanov and Sergei Troitsky for useful discussions. We are indebted to C. Csaki, J. Erlich and T. J. Hollowood for sending their paper prior to publication. R.G. was supported in part by the Royal Society, and V.R. and S.S. by the Russian Foundation for Basic Research, grant 990218410.
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# Introduction
## Introduction
Recent efforts to understand the fundamental physics of hydrodynamic turbulence have concentrated on the explanation of the observed violations of Kolmogorov’s scaling. These violations reflect the occurrence of large fluctuations in the velocity field on the small scales, a phenomenon referred to as intermittency. Some progress in the understanding of intermittency has been achieved recently through the study of simple model problems that include Burgers equation and the passive advection of a scalar by a velocity field of known statistics . This paper is a summary of the many interesting mathematical issues that arise in the problem of passive scalar advection together with our understanding of these issues. We put some of our results in the perspective of a new phenomenological model proposed recently by Barenblatt and Chorin using the formalism of incomplete self-similarity.
## Generalized Flows
Consider the transport equation for the scalar field $`\theta ^\kappa (x,t)`$ in $`^d`$:
$$\frac{\theta ^\kappa }{t}+(u(x,t))\theta ^\kappa =\kappa \mathrm{\Delta }\theta ^\kappa .$$
(1)
We will be interested in $`\theta ^\kappa `$ in the limit as $`\kappa 0`$. It is known from classical results that if $`u`$ is Lipschitz continuous in $`x`$, then as $`\kappa 0`$, $`\theta ^\kappa `$ converges to $`\theta `$, the solution of
$$\frac{\theta }{t}+(u(x,t))\theta =0.$$
(2)
Furthermore, if we define $`\{\phi _{s,t}(x)\}`$ as the flow generated by the velocity field $`u`$, satisfying the ordinary differential equations (ODEs)
$$\frac{d\phi _{s,t}(x)}{dt}=u(\phi _{s,t}(x),t),\phi _{s,s}(x)=x,$$
(3)
for $`s<t`$, then the solution of the transport equation in 2 for the initial condition $`\theta ^\kappa (x,0)=\theta _0(x)`$ is given by
$$\theta (x,t)=\theta _0(\phi _{0,t}^1(x))=\theta _0(\phi _{t,0}(x)).$$
(4)
This classical scenario breaks down when $`u`$ fails to be Lipschitz continuous in $`x`$, which is precisely the case for fully developed turbulent velocity fields. In this case Kolmogorov’s theory of turbulent flows suggests that $`u`$ is only Hölder continuous with an exponent roughly equal to $`\frac{1}{3}`$ for $`d=3`$. In such situations the solution of the ODEs in 3 may fail to be unique , and we then have to consider probability distributions on the set of solutions in order to solve the transport equation in 2. This is the essence of the notion of generalized flows proposed by Brenier (see also ).
There are two ways to think about probability distributions on the solutions of the ODEs in 3. We can either think of it as probability measures on the path-space (functions of $`t`$) supported by paths which are solutions of 3, or we can think of it as transition probability at time $`t`$ if the starting position at time $`s`$ is $`x`$. In the classical situation when $`u`$ is Lipschitz continuous, this transition probability degenerates to a point mass centered at the unique solution of 3. When Lipschitz condition fails, this transition probability may be non-degenerate and the system in 3 is intrinsically stochastic.
There is a parallel story for the case when $`u`$ is a white-in-time random process defined on a probability space $`(\mathrm{\Omega },,𝒫)`$. We will denote the elements in $`\mathrm{\Omega }`$ by $`\omega `$ and indicate the dependence on realization of the random velocity field by a super- or a subscript $`\omega `$. In connection with the transport equation in 2, it is most natural to consider the stochastic ODEs
$$d\phi _{s,t}^\omega (x)=u(\phi _{s,t}^\omega (x),t)dt,\phi _{s,s}^\omega (x)=x,$$
(5)
in Stratonovich sense. In this case, it is shown that if the local characteristic of $`u`$ is spatially twice continuously differentiable, then the system in 5 has a unique solution. Such conditions are not satisfied by typical turbulent velocity fields on the scale of interest. When the regularity condition on $`u`$ fails, there are at least two natural ways to regularize 3 or 5. The first is to add diffusion:
$$d\phi _{s,t}^{\omega ,\kappa }(x)=u(\phi _{s,t}^{\omega ,\kappa }(x),t)dt+\sqrt{2\kappa }d\beta (t),$$
(6)
and consider the limit as $`\kappa 0`$. We will call this the $`\kappa `$-limit. The second is to smooth out the velocity field. Let $`\psi _\epsilon `$ be defined as $`\psi _\epsilon (x)=\epsilon ^d\psi (x/\epsilon ),`$ where $`\psi `$ is a standard mollifier: $`\psi 0`$, $`_^d\psi 𝑑x=1`$, $`\psi `$ decays fast at infinity. Let $`u^\epsilon =u\psi _\epsilon `$ and consider
$$d\phi _{s,t}^{\omega ,\epsilon }(x)=u^\epsilon (\phi _{s,t}^{\omega ,\epsilon }(x),t)dt,$$
(7)
in the limit as $`\epsilon 0`$. We will call this the $`\epsilon `$-limit. Physically $`\kappa `$ plays the role of molecular diffusivity, $`\epsilon `$ can be thought of as a crude model of the viscous cut-off scale. The $`\kappa `$-limit corresponds to the situation when the Prandtl number, defined here as the ratio of $`\epsilon `$ and $`\kappa `$, tends to zero, $`Pr0`$, whereas the $`\epsilon `$-limit corresponds to the situation when the Prandtl number diverges, $`Pr\mathrm{}`$. The following questions naturally arise:
* How do the flows and the passive scalar behave statistically in the $`\kappa `$\- and $`\epsilon `$-limits?
* Does there exist a unique statistical steady state when the transport equation in 1 is suitably forced?
* What are the statistical and geometrical properties of solutions in the statistical steady state?
Below we address these questions on a specific model introduced by Kraichnan .
Before proceeding further, we relate the regularized flows in 6, 7 to the solutions of the transport equations. Consider the $`\kappa `$-regularization first. It is convenient to introduce the backward transition probability
$$g_\omega ^\kappa (x,t|dy,s)=𝐄_\beta \delta (y\phi _{t,s}^{\omega ,\kappa }(x))dy,s<t,$$
(8)
where the expectation is taken with respect to $`\beta (t)`$, and $`\phi _{t,s}^{\omega ,\kappa }(x)`$ is the flow inverse to $`\phi _{s,t}^{\omega ,\kappa }(x)`$ defined in 6 (i.e. $`\phi _{s,t}^{\omega ,\kappa }(x)`$ is the forward flow and $`\phi _{t,s}^{\omega ,\kappa }(x)`$ is the backward flow). The action of $`g_\omega ^\kappa `$ generates a semi-group of transformation
$$S_{t,s}^{\omega ,\kappa }\psi (x)=_^d\psi (y)g_\omega ^\kappa (x,t|dy,s),$$
(9)
for all test functions $`\psi `$. $`\theta _\omega ^\kappa (x,t)=S_{t,s}^{\omega ,\kappa }\psi (x)`$ solve the transport equation in 1 for the initial condition $`\theta _\omega ^\kappa (x,s)=\psi (x)`$. Similarly, for the flow in 7, define
$$S_{t,s}^{\omega ,\epsilon }\psi (x)=\psi (\phi _{t,s}^{\omega ,\epsilon }(x)),s<t.$$
(10)
$`\theta _\omega ^\epsilon (x,t)=S_{t,s}^{\omega ,\epsilon }\psi (x)`$ solves the transport equation
$$\frac{\theta ^\epsilon }{t}+(u^\epsilon (x,t))\theta ^\epsilon =0,$$
(11)
with initial condition $`\theta (x,s)=\psi (x)`$. Similar definitions can be given for forward flows but we will restrict attention to the backward ones since we are primarily interested in scalar transport. The results given below generalize trivially to forward flows.
## Kraichnan Model
In Kraichnan introduced one of the simplest model of passive scalar by considering the advection by a Gaussian, spatially non-smooth and white-in-time velocity field. The fact that white-in-time velocity fields may exhibit intermittency was first recognized by Majda . Definitive work on Kraichnan model has been done afterwards in .
We will consider a generalization of Kraichnan model introduced in (see also ). The velocity $`u`$ is assumed to be a statistically homogeneous, isotropic and stationary Gaussian field with mean zero and covariance
$$𝐄u_\alpha (x,t)u_\beta (y,s)=(C_0\delta _{\alpha \beta }c_{\alpha \beta }(xy))\delta (ts).$$
(12)
We assume that $`u`$ has a correlation length $`\mathrm{}_0`$, i.e. the covariance in 12 decays fast for $`|xy|>\mathrm{}_0`$. Consequently $`c_{\alpha \beta }(x)C_0\delta _{\alpha \beta }`$ as $`|x|/\mathrm{}_0\mathrm{}`$. On the other hand, we will be mainly interested in small scale phenomena for which $`|x|\mathrm{}_0`$. In this range, we take $`c_{\alpha \beta }(x)=d_{\alpha \beta }(x)+O(|x|^2/\mathrm{}_0^2)`$ with
$$d_{\alpha \beta }(x)=Ad_{\alpha \beta }^P(x)+Bd_{\alpha \beta }^S(x),$$
(13)
and
$$\begin{array}{c}d_{\alpha \beta }^P(x)=D\left(\delta _{\alpha \beta }+\xi \frac{x_\alpha x_\beta }{|x|^2}\right)|x|^\xi ,\hfill \\ d_{\alpha \beta }^S(x)=D\left((d+\xi 1)\delta _{\alpha \beta }\xi \frac{x_\alpha x_\beta }{|x|^2}\right)|x|^\xi .\hfill \end{array}$$
(14)
$`D`$ is a parameter with dimension $`[\text{length}]^{2\xi }[\text{time}]^1`$. The dimensionless parameters $`A`$ and $`B`$ measure the divergence and rotation of the field $`u`$. $`A=0`$ corresponds to incompressible fields with $`u=0`$. $`B=0`$ corresponds to irrotational fields with $`\times u=0`$. The parameter $`\xi `$ measures the spatial regularity of $`u`$. For $`\xi (0,2)`$, the local characteristic of $`u`$ fails to be twice differentiable and this fact has important consequences on both the transport equation in 2 and the systems of ODEs in 3 or 5.
Existing physics literature concentrates on the $`\kappa `$-limit for Kraichnan model. Let $`𝒮^2=A+(d1)B`$, $`𝒞^2=A`$, $`𝒫=𝒞^2/𝒮^2`$. $`𝒫[0,1]`$ is a measure of the degree of compressibility of $`u`$. The pioneering work of Gawȩdzki and Vergassola (see also ) identifies two different regimes for the $`\kappa `$-limit:
1. The strongly compressible regime when $`𝒫d/\xi ^2`$. In this regime $`g_\omega ^\kappa `$ converges to a flow of maps, i.e. there exists a two-parameter family of maps $`\{\phi _{t,s}^\omega (x)\}`$ such that
$$g_\omega ^\kappa (x,t|dy,s)\delta (y\phi _{t,s}^\omega (x))dy.$$
(15)
Moreover particles have finite probability to coalesce under the flow of $`\{\phi _{t,s}^\omega (x)\}`$. In other words the flow is not invertible.
2. When $`𝒫<d/\xi ^2`$, $`g_\omega ^\kappa `$ converges to a “generalized stochastic flow”
$$g_\omega ^\kappa (x,t|dy,s)g_\omega (x,t|dy,s),$$
(16)
and the limit $`g_\omega `$ is a nontrivial probability distribution in $`y`$. This means that the image of a particle under the flow defined by the velocity field $`u`$ is non-unique and has a non-trivial distribution. In other words, particle trajectories branch.
The same classification of the flows was obtained by Le Jan and Raimond using Wiener chaos expansion without explicit reference to the $`\kappa `$ limit. In contrast, our primary motivation is to study the limit of physical regularizations.
The following result answers the question Q1 and also points out that there are three different regimes if both the $`\kappa `$\- and the $`\epsilon `$-limits are considered.
###### Theorem 1
In the strongly compressible regime when $`𝒫d/\xi ^2`$, there exists a two-parameter family of random maps $`\{\phi _{t,s}^\omega (x)\}`$, such that for all smooth test functions $`\psi `$ and for all $`(s,t,x)`$, $`s<t`$,
$$𝐄\left(S_{t,s}^{\omega ,\kappa }\psi (x)\psi (\phi _{t,s}^\omega (x))\right)^20,$$
(17)
as $`\kappa 0`$, and
$$𝐄\left(\psi (\phi _{t,s}^{\omega ,\epsilon }(x))\psi (\phi _{t,s}^\omega (x))\right)^20,$$
(18)
as $`\epsilon 0`$. Moreover, the limiting flow $`\{\phi _{t,s}^\omega (x)\}`$ coalesces in the sense that for almost all $`(t,x,y)`$, $`xy`$, we can define a time $`\tau `$ such that $`\mathrm{}<\tau <t`$ a.s. and
$$\phi _{t,s}^\omega (x)=\phi _{t,s}^\omega (y)\text{for }s\tau .$$
(19)
In the weakly compressible regime when $`𝒫(d+\xi 2)/2\xi `$, there exists a random family of generalized flows $`g_\omega (x,t|dy,s)`$, such that for all test function $`\psi `$,
$$S_{t,s}^\omega \psi (x)=_^d\psi (y)g_\omega (x,t|dy,s),$$
(20)
satisfies
$$𝐄\left(S_{t,s}^{\omega ,\kappa }\psi (x)S_{t,s}^\omega \psi (x)\right)^20,$$
(21)
as $`\kappa 0`$ for all $`(s,t,x)`$, $`s<t`$, and
$$𝐄\left(_^d\eta (x)\left(\psi (\phi _{t,s}^{\omega ,\epsilon }(x))S_{t,s}^\omega \psi (x)\right)𝑑x\right)^20,$$
(22)
as $`\epsilon 0`$ for all $`(s,t)`$, $`s<t`$, and for all test functions $`\eta `$. Moreover, $`g_\omega (x,t|dy,s)`$ is non-degenerate in the sense that
$$S_{t,s}^\omega \psi ^2(x)\left(S_{t,s}^\omega \psi (x)\right)^2>0\text{a.s.}$$
(23)
In the intermediate regime when $`(d+\xi 2)/2\xi <𝒫<d/\xi ^2`$, there exists a random family of generalized flows $`g_\omega (x,t|dy,s)`$, such that for all test function $`\psi `$ and for all $`(s,t,x)`$, $`s<t`$,
$$𝐄\left(S_{t,s}^{\omega ,\kappa }\psi (x)S_{s,t}^\omega \psi (x)\right)^20$$
(24)
as $`\kappa 0`$. In the $`\epsilon `$-limit, the flows $`\phi _{t,s}^{\omega ,\epsilon }(x)`$ converges in the sense of distributions, i.e. there exists a family of probability densities
$$\{G_n(x_1,\mathrm{},x_n,t|y_1,\mathrm{},y_n,s)dy_1\mathrm{}dy_n\},$$
(25)
$`n=1,2,\mathrm{}`$, such that
$$\begin{array}{c}\hfill 𝐄\psi (\phi _{t,s}^{\omega ,\epsilon }(x_1),\mathrm{},\phi _{t,s}^{\omega ,\epsilon }(x_n))\underset{^d\times \mathrm{}\times ^d}{}\psi (y_1,\mathrm{},y_n)\\ \hfill \times G_n(x_1,\mathrm{},x_n,t|y_1,\mathrm{},y_n,s)dy_1\mathrm{}dy_n,\end{array}$$
(26)
as $`\epsilon 0`$ for any continuous function $`\psi `$ with compact support. Furthermore, the $`\epsilon `$limit coalesces in the sense that
$$\begin{array}{c}\hfill G_2(x_1,x_2,t|y_1,y_2,s)=\stackrel{~}{G}_2(x_1,x_2,t|y_1,y_2,s)\\ \hfill +A(y_1,x_1,x_2,t,s)\delta (y_1y_2),\end{array}$$
(27)
with $`A>0`$ when $`t>s`$. Here $`\stackrel{~}{G}_2`$ is the absolutely continuous part of $`G_2`$ with respect to the Lebesgue measure. Similar statements hold for the other $`G_n`$’s. In particular, the $`\{G_n\}`$’s differ from the moments of the $`\kappa `$-limit $`g_\omega `$ defined in 24.
Rephrasing the content of this result, we have strong convergence to a family of flow maps in the strongly compressible regime for both the $`\kappa `$-limit and the $`\epsilon `$-limit. In the weakly compressible regime, we have strong convergence to a family of generalized flows for the $`\kappa `$-limit, but weak convergence to the same limit for the $`\epsilon `$-regularization. In fact, using the terminology of Young measures , the limiting generalized flow $`\{g_\omega (x,t|dy,s)\}`$ is nothing but the Young measure for the sequence of flow maps $`\{\phi _{s,t}^{\omega ,\epsilon }(x)\}`$. Finally, in contrast to what is observed in the other two regimes, the $`\epsilon `$-limit and $`\kappa `$-limit are not the same in the intermediate regime. As we will see below, the structure functions of the passive scalar field scale differently in the two limits.
From Theorem 1, it is natural to define the solution of the transport equation in 2 for the initial condition $`\theta _\omega (x,s)=\theta _0(x)`$ as
$$\theta _\omega (x,t)=S_{t,s}^\omega \theta _0(x)=_^d\theta _0(y)g_\omega (x,t|dy,s),$$
(28)
for the weakly compressible and the intermediate regimes in the $`\kappa `$-limit (non-degenerate cases), and as
$$\theta _\omega (x,t)=\theta _0(\phi _{t,s}^\omega (x)),$$
(29)
for the strongly compressible regime. In the intermediate regime in the $`\epsilon `$-limit, it makes sense to look at the limiting moments of $`\theta _\omega ^\epsilon (x,t)`$ since we have as $`\epsilon 0`$
$$\begin{array}{c}\hfill 𝐄(\theta _\omega ^\epsilon (x_1,t)\mathrm{}\theta _\omega ^\epsilon (x_n,t))\underset{^d\times \mathrm{}\times ^d}{}\theta _0(y_1)\mathrm{}\theta _0(y_n)\\ \hfill \times G_n(x_1,\mathrm{},x_n,t|y_1,\mathrm{},y_n,s)dy_1\mathrm{}dy_n.\end{array}$$
(30)
It should be noted that when $`g_\omega `$ is non-degenerate, there exists an anomalous dissipation mechanism for the scalar, whereas no such anomalous dissipation is present in the coalescence cases . The presence of anomalous dissipation is the primary reason why the transport equation in 2 has a statistical steady state (invariant measure) if it is appropriately forced, as we will show later.
Details of the proof of Theorem 1 are given in . Crucial to the proof is the study of $`P(\rho |r,s)`$ defined through $`\epsilon `$-regularization as
$$_0^{\mathrm{}}\eta (r)P(\rho |r,st)𝑑r=\underset{\epsilon 0}{lim}𝐄\eta (|\phi _{t,s}^{\omega ,\epsilon }(y)\phi _{t,s}^{\omega ,\epsilon }(z)|),$$
(31)
where $`\eta `$ is a test function, and similarly through $`\kappa `$-regularization. Here $`\rho =|yz|`$ and $`s<t`$. $`P(\rho |r,s)`$ can be thought of as the probability density that two particles have distance $`r`$ at time $`s<t`$ if their final distance is $`\rho `$ at time $`t`$. For Kraichnan model, $`P`$ satisfies the backward equation
$$\frac{P}{s}=\frac{}{r}\left(b(r)P\right)+\frac{^2}{r^2}\left(a(r)P\right),$$
(32)
for the final condition $`lim_{s0}P(\rho |r,s)=\delta (r\rho )`$, and with $`a(r)`$, $`b(r)`$ such that
$$\begin{array}{c}a(r)=D(𝒮^2+\xi 𝒞^2)r^\xi +O(r^2/\mathrm{}_0^2),\\ b(r)=D((d1+\xi )𝒮^2\xi 𝒞^2)r^{\xi 1}+O(r/\mathrm{}_0^2).\end{array}$$
(33)
For $`r\mathrm{}_0`$, $`a(r)`$ tends to $`C_0`$, $`b(r)`$ to $`C_0(d1)/r`$, and the equation in 32 reduces to a diffusion equation with constant coefficient. The equation in 32 is singular at $`r=0`$. The proof of Theorem 1 is essentially reduced to the study of this singular diffusion equation. This is also the main step for which the white-in-time nature of the velocity field is crucial.
## Structure Functions
We now study some consequences of Theorem 1 for the passive scalar $`\theta _\omega `$ defined in 28 or 29. We note that the scaling of the second-order structure function is the same for the $`\kappa `$\- and the $`\epsilon `$-limits in the strongly and the weakly compressible cases , but it differs in the intermediate regime as a result of the difference between the limits in 24 and 26. For simplicity of presentation, we assume that $`\theta _0`$ is isotropic and Gaussian. Denote ($`n`$)
$$S_{2n}(|xy|,t)=𝐄(\theta _\omega (x,t)\theta _\omega (y,t))^{2n},$$
(34)
$$\text{or}S_{2n}(|xy|,t)=\underset{\epsilon 0}{lim}𝐄(\theta _\omega ^\epsilon (x,t)\theta _\omega ^\epsilon (y,t))^{2n},$$
$`(\mathrm{𝟑𝟒}^{})`$
in the intermediate regime in the $`\epsilon `$-limit. In the strongly compressible case, we have for both the $`\kappa `$\- and the $`\epsilon `$-limits
$$S_2(r,t)=O(r^\zeta ),$$
(35)
with
$$\zeta =\frac{2d\xi +2\xi 𝒫}{1+\xi 𝒫}.$$
(36)
In the weakly compressible case, we have for both the $`\kappa `$\- and the $`\epsilon `$-limits
$$S_2(r,t)=O(r^{2\xi }).$$
(37)
In the intermediate regime, the limits differ, and the $`\kappa `$-limit scales as in 37, whereas the $`\epsilon `$-limit scales as in 35. The equations in 35 and 37 can be derived upon expressing $`s_2`$ in terms of $`P`$; the details are given in Ref. .
It is interesting to discuss the higher order structure functions both in the non-degenerate and in the coalescence cases in 35 and 37 since their scalings highlight very different behavior of the scalar. We consider first the coalescence cases which are simpler. In these cases, because of the absence of dissipative anomaly, all higher order structure functions can again be expressed in terms of $`P`$, and it can be shown that
$$S_{2n}(r,t)=O(r^\zeta ),$$
(38)
with $`\zeta `$ given by 36 for all $`n2`$. In fact, coalescence implies that the temperature field $`\theta _\omega `$ tends to become flat except possibly on a zero-measure set where it presents shock-like discontinuities. Such a situation with two kinds of spatial structures for $`\theta _\omega `$ is usually refered to as bi-fractal, and, in simple cases, one may identify $`\zeta `$ with the codimension of the set supporting the discontinuities of $`\theta _\omega `$ .
The non-degenerate cases are more complicated. In these cases, one expects that $`\theta _\omega `$ presents a spatial behavior much richer than in the coalescence cases, with all kinds of scalings present. This is the multi-fractal situation for which the higher order structure functions behave as
$$S_{2n}(r,t)=O(r^{\zeta _{2n}}),$$
(39)
with $`\zeta _{2n}<n(2\xi )`$ for $`2n>2`$. The actual value of the $`\zeta _n`$’s cannot be obtained by dimensional analysis, and one has to resort to various sophisticated perturbation techniques (see ). We will consider again the scaling of the structure functions at statistical steady state in the section on incomplete self-similarity.
## One Force, One Solution Principle for Temperature
We now turn to question Q3 and consider the existence of a statistical steady state for the transport equation with appropriate forcing. We restrict attention to the non-degenerate cases which include the weakly compressible regime and the intermediate regime in the $`\kappa `$-limit. Indeed, in these regimes the non-degeneracy of $`g_\omega (x,t|dy,s)`$ as a probability distribution in $`y`$ implies dissipation of energy or, phrased differently, decay in memory in the semi-group $`S_{t,s}`$ generated by $`\{g_\omega \}`$. We show that the anomalous dissipation is strong enough in order that the forced transport equation has a unique invariant measure for both the weakly compressible regime and the intermediate regime in the $`\kappa `$-limit. This result, however, depends on the finiteness of $`\mathrm{}_0`$. In limit as $`\mathrm{}_0\mathrm{}`$ an invariant measure exists only for the weakly compressible regime.
We will consider (compare with 2)
$$\frac{\theta }{t}+(u(x,t))\theta =b(x,t).$$
(40)
where $`b`$ is a white-noise forcing such that
$$𝐄b(x,t)b(y,s)=B(|xy|)\delta (ts).$$
(41)
$`B(r)`$ is assumed to be smooth and rapidly decaying to zero for $`rL`$; $`L`$ will be referred to as the forcing scale. The solution of 40 for the initial condition $`\theta _\omega (x,s)=\theta _0(x)`$ is understood as
$$\theta _\omega (x,t)=S_{t,s}^\omega \theta _0(x)+_s^tS_{t,\tau }^\omega b(x,\tau )𝑑\tau .$$
(42)
Define the product probability space $`(\mathrm{\Omega }_u\times \mathrm{\Omega }_b,_u\times _b,𝒫_u\times 𝒫_b)`$, and the shift operator $`T_\tau \omega (t)=\omega (t+\tau )`$, with $`\omega =(\omega _u,\omega _b)`$. We have
###### Theorem 2 (One force–one solution I)
For $`d>2`$, in the weakly compressible regime and in the intermediate regime in the $`\kappa `$-limit, for almost all $`\omega `$, there exists a unique solution of 40 defined on $`^d\times (\mathrm{},\mathrm{})`$. This solution can be expressed as
$$\theta _\omega ^{}(x,t)=_{\mathrm{}}^tS_{t,s}^\omega b(x,s)𝑑s.$$
(43)
Furthermore the map $`\omega \theta _\omega ^{}`$ satisfies the invariance property
$$\theta _{T_\tau \omega }^{}(x,t)=\theta _\omega ^{}(x,t+\tau ).$$
(44)
Theorem 2 is the “one force, one solution” principle articulated in . Because of the invariance property 44, the map in 43 leads to a natural invariant measure. As a consequence we have
###### Corollary 3
For $`d>2`$, in the weakly compressible regime and in the intermediate regime in the $`\kappa `$-limit, there exists a unique invariant measure on $`L_{\text{loc}}^2(^d\times \mathrm{\Omega })`$ for the dynamics defined by 40.
The connection between the map 43 and the invariant measure, together with uniqueness, is explained in . The restriction on the dimensionality in Theorem 2 arises because the velocity field has finite correlation length $`\mathrm{}_0`$: Theorem 2 is changed into Theorem 4 below in the limit as $`\mathrm{}_0\mathrm{}`$ which can be considered after appropriate redefinition of the velocity field.
We sketch the proof of Theorem 2. Basically, it amounts to verifying that the dissipation in the system is strong enough in the sense that
$$𝐄\left(_{T_1}^{T_2}_^dS_{t,t+s}^\omega b(x,s)𝑑s\right)^20,$$
(45)
as $`T_1`$, $`T_2\mathrm{}`$ for fixed $`x`$ and $`t`$. The average in 45 is given explicitly by
$$_{T_1}^{T_2}_0^{\mathrm{}}B(\rho )P(0|\rho ,s)𝑑r𝑑s,$$
(46)
where $`P`$ satisfies 32. The convergence of the integral in 45 depends on the rate of decay in $`|s|`$ of $`P(0|\rho ,s)`$. The latter can be estimated by studying the equation in 32 , which yields $`P(0|\rho ,s)C\rho ^\alpha |s|^{d/2}`$ with $`\alpha =(d1\xi (\xi +1)𝒫)/(1+\xi 𝒫)`$ for $`|s|`$ large and $`\rho \mathrm{}_0`$. Hence, the integral in $`s`$ in 46 tends to zero as $`T_1`$, $`T_2\mathrm{}`$ if $`d>2`$. It follows that the invariant measure in 43 exists provided $`d>2`$.
We now ask what happens if we let $`\mathrm{}_0\mathrm{}`$ in order to emphasizes the effect of the inertial range of the velocity? This question, however, has to be considered carefully because the velocity field with the covariance in 12 diverges as $`\mathrm{}_0\mathrm{}`$. The right way to proceed is to consider an alternative velocity $`v`$, taken to be Gaussian, white-in-time, but non-homogeneous, with covariance
$$\begin{array}{c}𝐄v_\alpha (x,t)v_\beta (y,s)\hfill \\ =(c_{\alpha \beta }(xa)+c_{\alpha \beta }(ay)c_{\alpha \beta }(xy))\delta (ts).\hfill \end{array}$$
(47)
For finite $`\mathrm{}_0`$, one has $`v(x,t)=u(x,t)u(a,t)`$, where $`a`$ is arbitrary but fixed. However, $`v`$ makes sense in the limit as $`\mathrm{}_0\mathrm{}`$. Denote by $`\vartheta _\omega (x,t)`$ the temperature field advected by $`v`$, i.e. the solution of the transport equation 40 with $`u`$ replaced by $`v`$:
$$\frac{\vartheta }{t}+(v(x,t))\vartheta =b(x,t).$$
(48)
Restricting to zero initial condition, it follows from the homogeneity of the forcing that the single-time moments of $`\theta _\omega `$ and $`\vartheta _\omega `$ coincide for finite $`\mathrm{}_0`$, but in contrast to $`\theta _\omega `$, $`\vartheta _\omega `$ makes sense as $`\mathrm{}_0\mathrm{}`$. Thus, $`\vartheta _\omega `$ is a natural process to study the limit as $`\mathrm{}_0\mathrm{}`$, and we have
###### Theorem 4 (One force–one solution II)
In the limit as $`\mathrm{}_0\mathrm{}`$ in the weakly compressible regime, for almost all $`\omega `$, there exists a unique solution of 48 defined on $`^d\times (\mathrm{},\mathrm{})`$. This solution can be expressed as
$$\vartheta _\omega ^{}(x,t)=_{\mathrm{}}^tS_{t,s}^\omega b(x,s)𝑑s,$$
(49)
where $`S_{s,t}^\omega `$ is the semi-group for the generalized flow associated with the velocity defined in 47 in the limit as $`\mathrm{}_0\mathrm{}`$. Furthermore the map $`\omega \vartheta _\omega ^{}`$ satisfies the invariance property
$$\vartheta _{T_\tau \omega }^{}(x,t)=\vartheta _\omega ^{}(x,t+\tau ).$$
(50)
As a direct result we also have
###### Corollary 5
In the limit as $`\mathrm{}_0\mathrm{}`$, in the weakly compressible regime there exists a unique invariant measure on $`L_{\text{loc}}^2(^d\times \mathrm{\Omega })`$ for the dynamics defined by 48.
Notice that, as $`\mathrm{}_0\mathrm{}`$, the anomalous dissipation is not strong enough in the intermediate regime in the $`\kappa `$-limit, for which no statistical steady state with finite energy exists.
The proof of Theorem 4 proceeds as the one for Theorem 2, but the estimate for $`P`$ in 46 changes as $`P(0|\rho ,s)C\rho ^\alpha |s|^{(\alpha +1)/(2\xi )}`$ with $`\alpha =(d1\xi (\xi +1)𝒫)/(1+\xi 𝒫)`$ for $`|s|`$ large and $`\rho \mathrm{}_0`$. It follows that the integral in $`s`$ in 46 converges as $`T_1`$, $`T_2\mathrm{}`$ in the weakly compressible regime only.
## One Force, One Solution Principle for the Temperature Difference
Since no anomalous dissipation is present in the coalescence cases, i.e the strongly compressible regime and the intermediate regime in the $`\epsilon `$-limit, no invariant measure for the temperature field exists in these regimes. It makes sense, however, to ask about the existence of an invariant measure for the temperature difference, i.e. to consider
$$\delta \theta _\omega (x,y,t)=_T^tS_{t,s}^\omega (b(x,s)b(y,s))𝑑s,$$
(51)
in the limit as $`T\mathrm{}`$. When $`\theta _\omega ^{}`$ exists, one has
$$\delta \theta _\omega ^{}(x,y,t)=\underset{T\mathrm{}}{lim}\delta \theta _\omega (x,y,t)=\theta _\omega ^{}(x,t)\theta _\omega ^{}(y,t),$$
(52)
but it is conceivable that $`\delta \theta _\omega ^{}`$ exists in the coalescence cases even though $`\theta _\omega ^{}`$ is not defined. The reason is that coalescence of the generalized flow implies that the temperature field flattens with time, which is a dissipation mechanism as far as the temperature difference is concerned. Of course, this effect has to overcome the fluctuations produced by the forcing, and the existence of an invariant measure such as 51 will depend on how fast particles coalesce under the flow, which happens only in the limit as $`\mathrm{}_0\mathrm{}`$ (i.e. for the alternate velocity defined in 47) as we show now.
For finite $`\mathrm{}_0`$, if we were to consider two particles separated by much more than the correlation length $`\mathrm{}_0`$, the dynamics of their distance under the flow is governed by the equation in 32 for $`r\mathrm{}_0`$, i.e. by a diffusion equation with constant diffusion coefficient on the scale of interest. It follows that no tendency of coalescence is observed before the distance becomes smaller than $`\mathrm{}_0`$, which, as shown below, does not happen fast enough in order to overcome the the fluctuations produced by the forcing. In other words,
###### Lemma 6
In the coalescence cases, for finite $`\mathrm{}_0`$, there is no invariant measure with finite energy for the temperature difference.
Consider now the limit as $`\mathrm{}_0\mathrm{}`$, and let $`\delta \vartheta _\omega (x,y,t)=\vartheta _\omega (x,t)\vartheta _\omega (y,t)`$ where $`\vartheta _\omega `$ solves the equation in 48. The temperature difference $`\delta \vartheta _\omega `$ satisfies the transport equation
$$\frac{\delta \vartheta }{t}+(v(x,t)_x+v(y,t)_y)\delta \vartheta =b(x,t)b(y,t).$$
(53)
We have
###### Theorem 7 (One force–one solution III)
In the limit as $`\mathrm{}_0\mathrm{}`$, for almost all $`\omega `$, in the strongly and the weakly compressible regimes, as well as in the intermediate regime if the flow is non-degenerate, there exists a unique solution of 53 defined on $`^d\times (\mathrm{},\mathrm{})`$. This solution can be expressed as
$$\delta \vartheta _\omega ^{}(x,y,t)=_{\mathrm{}}^tS_{t,s}^\omega (b(x,s)b(y,s))𝑑s,$$
(54)
where $`S_{s,t}`$ is the semi-group for the generalized flow associated with the velocity defined in 47 in the limit as $`\mathrm{}_0\mathrm{}`$. Furthermore the map $`\omega \delta \vartheta _\omega ^{}`$ satisfies the invariance property
$$\delta \vartheta _{T_\tau \omega }^{}(x,y,t)=\delta \vartheta _\omega ^{}(x,y,t+\tau ).$$
(55)
An immediate consequence of this theorem is
###### Corollary 8
In the limit as $`\mathrm{}_0\mathrm{}`$, in the strongly and the weakly compressible regimes, as well as in the intermediate regime if the flow is non-degenerate, there exists a unique invariant measure on $`L_{\text{loc}}^2(^d\times \mathrm{\Omega })`$ for the dynamics defined by 53.
The proof of Theorem 7 proceeds similarly as the proof of Theorem 2. In the non-degenerate cases, one study the convergence of (compare 45)
$$𝐄\left(_{T_1}^{T_2}_^dS_{t,t+s}^\omega (b(x,s)b(y,s))𝑑s\right)^20,$$
(56)
as $`T_1`$, $`T_2\mathrm{}`$ for fixed $`x`$ and $`t`$. The average in 56 can be expressed in terms of $`P`$, and it can be shown that the expression in (56) converges as $`T_1`$, $`T_2\mathrm{}`$ in the non-degenerate cases. In the strongly compressible regimes, because of the existence of a flow of maps, 56 is replaced by
$$𝐄\left(_{T_1}^{T_2}(b(\phi _{t,s}^\omega (x),s)b(\phi _{t,s}^\omega (y),s))𝑑s\right)^2.$$
(57)
This average can again be expressed in terms of $`P`$, and it can be shown that the convergence of the time integral in 57 depends on the rate at $`P`$ looses mass at $`r=0+`$ (i.e. the rate at which particles coalesce). The analysis of the equation in 32 shows that the process is fast enough in order that the integral over $`s`$ in 57 tends to zero as $`T_1`$, $`T_2\mathrm{}`$ in the strongly compressible regime. In contrast, the equivalent of 57 in the intermediate regime in the $`\epsilon `$-limit can be shown to diverge as $`T_1`$, $`T_2\mathrm{}`$.
It can be shown that the invariant measure has finite correlation functions of all order, even though these results do not by themselves imply uniqueness of stationary solutions to the $`n`$-point Fokker-Planck equation. The task of studying the passive scalar is now changed to the study of the short distance behavior of these correlation functions.
## Incomplete self-similarity
We finally turn to question Q4 and consider the scaling of the structure functions based on the invariant measure $`\delta \vartheta ^{}`$ defined in 54. Denote
$$𝒮_n(|xy|)=𝐄|\delta \vartheta _\omega ^{}(x,y,t)|^n.$$
(58)
The dimensional parameters are $`B_0=B(0)`$ ($`[`$temperature$`]^2`$$`[`$time$`]^1`$),
$`D`$ ($`[`$length$`]^{2\xi }`$$`[`$time$`]^1`$), $`L`$ ($`[`$length$`]`$). It follows that
$$𝒮_n(r)=\left(\frac{B_0r^{2\xi }}{D}\right)^{n/2}f_n\left(\frac{r}{L}\right),$$
(59)
where the $`f_n`$’s are dimensionless functions which cannot be obtained by dimensional arguments. For instance, the scalings in 38, 39 correspond to different $`f_n`$. It is however obvious from the equation 59 that, provided the limit exists and is non-zero
$$\underset{L\mathrm{}}{lim}𝒮_n(r)=C_n\left(\frac{B_0r^{2\xi }}{D}\right)^{n/2}=O(r^{n(2\xi )/2}).$$
(60)
where $`C_n=lim_r\mathrm{}f_n(r/L)`$ are numerical constants. The scaling in 60 is usually referred to as the normal scaling since, consistent with Kolmogorov’s picture, it is independent of the forcing or the dissipation scales. In contrast, anomalous scaling is a statement that the structure functions diverge in the limit of infinite forcing scale, $`L\mathrm{}`$. In the spirit of Barenblatt-Chorin , we may say that normal scaling holds in case of complete self-similarity, whereas anomalous scaling is equivalent to incomplete self-similarity.
It is interesting to discuss the existence or non-existence of the limit in 60 for both the coalescence and the non-degenerate cases. When the flow coalesces, because of the existence of a flow of maps and the absence of dissipative anomaly, the $`𝒮_{2n}`$’s of even order $`2n2`$ can be computed exactly . It gives $`𝒮_{2n}(r)=\mathrm{}`$ for $`n\zeta /(2\xi )`$, whereas
$$𝒮_{2n}(r)=O(r^\zeta )\text{for }n<\frac{\zeta }{2\xi },$$
(61)
where $`\zeta `$ is given in 36. Thus, for $`n<\zeta /(2\xi )`$,
$$f_{2n}(r)=O\left((r/L)^{\zeta n(2\xi )}\right).$$
(62)
It follows that $`f_{2n}`$ and, hence, $`𝒮_{2n}`$ tend to zero as $`L\mathrm{}`$ for $`2n<\zeta /(2\xi )`$, whereas they are infinite for all $`L`$ for $`n\zeta /(2\xi )`$. In fact, in the coalescence case, it can be shown that on scales much larger than the forcing scale $`L`$, the structure functions of order $`n<\zeta /(2\xi )`$ behave as
$$𝒮_{2n}(r)C_{2n}r^{n(2\xi )}\text{as }r/L\mathrm{}.$$
(63)
Thus in the coalescence case, it is more natural to consider the limit as $`L0`$ of the structure functions, for which the expression in 63 shows the absence of intermittency corrections.
In the non-degenerate case, one has
$$𝒮_2(r)=O(r^{2\xi }),$$
(64)
while perturbation analysis gives for the higher order structure functions
$$𝒮_{2n}(r)=O(r^{\zeta _{2n}}),$$
(65)
with $`\zeta _{2n}<n(2\xi )`$ for $`2n>2`$. It follows that $`f_2(r)=O(1)`$, while
$$f_{2n}(r)=O\left((r/L)^{\zeta _nn(2\xi )}\right),2n>2.$$
(66)
In other words, as $`L\mathrm{}`$, $`𝒮_2`$ has a limit which exhibits normal scaling, whereas the $`𝒮_{2n}`$’s, $`2n>2`$, diverge. This may be closely related to the argument in that, in appropriate limits, intermittency corrections may disappear and higher than fourth order structure functions may not exist. We note, however, that Barenblatt and Chorin were discussing the case of infinite Reynolds number (here infinite Peclet number, $`\kappa 0`$) at finite $`L`$, whereas we require $`L\mathrm{}`$.
We thank many people for helpful discussions, including G. Barenblatt, E. Balkovsky, M. Chertkov, A. Chorin, U. Frisch, J. Goodman, A. Majda, and S. R. S. Varadhan. We are particularly grateful to K. Gawȩdzki for pointing out an error in the first draft of this paper. W. E is partially supported by a Presidential Faculty Fellowship from NSF. E. Vanden Eijnden is partially supported by NSF Grant DMS-9510356.
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# 1 Introduction
## 1 Introduction
It seems quite natural to expect that the harder a high-energy collision is, the higher is the number of fragments. One of the most tractable and widely explored processes where this phenomenon can be seen is deeply inelastic scattering (DIS) with a possibility to change a hard scale ($`Q^2`$, gauge boson virtuality) and to detect its influence (if any) on the hadronic invariant mass ($`W`$) ”efficiency”. The simplest measure of this efficiency is the multiplicity of secondaries.
In 1969 Yang and his collaborators , basing themselves on the “fragmentation picture“ of violent collisions, made a qualitative prediction: “…for larger values of the momentum transfer $`t`$, the breakup process favors larger multiplicities of hadrons“ (at fixed hadronic mass). Early searches for this effect were inconclusive in both theory and experiment .
In the framework of QCD a quantitative result has been obtained in Refs. : it appears that QCD gluon bremsstrahlung leads to an increase of the hadron multiplicities in DIS, with an increase of $`Q^2`$ at fixed hadronic mass $`W`$, but that this increase is very slow. The distinctive feature of this result is that $`n^{DIS}(W,Q^2)`$ has a finite limit at $`Q^2\mathrm{}`$ and $`W`$ fixed.
Later, another result has been claimed in Ref. , which predicted an infinite and quite rapid growth of $`n^{DIS}(W,Q^2)`$ with $`Q^2`$. In the course of the inference of this result it was supposed that the influence of the (non-perturbative) composite structure of the nucleon is negligible while in it plays a key role in the slowness of the $`Q^2`$ dependence of the multiplicity at fixed $`W`$. There is even more trivial objection. On general grounds, the infinite growth with $`Q^2`$ at fixed $`W`$ is impossible because of the apriori kinematical bound
$$n^{DIS}(W,Q^2)\frac{W}{m},$$
(1)
where $`m`$ is some effective mass.
Quite recently, a weak dependence on $`Q^2`$ in the framework of the Dual Parton Model was mentioned in .
Experimentally a statistically significant effect of the slow growth of $`n^{DIS}(W,Q^2)`$ was established in $`\nu (\overline{\nu })p`$ interactions and in $`\mu ^+p`$ interactions . The results of the EMC have been described in the framework of QCD in Ref. .
However, subsequent measurements at HERA (H1) were interpreted as a practical $`Q^2`$ independence of $`n^{DIS}(W,Q^2)`$ (for the current hemisphere in the hadronic c.m.s.), while H1 and ZEUS reported quite fast $`Q^2`$ dependence for the current hemisphere in the Breit frame. It should be noted, however, that this last result concerns different bins in $`W`$ for changing $`Q^2`$ values. Anyway, the situation is controversial and therefore very interesting.
In this paper we give our own interpretation of the HERA data on charged hadron multiplicities in the current fragmentation region; as will be seen in the text below, these are in agreement with Yang’s general hypothesis and our early QCD results (see also the review ).
## 2 Hadronic Spectrum and Multiplicity in DIS
According to the factorization for inclusive spectra in DIS, the hadronic spectrum in DIS is represented by two terms:
$$\frac{dn}{dy}^{DIS}(W,Q^2)=_{x_0}^1\frac{dz}{z}w(x,z,Q^2)\frac{d\widehat{n}}{dy}(W_{eff},Q^2)+\frac{dn_0}{dy},$$
(2)
where $`x_0=x+(1x)(m_h/W)\mathrm{exp}(y)`$, $`y`$ is the rapidity of the detected hadron and $`m_h`$ is its mass. In Eq. (2) $`d\widehat{n}/dy`$ defines the hadronic spectrum in partonic subprocess, while the quantity $`dn_0/dy`$ describes the spectrum of the proton remnant. The latter does not contribute to the current fragmentation region of DIS at HERA energies.
Correspondingly, the average hadronic multiplicity in DIS is represented by
$$n^{DIS}(W,Q^2)=_{x_0}^1\frac{dz}{z}w(x,z,Q^2)\widehat{n}(W_{eff},Q^2)+n_0.$$
(3)
For small $`x`$ the weight $`w(x,z,Q^2)`$ in (2), (3) is of the form:
$`w(x,z,Q^2)`$ $`=`$ $`D_g^q({\displaystyle \frac{x}{z}},Q^2,Q_0^2)f_g(z,Q_0^2)`$ (4)
$`\times `$ $`\left({\displaystyle _{x_0}^1}{\displaystyle \frac{dz}{z}}D_g^q({\displaystyle \frac{x}{z}},Q^2,Q_0^2)f_g(z,Q_0^2)\right)^1.`$
As can be seen, the hadronic spectrum in partonic subprocess, $`d\widehat{n}/dy`$, and the hadronic multiplicity $`\widehat{n}`$ depend on the effective energy, which is smaller than $`W`$:
$$W_{eff}^2=\frac{zx}{1x}W^2.$$
(5)
In what follows, we shall work in the c.m.s. of the final hadrons. In terms of rapidity, the current region in the c.m.s. corresponds to
$$Y<y<0$$
(6)
(it is assumed that the proton goes in the positive direction).
In our papers it has been established that the total hadronic multiplicity in the partonic subprocess of DIS is related to the hadronic multiplicity in $`e^+e^{}`$ annihilation:
$$\widehat{n}(W,Q^2)n^{e^+e^{}}(W)$$
(7)
(up to small NLO corrections, which decrease in $`Q^2`$).
In the partonic subprocess, the rapidity varies in the range
$$\widehat{Y}<y+y_0<\widehat{Y},$$
(8)
where $`\widehat{Y}=\mathrm{ln}(W_{eff}/m_h)`$ and
$$y_0=\frac{1}{2}\mathrm{ln}\left(\frac{1x}{1z}\right).$$
(9)
The quantity $`y_0`$ determines the rapidity of the centre of mass of the partonic subprocess in the centre of mass of the complete process. On integration over $`z`$, the region (8) is ”smeared” into the region
$$Y<y<Y,$$
(10)
with $`Y=\mathrm{ln}(W/m_h)`$.
The average value of the effective energy in (2), (3), available for particle production, appeared to be dependent on both $`W`$ and $`Q^2`$ :
$$W_{eff}^2\kappa (Q^2)W^2.$$
(11)
The efficiency factor $`\kappa (Q^2)`$, which stands in front of $`W^2`$ in (11), is much less than 1 and grows slowly in $`Q^2`$.
From formulas (3), (7) and (11), one can see that the rise of the average hadronic multiplicity in DIS has the same physical nature as in $`e^+e^{}`$ annihilation. For the first time this behaviour has been experimentally established by H1 in 1996 .
However, the QCD growth of $`\widehat{n}`$ is delayed in DIS by the bound-state effects and the slow QCD evolution of the structure function. This is why we predicted that the $`Q^2`$ dependence of $`\widehat{n}`$ at fixed $`W`$ should remain numerically weak at HERA energies .
It follows from (8), (9) that the centre of the spectrum is shifted to the region of positive rapidities and tends to zero at asymptotically high $`Q^2`$ :
$$y_0|{}_{Q^2\mathrm{}}{}^{}\frac{1}{\mathrm{ln}(\mathrm{ln}Q^2))}.$$
(12)
The hadronic spectrum in partonic subprocess in the c.m.s. of DIS has the form
$$\frac{d\widehat{n}^h}{dy}=n^{e^+e^{}}(W_{eff},Q^2)\overline{D}^h(W_{eff},y).$$
(13)
Normalization in the RHS of Eq. (13) is done in agreement with formula (3).
## 3 Hadronic Multiplicities in the Current <br>Fragmentation Region
To calculate the multiplicity of charged hadrons in the current fragmentation region, we have to define expressions of the quark distribution at small $`x`$, of the hadronic spectrum in the partonic subprocess as well as of the multiplicity of charged hadrons in $`e^+e^{}`$ annihilation.
For the quark distribution, we use an analytical expression from Ref. , in the case of soft initial conditions. As was shown in , at small $`x`$ it is in good agreement with the data on the structure function from HERA in the wide ranges of $`Q^2`$. Namely, at high $`Q^2`$, we have
$$D_g^q(z,Q^2)rI_1(t)\mathrm{exp}(d\xi /2),$$
(14)
with $`d=\beta _0+20N_f/27`$. The variable
$$t=2\sqrt{6\xi \mathrm{ln}\left(\frac{1}{z}\right)}$$
(15)
is related to the QCD evolution parameter
$$\xi =\frac{2}{\beta _0}\mathrm{ln}\left(\frac{\alpha (Q_0^2)}{\alpha (Q^2)}\right),$$
(16)
where $`\beta _0=112N_f/3`$ is the $`\beta `$-function in lowest order and
$$r=\frac{t}{2\mathrm{ln}(1/z)}.$$
(17)
The quark and gluon distributions from Ref. obey the GLAP evolution equations .
The expression of the initial gluon distribution at $`z`$ closed to $`1`$ is chosen to have the following form
$$f_g(z,Q_0^2)|{}_{z1}{}^{}(1z)^{n_g}.$$
(18)
We have omitted constant factors in the RHS of Eqs. (14) and (18) as they do not influence our final results.
The spectrum of hadrons in the partonic process $`\overline{D}^h`$ was calculated by many authors. We use the expression from Refs. ($`N`$ is a normalization factor):
$$\overline{D}^h(W,\zeta )=\frac{N}{\sigma \sqrt{2\pi }}\mathrm{exp}\left[\frac{1}{8}k\frac{1}{2}s\delta \frac{1}{4}(2+k)\delta ^2+\frac{1}{6}s\delta ^3+\frac{1}{24}k\delta ^4\right]$$
(19)
calculated in the variable
$$\zeta =\mathrm{ln}\left(\frac{W}{E_h}\right).$$
(20)
Here $`E_h`$ is the energy of the detected hadron.
The average value of $`\zeta `$, $`\zeta _0`$, and its dispersion $`\sigma `$ are given by the formulas:
$$\zeta _0=\frac{1}{2}\tau \left(1+\frac{\rho }{24}\sqrt{\frac{48}{\beta _0\tau }}\right)\left(1\frac{\omega }{6\tau }\right),$$
(21)
$$\sigma =\sqrt{\frac{\tau }{3}}\left(\frac{\beta _0\tau }{48}\right)^{1/4}\left(1\frac{\beta _0}{64}\sqrt{\frac{48}{\beta _0\tau }}\right)\left(1+\frac{\omega }{8\tau }\right),$$
(22)
where
$$\tau =\mathrm{ln}\left(\frac{W}{\mathrm{\Lambda }}\right)$$
(23)
and
$$s=\frac{\rho }{16}\sqrt{\frac{3}{\tau }}\left(\frac{48}{\beta _0\tau }\right)^{1/4}\left(1+\frac{\omega }{4\tau }\right),$$
(24)
$$k=\frac{27}{5\tau }\left(\sqrt{\frac{\beta _0\tau }{48}}\frac{\beta _0}{24}\right)\left(1+\frac{5\omega }{12\tau }\right),$$
(25)
$$\delta =\frac{\zeta \zeta _0}{\sigma }.$$
(26)
Here $`\rho =11+2N_f/27`$, $`\omega =1+N_f/27`$.
At low (effective) energies we use the fit of the low-energy data on multiplicity of charged hadrons in $`e^+e^{}`$ annihilation from Ref. :
$$n^{e^+e^{}}=2.67+0.48\mathrm{ln}W^2,$$
(27)
while for high energies ($`W_{eff}>10`$ GeV) we apply the fit from Ref. , which well decsribes $`e^+e^{}`$ data up to LEP energies:
$$n^{e^+e^{}}=1.66+0.866\mathrm{exp}(1.047\sqrt{\mathrm{ln}W^2}).$$
(28)
We have corrected (27) for a fraction of the charged particles from $`K_s^0`$ and $`\mathrm{\Lambda }(\overline{\mathrm{\Lambda }})`$ decays.
Figure 1 represents the result of calculations of charged multiplicity in current hemisphere of the c.m.s. by using formulas (2) and (14) (solid curves) in comparison with the H1 data from Ref. . As can be seen, our QCD predictions are in very good agreement with the data. These are quite compatible with a slow growth of $`n^{DIS}(W,Q^2)`$ in $`Q^2`$ at fixed $`W`$.
Figure 2 demonstrates a rapid rise of $`n^{DIS}(W,Q^2)`$ in the variable $`W`$ for different values of $`Q^2`$, which was predicted many years ago in Refs. and seen previously in $`e^+e^{}`$ annihilation (the very values of $`Q^2`$ taken from ). Let us note that the H1 data presented in Fig. 2 (see Table 4 in ) do not correspond to some fixed values of $`Q^2`$, in contrast with the experimental points in Fig. 1.
In order to obtain multiplicity of charged hadrons in current region of the Breit frame, the c.m.s. spectrum (13) must be integrated in the region
$$Y<y<y_B,$$
(29)
where
$$y_B=\frac{1}{2}\mathrm{ln}\left(\frac{1+v}{1v}\right)\frac{1}{2}\mathrm{ln}\left(\frac{1}{x}\right).$$
(30)
The quantity
$$v=\sqrt{14x(1x)}$$
(31)
in the RHS of Eq. (30) is the velocity of the Breit frame in the c.m.s. So, $`y_B`$ corresponds to zero rapidity in this frame.
The results of our calculations of the multiplicity of charged hadrons in the current region of the Breit frame are presented in Fig. 3 as a function of $`Q^2`$, in comparison with the H1 data (solid squares) and ZEUS data (solid circles).
The theoretical curves in Figs. 1–3 correspond to the following values of parameters:
$$Q_0^2=1\text{ GeV}^2,\mathrm{\Lambda }=0.25\text{ GeV}.$$
(32)
The parameter $`n_g=6.1`$ in Eq. (19) is taken from one of the MRST sets of parton distributions .
It should be noted that the strong $`Q^2`$ dependence seen by H1 and ZEUS in the Breit frame has nothing to do with the $`Q^2`$ dependence of $`n^{DIS}(W,Q^2)`$ in the c.m.s. (see Fig. 1); to a large extent it has a kinematical origin. The point is that an increase of $`Q^2`$ at fixed $`W`$ is equivalent to an increase of $`x`$. As a result, the current region of the Breit frame (29) enlarges. Thus, the rapid growth of hadronic multiplicity in the Breit frame in $`Q^2`$ (at fixed W) reflects a strong increase of the hadronic spectrum towards the central region.
As for the increase of $`n^{DIS}(x,Q^2)`$ in $`Q^2`$ at fixed $`x`$ in the Breit frame, it has been found that it is similar to that in $`e^+e^{}`$ annihilation at high $`Q^2`$, while there is a discrepancy between DIS and $`e^+e^{}`$ data at low $`Q^2`$ . Within our approach, it can be understood as follows. On the one hand, the increase of $`Q^2`$ results in an increase of the height of the spectrum (because $`W`$ grows). On the other hand, the position of the spectrum, $`y_0`$ as defined in (9), tends towards the region of positive rapidities.
These two phenomena go in opposite directions. At $`Q^2`$ high enough, $`y_0`$ varies very slowly with $`Q^2`$ (12). As a result, the rise of hadronic multiplicity in the Breit frame is analogous to that in $`e^+e^{}`$ annihilation. At low $`Q^2`$, $`y_0`$ changes more significantly . This effect partially compensates the growth of the spectrum and there appears a significant difference between DIS and $`e^+e^{}`$ data.
Finally, it is interesting to analyse the case when $`x`$ increases while $`Q^2`$ remains fixed. In this, $`W`$ decreases, which results in a rapid decrease of the spectrum. At the same time, however, the current region in the Breit frame becomes larger in accordance with formulas (29), (30). The effects are of the same order but opposite in sign. The ZEUS data in the Breit frame (see Table 2 in Ref. ) show that there is a slow rise of the hadronic multiplicity in the current hemisphere in the variable $`x`$ at different fixed values of $`Q^2`$.
## 4 Conclusions
In this paper we have presented the results of QCD calculations of the multiplicity of charged hadrons in the current hemisphere of DIS. Both the c.m.s. and the Breit frame are considered. We have shown that the H1 data are in agreement with Yang’s hypothesis and our QCD predictions. Namely, the efficiency of high energy collisions does weakly depend on the hard scale of the process (momentum transfer $`Q^2`$). It means that at fixed energy the efficiency of a particle production in hard processes increases with the shrinking of the interaction region ($`1/Q`$, in DIS).
The observed rise of the hadronic multiplicity with $`Q^2`$ in the current region of the Breit frame has both a dynamical and kinematical origin. This is why a direct comparison of available DIS data in the Breit frame with $`e^+e^{}`$ data is not completely correct.
New data from HERA on the hadronic multiplicity in the current region of the c.m.s. as well as measurements of the total multiplicity in DIS as functions of two variables ($`Q^2`$ and $`W`$/$`x`$) would be very important.
## Acknowledgements
We would like to thank Professors G. Altarelli, S. Catani, A.B. Kaidalov, M. Mangano and G. Veneziano for useful discussions. One of us (A.V.K.) is indebted to Professors G. Altarelli, S. Catani and M. Mangano for their support. We also thank Professor L.M. Shcheglova for stimulating discussions of the HERA data and Professor A.L. Kataev for sending us information on parton distributions.
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# Detection of weak gravitational lensing distortions of distant galaxies by cosmic dark matter at large scales
## Wide-field imaging with control of systematic shape errors
We observed three “blank” (i.e., not containing any known mass concentrations) fields at 23<sup>h</sup>48<sup>m</sup>, +0057 J2000; 04<sup>h</sup>29<sup>m</sup>, -3618; and 11<sup>h</sup>38<sup>m</sup>, -1233 over a period of several years, using the Big Throughput Camera, an array of four large, blue-sensitive CCDs at the Cerro Tololo Inter-American Observatory’s 4-m Blanco telescope. Constructed originally for weak lens observations, this camera covers a 35 arcminute field of view with 0.43 arcsecond pixels. We took multiple 500 second exposures shifted by 5-7 arcminutes and combined them to cover a 43 arcminute field. Before combining, we took several steps to reduce systematic errors arising from the optical system. First, we registered all the images onto a common linear coordinate system, free of the known radial distortion of the telescope optics. We then used the shapes of stars, which are foreground point sources free of the gravitational lensing effect, to correct any additional anisotropies in the point-spread function (the response of the optical system to point sources), such as those due to astigmatism and guiding errors. As described below, our observations covered multiple wavelengths; this enables filtering for certain types of stars and for distant galaxies.
The shape of a star or galaxy can be described by its second central moments, $`I_{xx}\mathrm{\Sigma }Iwx^2`$, $`I_{yy}\mathrm{\Sigma }Iwy^2`$, and $`I_{xy}\mathrm{\Sigma }Iwxy`$, where $`I(x,y)`$ is the intensity distribution above the night sky level, $`w(x,y)`$ is a weight function, the sum is over a contiguous set of pixels defined as belonging to the galaxy, and the coordinate system has been translated so that the first moments vanish. The second moments can be combined to form a size, $`I_{xx}+I_{yy}`$, and two components of a pseudo-vector ellipticity, $`e_1(I_{xx}I_{yy})/(I_{xx}+I_{yy})`$ and $`e_22I_{xy}/(I_{xx}+I_{yy})`$, which vary in the range $`[1,1]`$ (ellipticity in its colloquial sense is the amplitude of this pseudo-vector, $`ϵ\sqrt{e_1^2+e_2^2}`$ with its range $`[0,1]`$). Traditional intensity-weighted moments are calculated with $`w=1`$, but this produces ellipticity measurements with noise properties that are far from optimal or even divergent. In cases of white noise the formal optimal weight for an elliptical source is a noise-free image of that elliptical source. In the absence of such an image, weak lensing measurements are generally made with circular Gaussian weights. We use an elliptical Gaussian as the weight function, which places more weight on the high-signal-to-noise inner parts of the galaxy image, and is nearly optimal for most point-spread functions and for typical exponential galaxy profiles.
The moments of the Gaussian weight ellipse are iterated (from initial values provided by unweighted moments) to match the size and shape of the object, in order to obtain the highest possible signal-to-noise and to insure that the measured ellipticity is not biased toward the shape of the weight function. This “adaptively weighted moments” scheme has been extensively tested on simulated and real data, and has been shown to be unbiased. On simulated data, the algorithm recovers a somewhat higher fraction of the artificially induced shear than does simple intensity weighting. However, our final results do not depend on this particular weighting scheme. Its real benefit lies in rejection of peculiar objects, the vast majority of which are overlapping galaxies seen in projection. If the final centroid of a galaxy differs significantly from the starting centroid, that galaxy is rejected. If the centroid does not shift, the galaxy is accepted (and probably suffers little contamination by its neighbor). Object candidates are also rejected if the centroid or the ellipticity fails to converge; if they are too near the edge of the image; if the size grows too large; or if the moments are negative. About one-third of candidates found by the detection software (which provides the initial unweighted moments, and can “detect” occasional noise peaks) are rejected. For candidates which survive, the measurement error in the final ellipticity is accurately estimated by propagating the Poisson photon noise through the moment equations.
We used foregound stars at many positions in the field of view to measure and correct for systematic ellipticity error. Stars are distinguished from galaxies by their clear separation at the bright end of a size–flux density diagram. We identified roughly 100 such stars on each exposure of each CCD and made a least-squares fit (with 3$`\sigma `$ clip) of a second-order polynomial to the spatial variation of their ellipticity components $`e_1^{}(x,y)`$ and $`e_2^{}(x,y)`$, which would be zero at all points in an ideal observation free of point-spread function anisotropy. Fischer and Tyson have shown that nonzero $`e_1^{}`$ and $`e_2^{}`$ can be cancelled by convolution with a small (three pixel square) flux-conserving kernel with ellipticity components equal and opposite to those of the stars. Simulations as well as weak lensing data on clusters of galaxies show that faint galaxy induced systematics are also removed in this process of circularising stars. We convolved each image with its resulting position-dependent circularising kernel, after which the stellar ellipticities show little variation as a function of position. We then combined the images by averaging with a 3$`\sigma `$ clip, and repeated the point-spread function rounding on the combined image (using roughly 1000 stars and a fourth-order polynomial in this case). Fig. 2 depicts the evolution of one of our worst raw images through this process.
## Catalogues of distant galaxies
We repeated the observing and image processing for each field in three wavelength bands centered on 450 nm, 650 nm, and 850 nm, and for two of the fields we also took 550 nm images. The mean exposure time at each wavelength was 3400 sec. In each field, we used standard software to identify object positions and fluxes on the 650 nm image (which is the deepest image in each field), yielding roughly 150,000 objects per field. We have confirmed the robustness of the weighted intensity moments in our detected object catalogues by using different (FOCAS with adaptive circular kernel ) detection and evaluation software. At each object’s position, we evaluated the weighted moments at each wavelength, retaining only the measurements which the iterative weighted moment algorithm did not flag as suspect. Measurements with small sizes ($`I_{xx}<1.0`$ or $`I_{yy}<1.0`$) were also excluded as suspect. The result is a list of multiple independent ellipticity measurements (and corresponding estimated measurement errors $`\sigma _i`$) for each object.
We then computed the best estimate of each galaxy’s ellipticity by averaging the remaining measurements, weighted inversely by their estimated errors. If either of the ellipticity components at any wavelength deviated from this mean by more than $`3\sigma _i`$, that wavelength was eliminated and the process repeated. This step thus eliminates individual galaxy ellipticity measurements at wavelengths at which objects were noisy or blended, and it also reduces the systematic errors because the images at different wavelengths do not share the same residual point-spread function anisotropy. Finally we rejected objects with $`ϵ>0.6`$ as likely to be blends of more than one object, and applied flux density criteria ($`1.8\mu Jy>F_\nu >0.11\mu Jy`$ through the 650 nm filter, 23–26 R magnitude) to select objects likely to be distant galaxies. We use these same selection criteria in calibrating the typical redshift of the background galaxies (below). The final catalogues contain about 45,000 galaxies in each field. A visual inspection of the final catalogues indicates that they are free of spurious objects such as bits of scattered light around bright stars.
These observed ellipticities must be corrected for the overall broadening effect of the point-spread function, which makes elliptical galaxies appear more circular even if the point-spread function itself is perfectly isotropic. To calibrate this effect, we took a deep image with a very small point spread (the Hubble Deep Field South) and convolved it to the resolution of our final images, which is 1.07–1.25 arcsec as measured by the full width at half-maximum at 650 nm. (The resolution on individual exposures, or “seeing” was better, but the stellar size is larger in the final image with systematic shape errors removed from the point-spread function.) While some isolated galaxies became broader and less elliptical as predicted, most merged with their neighbors, producing many more elliptical objects than predicted and preventing the construction of a clear relationship between observed and true ellipticity for individual galaxies.
Instead we calibrated the fraction of cosmic shear recovered, as a function of resolution, from the ensemble of galaxies matching our selection criteria. We induced a known shear into the Hubble Deep Field South, convolved to the desired resolution, applied the same galaxy measurement and selection routines (at 650 nm only), and measured the mean ellipticity of the resulting sample. We averaged over repeated shears in several different directions to assess the measurement errors. The ratio of induced to recovered ellipticity was $`4.5\pm 0.5`$, with no clear trend as a function of resolution. The lack of such a trend would be quite surprising for isolated galaxies, but the coalescence of galaxy images appears to be the dominant effect. In the fairly small range of 1.07–1.25 arcsec resolution, this effect does not change the recovery factor by more than the measurement error of 0.5, so we adopt 4.5 as an overall ellipticity recovery factor.
## Ellipticity correlations of distant galaxies
Miralda-Escudé has defined two physically revealing ellipticity correlation functions. In this approach, the ellipticity components of a galaxy $`i`$ are calculated not with respect to the arbitrary $`x`$ and $`y`$ axes of the image, but with respect to the line joining it to another galaxy $`j`$ (Fig. 3). Averaging over all galaxies $`i`$ and $`j`$ separated by angle $`\theta `$ on the sky, the correlations $`\xi _1(\theta )e_{1i}e_{1j}`$ and $`\xi _2(\theta )e_{2i}e_{2j}`$ have a unique signature in the presence of gravitational lensing, explained in detail in Fig. 3. We recently reported the detection of a cosmic shear signal in the quadrature sum of these correlation functions.
Fig. 4 shows the ellipticity correlations for each of the three fields in the angular separation range 2-36 arcmin (top panels). The plotted errors indicate 68% confidence intervals determined from 200 bootstrap-resampled realizations of the final galaxy catalogue in each field. Note that the measurements in different angular separation bins are not statistically independent, but $`\xi _1`$ and $`\xi _2`$ are independent from each other, as are the three fields. At $`\theta =6.1`$ arcmin, the confidence that $`\xi _1>0`$ is 97%, 99.5%, and 99.5% for the three fields in the order shown in Fig. 4. Similarly, the confidence that $`\xi _2>0`$ at the same angular scale is 87.5%, $`>99.5\%`$, and 97% respectively. Some “cosmic variance”, or real systematic differences among fields of this size, is expected, but the statistical errors are too large to examine this effect. We plot the average over the three fields in the lower panels of Fig. 4, with $`1\sigma `$ errors in the mean derived from the variance among the fields (black points and errors). The signature of gravitational lensing by large-scale structure is evident: $`\xi _1`$ declines as the angular scale increases, but is positive at all scales, while $`\xi _2`$ matches $`\xi _1`$ at small scales but drops below zero at large scales. This result is robust: Similar, but lower signal-to-noise, profiles are obtained if we use unweighted moments or moments from the 650 nm images only.
We performed several tests for systematic errors. The effects of residual point-spread function anisotropy are demonstrated by plotting the correlation functions of the stars (blue in Fig. 4). These are far closer to zero than are the galaxy correlations. Only $`\xi _1`$ in the innermost bin has an apparently significant stellar correlation. To test the effect that this might have on the galaxy correlations, we computed the star-galaxy correlations $`e_{1,star}e_{1,gal}`$ and $`e_{2,star}e_{2,gal}`$ (green in Fig. 4). The star-galaxy correlations are extremely close to zero in this bin. There are also tests involving the galaxy sample alone. The cross correlation $`\xi _3\frac{1}{2}(e_{1i}e_{2j}+e_{2i}e_{1j})`$ should vanish in the absence of systematic errors (red in Fig. 4). The result is reassuringly close to zero. The plotted errors for $`\xi _3`$ can also be taken as an indicator of the statistical error associated with the number and distribution of galaxies included in the catalogues (but reduced due to the averaging of two functions in $`\xi _3`$). This estimate of statistical error agrees roughly with that shown for $`\xi _1`$ and $`\xi _2`$. Finally, the weak lensing signature disappears if we randomise the galaxy positions.
Apart from these null tests, there are also affirmative tests. One test is to take similar data centered on a cluster of galaxies of known mass. A 650-nm image of massive cluster at redshift 0.45, taken with the same camera and processed in the same way, exhibits correlation functions ($`\xi _1`$ and $`\xi _2`$) of the expected angular dependence and of larger amplitude than in any of the three blank fields, despite likely contamination of the galaxy sample by cluster members. $`\xi _3`$ also vanishes in this field. Another test involves inverting the background galaxy ellipticity distribution to yield a map of projected mass in each “blank” field. We find occasional mass concentrations which can often be identified with likely foreground clusters, but no linear features or pileups at the edges of the image which might indicate problems in the background galaxy catalogues. Furthermore, when a mass map is made using only those galaxies likely to be behind a serendipitous cluster (based on colour information), the lensing signal from that cluster increases markedly. This corroborates the idea that the correlation functions are accumulating over many sources and many overdensities spread throughout the line of sight. All these tests indicate that we have indeed measured cosmic shear in our “blank” fields and that contamination from surviving systematic error is low. We now turn to comparisons with theoretical predictions of this effect.
## Comparison with theoretical predictions
Ellipticity correlations increase strongly with background galaxy redshift, so we must first constrain the source redshift distribution $`N(z)`$. Very little is known about the redshift distribution of galaxies as faint as those used here, so we assume a simple model $`N(z)z^2exp(z/z_0)`$, and adjust $`z_0`$ to match weak gravitational lensing observations of a high-redshift galaxy cluster of known velocity dispersion (MS1054 at $`z=0.83`$). We observed this cluster with the same camera and telescope and reduced the data in the same way as for the blank fields, and compared the faint galaxy ellipticities (tangential to the cluster center) to that expected for a range of $`z_0`$. We found that $`z_0=0.5`$ was the best match.
This model $`N(z)`$ was used as input to a cold dark matter simulation code by W. Hu and J. Miralda-Escudé (see ref ), which computes the shear power spectrum and correlation function for any given cosmology, using the prescription of Hamilton et al. and Peacock & Dodds to calculate the mass power spectrum in the non-linear regime when the growth of gravitationally collapsed dark matter structures modified the mass spectrum. Results were obtained for three cosmological models and are plotted along with our seeing-corrected measurements on a logarithmic scale in Figure 5. Two current models were normalised to the microwave background fluctuations (COBE) at large angle and to local galaxy cluster abundance (assuming mass traces light) at small angle: an open universe with $`\mathrm{\Omega }_{\mathrm{matter}}=0.45`$ (orange in Figure 5), and a flat universe dominated by a cosmological constant $`\mathrm{\Lambda }=0.67`$ (green, solid line). The agreement between the data and the two viable cosmological models is impressive for a first measurement. For comparison purposes we also show the old standard cold dark matter flat cosmology (blue), which is only COBE normalised. (A full listing of the parameters used in these models is shown in Table 1.) To illustrate the effect of varying $`N(z)`$, we also plot the $`\mathrm{\Lambda }`$-dominated cosmology with $`z_0=0.3`$ (green, dotted line). Since our model $`N(z)`$ peaks at $`z=2z_0`$, this lowers the typical redshift from 1.0 to 0.6. Decreasing $`z_0`$ decreases the amplitude of the correlations, but has little effect on their shapes. The uncertainty in $`N(z)`$ implies a factor of several uncertainty in the amplitude of the correlation, and is by far the dominant calibration error.
Despite this uncertainty, COBE-normalised standard cold dark matter is ruled out by the measured values of $`\xi _1`$. While this is not surprising, it is the first cosmological constraint from wide-field weak lensing, and it agrees with several other methods which disfavor this model. The other two models are consistent with the data at the $`3\sigma `$ level. The indication of a low $`\mathrm{\Omega }_{\mathrm{matter}}`$ universe here is in agreement with a remarkable array of independent methods, including type Ia supernovae, cosmic microwave background anisotropies, cluster baryon fraction together with cluster mass (lensing) and primeval deuterium, and the age of the oldest stars coupled with the Hubble constant. However, the shape of $`\xi _2`$ is not a good fit to either of these two model cosmologies, which are based on a single power-law mass spectrum. If confirmed by further data, this would suggest the need for a more complicated mass spectrum.
This technique can further distinguish between open and $`\mathrm{\Lambda }`$-dominated universes if extended to the somewhat larger angular scales where those cosmologies predict $`\xi _2`$ will cross zero as shown in Figure 5. A survey of many $`2^{}\times 2^{}`$ fields now underway will rule out one or more of these cosmologies at the $`8\sigma `$ level at 10 arcmin angles ($`3\sigma `$ level for a differential measure of the slope of the power spectrum). Separating the background galaxies into discrete redshift bins based on multi-colour photometry will enable measurement of the ellipticity correlation (or equivalently the dark matter power spectrum) as a function of cosmic time; wide-field weak lensing surveys deep enough to identify galaxies at $`z2`$ and measure their shapes will constrain several cosmological parameters. Ultimately, the combination of all the power spectrum probes (lensing, cosmic microwave background, galaxy distributions, and peculiar velocities) will tightly constrain theories of the origins of fluctuations in the early universe and their growth into galaxies and large-scale structure.
### Acknowledgements
We gratefully acknowledge help from Wayne Hu and Jordi Miralda-Escudé on theoretical predictions of several cosmological models. We thank Steven Gentile for his artwork, and the staff of CTIO for their help with the BTC project and for their upgrading and maintenance of the delivered image quality of the Blanco telescope. Cerro Tololo Inter-American Observatory is a division of National Optical Astronomy Observatory (NOAO), which is operated by the Association of Universities for Research in Astronomy, Inc., under Cooperative Agreement with the National Science Foundation. BTC construction was partially funded by the NSF.
| Model (Fig 5 colour) | $`\mathrm{\Omega }_b`$ | $`\mathrm{\Omega }_{\mathrm{matter}}\mathrm{\Omega }_b`$ | $`\mathrm{\Omega }_\mathrm{\Lambda }`$ | $`H_0`$ | n | $`\sigma _8`$ | normalization |
| --- | --- | --- | --- | --- | --- | --- | --- |
| Standard cold dark matter (blue) | 0.05 | 0.95 | 0 | 50 | 1.0 | 1.17 | COBE only |
| $`\mathrm{\Lambda }`$-dominated, flat (green) | 0.039 | 0.291 | 0.67 | 70 | 0.94 | 0.84 | COBE+clusters |
| Open universe (orange) | 0.045 | 0.405 | 0 | 65 | 1.01 | 0.71 | COBE+clusters |
These cosmological models were chosen in order to put our ellipticity correlation measurements in context. The old standard cold dark matter model in which the universe is nearly closed by cold dark matter, is also disfavored in other observations. Its rms mass contrast is normalized to that found 300,000 years after the Big Bang via the cosmic microwave background radiation fluctuations observed with the COsmic Background Explorer satellite ($`COBE`$). The other two models agree with a wide variety of observations, but only the cosmological constant ($`\mathrm{\Lambda }`$-dominated, flat) cosmology also agrees with the recent evidence from supernova studies for accelerated expansion. $`\mathrm{\Omega }_b`$ is the fraction of critical density in to ordinary (baryonic) matter. $`\mathrm{\Omega }_{\mathrm{matter}}`$ is the fraction in all matter (mostly dark matter), and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ is the fraction in dark energy (the cosmological constant). $`H_0`$ is the Hubble constant in units of km s<sup>-1</sup> Mpc<sup>-1</sup>. The power spectrum $`P(k)`$ of mass density fluctuations is often plotted in terms of inverse size: the wave number k is inversely proportional to length. The parameter n is the slope of the primeval density power spectrum as a function of k: $`P(k)k^n`$. For a scale-free power spectrum of density fluctuations, n = 1. The parameter $`\sigma _8`$ is the current rms mass contrast in a random sphere of radius 8($`100/H_0`$) Mpc compared with that for numbers of galaxies. The choice of n=1 and $`COBE`$ normalization for standard cold dark matter results in too much mass fluctuation on galaxy cluster scales. By adjusting the slope n and current rms mass contrast $`\sigma _8`$, models can be forced to fit the rms mass contrast now on galaxy cluster scales as well as the $`COBE`$ normalization.
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# Systematic uncertainties in gravitational lensing models: a semi-analytical study of PG1115+080
## 1 Introduction
Gravitationally lensed systems are powerful probes of galactic potentials and the scale of the universe. The advantage over the traditional stellar dynamical method is that from the image positions we can measure the shape and the mass of the dark halo well beyond the half-light-radius of a faraway lens galaxy, whether it is virialized or not. The time delay between two images is a measurement of the difference in the length of the two bent light paths, and scales with the distances to the lens and source. So the Hubble constant can be constrained once the redshift of the lens and the source and the time delay are measured. This way of getting $`H_0`$ has the advantage that the underlying physics (general relativity) can be rigorously modeled. The limitation is that there is often a sequence of lens models that can fit the image positions.
Presently about twenty strongly lensed systems are known, half of them being quadruple-imaged quasars and half being double-imaged quasars (e.g., Keeton & Kochanek 1996). We will concentrate on quadruple-imaged systems. They are better constrained than double-imaged systems, since the lens model needs to fit more image positions and also the ratios of time delays between any two pairs of images. Quadruple systems typically involve a quasar source well-aligned with the center of the lens potential well. Presently only two such systems (PG1115+080 and B1608+656) have accurately determined image positions and time delays.
A significant amount of numerical computation is usually required to invert image positions to intrinsic parameters of the potential (cf. Schneider, Ehlers & Falco 1992). The degeneracy of the resulting potential is often not fully explored because of the need to cover a large parameter space, particularly for flattened potentials. Previous authors have often restricted their studies to isothermal or power-law spherical models (Evans & Wilkinson 1998 and references therein), and elliptical models (Witt & Mao 1997 and references therein) and other simple models (Kassiola & Kovner 1993, 1995) with or without external shear. The fully general non-parametric method, e.g., the pixelated lens method of Williams & Saha (2000), is very powerful in demonstrating the complete range of the degeneracy in the lens models, but it involves significant amount of numerical computation and does not provide a clear insight to the relations between the characteristic parameters of the lens. For these reasons it is still desirable to find analytical, yet general potentials, which allow a quick exploration of the model parameter space. For example, it is of interest to generalize the analytical work of Witt & Mao (1997) to non-elliptical lenses. It would also be interesting to find analytical expressions for the time delay in these general lenses, which could help us to understand how the radial profile and lens shape affect the predictions on the Hubble constant.
Here we study a broad class of analytical models with non-axisymmetric, non-elliptical shape and semi-power-law radial profile (§2), and show how to calculate the lens shape and radial profile parameters directly from the image positions (§3). We apply the models to PG1115+080 (§4) and show that the images can be fit perfectly by a large range of lens models (§5). We summarize our results in §6 and conclude with the implications on the Hubble constant.
## 2 Lens equation in a general class of models
### 2.1 Decoupling of angular dependence and radial profile
Any two-dimensional lens potential can be cast in the following form,
$$\psi (x,y)=\psi (\omega ),\omega =\omega (r,\theta ),$$
(1)
where $`\psi `$ has the dimension of square arcsec, $`(x,y)`$ defines a rectangular coordinate system (in units of arcsecs to the West and North of the lens galaxy center) and $`(r,\theta )`$ the corresponding polar coordinate system with
$$(x,y)=(r\mathrm{sin}\theta ,r\mathrm{cos}\theta ),$$
(2)
where $`\theta `$ is the position angle, counterclockwise from North. Unless otherwise specified we shall follow the notations of Schneider et al. (1992). Here $`\omega `$ is defined to have the dimension of radius, so that constant $`\omega `$ curves correspond to equal-potential contours, and define the shape and flattening of the potential. The radial profile of the potential is a smooth function $`\psi (\omega )`$ of the radius $`\omega `$. We can also define
$$m(\omega )\omega \frac{d\psi }{d\omega }$$
(3)
as the mass (in units of square arcsec) enclosed inside the radius $`\omega `$ (in units of arcsec). For example, we have $`m=\alpha \psi `$ for a power-law model $`\psi \omega ^\alpha `$ with slope $`\alpha `$. For a source at redshift $`z_s`$ and angular distance $`D_{os}(z_s)`$, and a lens at redshift $`z_l`$ and distance $`D_{ol}(z_l)`$, the physical mass $`M(\omega )`$ is related to $`m(\omega )`$ by the scaling
$$M(\omega )\pi \mathrm{\Sigma }_cm(\omega ),$$
(4)
where
$$\mathrm{\Sigma }_c\frac{c^2}{4\pi G}\frac{D_{ol}D_{os}}{D_{ls}}\left(\frac{1\mathrm{radian}}{206265\mathrm{arcsec}}\right)^239\frac{D_{ol}D_{os}}{D_{ls}}$$
(5)
is the critical density in units of $`M_{}/\mathrm{arcsec}^2`$, and $`D_{ls}(z_l,z_s)`$ is the relative distance of the lens and source and all distances are in units of parsec.
A light ray from a source at $`(r_s,\theta _s)`$, being deflected to a direction $`(r,\theta )`$ by the lensing galaxy with potential $`\psi (\omega )`$ located at the origin, will experience a time delay $`\mathrm{\Delta }t`$ given by
$$\mathrm{\Delta }t(r,\theta )=h^1\tau _{100}(z_l,z_s)\left[\frac{1}{2}r^2rr_s\mathrm{cos}(\theta \theta _s)+\frac{1}{2}r_s^2\psi (\omega )\right],$$
(6)
where $`h`$ is the Hubble constant $`H_0`$ rescaled to $`100`$ km/s/Mpc. A characteristic value for the time delay is
$$\tau _{100}(z_l,z_s)\frac{D_{ol}D_{os}}{D_{ls}}\frac{1+z_l}{c},\text{for }H_0=100\text{ km/s/Mpc}.$$
(7)
According to Fermat’s principle, the images lie at the minimum of $`\mathrm{\Delta }t`$, so the lens equation is given by
$`{\displaystyle \frac{}{r}}\mathrm{\Delta }t`$ $`=`$ $`rr_s\mathrm{cos}(\theta \theta _s)m(\omega )_r\mathrm{ln}\omega =0`$ (8)
$`{\displaystyle \frac{}{r\theta }}\mathrm{\Delta }t`$ $`=`$ $`r_s\mathrm{sin}(\theta \theta _s){\displaystyle \frac{m(\omega )}{r}}_\theta \mathrm{ln}\omega =0.`$ (9)
Interestingly, the radial part $`m(\omega )`$ can be eliminated by simply combining the two equations,
$$\frac{r_s\mathrm{sin}(\theta \theta _s)}{rr_s\mathrm{cos}(\theta \theta _s)}=\frac{_\theta \mathrm{ln}\omega }{r_r\mathrm{ln}\omega },$$
(10)
similar as in Witt & Mao (1997). This implies that there is a relation linking the image positions to the shape of the potential directly, independent of the radial profile.
### 2.2 Property of the image positions: the semi-hyperbolic curve
For simplicity we shall concentrate on self-similar models with
$$\omega (r,\theta )=rf(\theta ),$$
(11)
where $`f(\theta )`$ defines the shape of the equal potential contours. In principle the shape function can be bi-symmetric or lopsided as long as the corresponding surface density is positive everywhere in the lens plane. To be specific, we will restrict our discussions to bi-symmetric potentials with an angular part
$$f(\theta )=\left|1\delta \mathrm{cos}2\theta ^{}\right|^\beta ,$$
(12)
which is a three-parameter $`(\beta ,\delta ,\theta _p)`$ function of the angle $`\theta `$, where $`\beta `$ is a constant, the parameter $`\delta `$ is a flattening indicator, and $`\theta _p`$ is the position angle of a principal axis of the potential. The angle $`\theta ^{}`$ is the azimuthal angle $`\theta `$ except for a rotation with
$$\theta ^{}\theta \theta _p,(x^{},y^{})=(r\mathrm{cos}\theta ^{},r\mathrm{sin}\theta ^{})$$
(13)
so that $`(x^{},y^{})`$ defines a rectangular coordinate system with the axes coinciding with the principal axes of the lens. For example, the source at
$$(x_s,y_s)=(r_s\mathrm{sin}\theta _s,r_s\mathrm{cos}\theta _s)$$
(14)
in the original rectangular coordinate system would be at
$$(x_s^{},y_s^{})=(r_s\mathrm{cos}\theta _s^{},r_s\mathrm{sin}\theta _s^{}),\theta _s^{}\theta _s\theta _p,$$
(15)
in the rotated rectangular coordinate system.
With these we can compute the right-hand part of eq. 10,
$$\frac{_\theta \mathrm{ln}\omega }{r_r\mathrm{ln}\omega }=\frac{d\mathrm{ln}f(\theta )}{d\theta }=\frac{1}{v\mathrm{cot}\theta ^{}+u\mathrm{tan}\theta ^{}},$$
(16)
where
$$u\frac{1+\delta }{2\mathrm{\Delta }},v\frac{1\delta }{2\mathrm{\Delta }},\mathrm{\Delta }=2\beta \delta ,$$
(17)
are shape indicators just like $`\beta `$ and $`\delta `$. Substituting in eq. 10, and rewriting the image radius $`r`$ as a function of $`\theta ^{}`$, we find the images fall on a family of curves $`r=r(\theta ^{})`$ defined by
$$r(\theta ^{})=r_s\left[\mathrm{cos}(\theta ^{}\theta _s^{})+\mathrm{sin}(\theta ^{}\theta _s^{})\left(v\mathrm{cot}\theta ^{}+u\mathrm{tan}\theta ^{}\right)\right].$$
(18)
An alternative expression for these curves can be obtained by expanding the sinusoidal terms in eq. 18 so that
$$r=(X_a\mathrm{cos}\theta ^{}+Y_a\mathrm{sin}\theta ^{})+\left(\frac{X_b}{\mathrm{cos}\theta ^{}}\frac{Y_b}{\mathrm{sin}\theta ^{}}\right),$$
(19)
which is now linear in four new parameters $`(X_a,Y_a,X_b,Y_b)`$; these parameters are related to the lens shape parameters $`(v,u)`$ and the source position $`(x_s^{},y_s^{})`$ by
$$X_a\left(1+vu\right)x_s^{},Y_a\left(1u+v\right)y_s^{},X_bux_s^{},Y_bvy_s^{}.$$
(20)
These curves (cf. eq. 18) have the nice property that they go through all image positions independent of the radial profile of the lensing galaxy. An example is the semi-hyperbolic curve in Fig. 1. The curve is determined by the source parameters $`(r_s,\theta _s^{})`$, the lens shape parameters $`(v,u)`$ and the lens position angle $`\theta _p`$. The radial profile can take any general physical profile, isothermal or power-law.
The boxiness parameter $`\beta `$ is such that the shape function $`f(\theta )`$ reduces to the usual elliptical form when $`\beta =\frac{1}{2}`$ (Witt & Mao 1997). In the case that $`\psi (\omega )`$ is linear in $`\omega `$, the models reduce to the simple models of Kassiola & Kovner (1995) when $`\beta =1`$. Interestingly, elliptical models with $`\beta =\frac{1}{2}`$ have $`uv=\frac{\delta }{\mathrm{\Delta }}=\frac{1}{2\beta }=1`$ (cf. eq. 17), hence $`X_a=Y_a=0`$, and eq. 19 reduces to
$$1=\frac{X_b}{x^{}}\frac{Y_b}{y^{}}$$
(21)
after applying $`x^{}=r\mathrm{cos}\theta ^{},y^{}=r\mathrm{sin}\theta ^{}`$. This equation prescribes a hyperbolic curve, which is consistent with Witt (1996) and Witt & Mao’s (1997) finding that all four image positions and the source position lie on a certain hyperbolic curve. A hyperbolic curve has a maximum of 5 free parameters, thus they cannot fit image positions of a general quadruple system; four image positions yield a minimal of eight constraints. Experience with fitting several quadruple systems (G2237+0305, CLASS1608+656, HST12531-2914) shows that models with $`|\beta |1/2`$ often give unphysical mass-radius relations $`m(\omega )`$. We find that $`\beta =\frac{1}{8}`$ gives a fair approximation to realistic models. These models have non-elliptical contours, and often yield physical density distributions. Fig. 2 shows that they also cover a sufficiently wide range of axis ratios for the potential and the density so that we can explore the shape of the lens galaxy in fitting the image positions. The expressions for the axis ratios of a power-law lens are given in Appendix A.
## 3 Results
### 3.1 Lens shape directly from fitting image positions
Our models can be used to fit image positions and derive lens shape parameters ($`\delta ,\beta `$) and the source position $`(r_s,\theta _s)`$, free from assumptions of the lens radial profile, but subject to the assumption that the lens’ angular profile obeyes eq. 12. The four unknowns can be derived from the four observed image positions, i.e., the *eight* observables $`(r_i,\theta _i)`$ with $`i=1,2,3,4`$. The position angle $`\theta _p`$ of the lens principal axis is treated as a free variable.
The procedure is simple. First substitute the four observed image positions in eq. 19 to obtain the following four linear equations
$$\frac{\mathrm{cos}\theta _i^{}}{r_i}X_a+\frac{\mathrm{sin}\theta _i^{}}{r_i}Y_a+\frac{1}{r_i\mathrm{cos}\theta _i^{}}X_b\frac{1}{r_i\mathrm{sin}\theta _i^{}}Y_b=1,i=1,2,3,4,$$
(22)
of the four unknown parameters $`(X_a,Y_a,X_b,Y_b)`$. After solving these, either analytically or numerically, these parameters are substituted in eqs. 2015, and 17 to yield the source position and the lens shape in terms of $`(r_s,\theta _s,\delta ,\beta )`$. In fact, we can recast eq. 20 to a set of four simple linear equations of four new unknowns $`(u,v,1/x_s^{},1/y_s^{})`$ by moving the terms $`x_s^{}`$ and $`y_s^{}`$ to the left hand side of the equations, i.e.,
$$X_a/x_s^{}v+u=1,Y_a/y_s^{}v+u=1,X_b/x_s^{}u=0,Y_b/y_s^{}v=0.$$
(23)
The lens shape parameters $`(\delta ,\beta )`$ can then be computed from
$$\delta =\frac{uv}{u+v},\beta =\frac{1}{2(uv)},$$
(24)
and the source position $`(r_s,\theta _s)`$ from
$$r_s=\sqrt{x_{s}^{}{}_{}{}^{2}+y_{s}^{}{}_{}{}^{2}},\theta _s=\mathrm{arctan}(y_s^{},x_s^{})+\theta _p.$$
(25)
Thus we have effectively reduced the problem of fitting image positions to successively solving linear equations, which is a straightforward task.
### 3.2 Non-parametric radial profile and power-law slope
The radial part of the potential can also be extracted from eq. 8 without parameterization. At the positions of the images we have
$$m(\omega )=r^2rr_s\mathrm{cos}(\theta \theta _s)=x^2+y^2xx_syy_s,\omega =rf(\theta ),$$
(26)
where we have used $`_r\mathrm{ln}\omega =r^1`$. Thus we have obtained a mass-radius relation directly from the observed image positions, assuming that the source position $`(x_s,y_s)`$ and the flattening and position angle of the potential $`(\delta ,\theta _p)`$ have been determined by fitting the curve (cf. eq. 18). Interestingly, in the limit that the source is at the center of a circular lens, we have $`r_s0`$, $`\beta 0`$ and $`mr^2`$.
It is useful to characterize the radial profile of a lens galaxy, which is generally not a power-law, by an effective power-law slope $`\alpha (\omega )`$, which varies with the radius $`\omega `$ except for scale-free models. There are several ways of estimating the characteristic power-law slope. Taking any two images $`i`$ and $`j`$, we can form a characteristic power-law slope $`\alpha _{ij}`$ from the mass $`m_i`$ and $`m_j`$ at the image radii $`\omega _i`$ and $`\omega _j`$ with
$$\alpha _{ij}\frac{\mathrm{log}m_i\mathrm{log}m_j}{\mathrm{log}\omega _i\mathrm{log}\omega _j}=\frac{2\mathrm{log}\frac{r_i}{r_j}+\mathrm{log}\frac{1r_sr_i^1\mathrm{cos}(\theta _i\theta _s)}{1r_sr_j^1\mathrm{cos}(\theta _i\theta _s)}}{\mathrm{log}\frac{r_i}{r_j}+\beta \mathrm{log}\frac{1\delta \mathrm{cos}(\theta _i\theta _s)}{1\delta \mathrm{cos}(\theta _j\theta _s)}},$$
(27)
where we have used the mass-radius eq. 26.
Alternatively we can estimate the power-law slope from the observed time delay. First we rewrite the time delay eq. 6, so that the time delay, $`t_{ij}`$, between any two images $`i`$ and $`j`$ is given by
$$\frac{ht_{ij}}{\tau _{100}}=\frac{1}{2}(x_i^2+y_i^2)\frac{1}{2}(x_j^2+y_j^2)(x_ix_j)x_s(y_iy_j)y_s(\psi _i\psi _j),$$
(28)
where $`\psi _i\psi _j`$ is the difference in the lens potential between the two images. Rewriting eq. 28 we form a new estimator $`\alpha _{ij}^t`$ for the power-law slope with
$$\alpha _{ij}^t\frac{m_im_j}{\psi _i\psi _j}=\frac{(x_i^2+y_i^2)(x_j^2+y_j^2)(x_ix_j)x_s(y_iy_j)y_s}{\frac{1}{2}(x_i^2+y_i^2)\frac{1}{2}(x_j^2+y_j^2)(x_ix_j)x_s(y_iy_j)y_s\frac{ht_{ij}}{\tau _{100}}},$$
(29)
where we have used eq. 26.
Thus we have two direct estimators $`\alpha _{ij}`$ and $`\alpha _{ij}^t`$ of the radial profile, computed from the observed images and time delays. In the limit of scale-free power-law models $`\alpha _{ij}=\alpha _{ij}^t=\alpha (\omega )=cst`$. So the deviation from scale-freeness can be estimated by taking the differences such as $`\alpha _{12}\alpha _{34}`$, $`\alpha _{12}^t\alpha _{34}^t`$, or $`\alpha _{14}\alpha _{14}^t`$.
### 3.3 Hubble constant and time delay ratios
To apply the above estimates of the power-law slopes, we have assumed that we know the rescaled Hubble constant $`h`$ from independent observations. Alternatively the Hubble constant $`H_0=100h`$ can be estimated from the time delay $`t_{ij}`$ between two images. Letting $`\alpha _{ij}^t=\alpha _{ij}`$, we get
$$h=\frac{\tau _{100}}{t_{ij}}\left\{\left(\frac{1}{2}\frac{1}{\alpha _{ij}}\right)\left[(x_i^2+y_i^2)(x_j^2+y_j^2)\right]\left(1\frac{1}{\alpha _{ij}}\right)\left[(x_ix_j)x_s+(y_iy_j)y_s\right]\right\}.$$
(30)
The time delay ratio can also be predicted with
$$\frac{t_{ij}}{t_{kl}}=\frac{(\frac{1}{2}\frac{1}{\alpha _{ij}})[(x_i^2+y_i^2)(x_j^2+y_j^2)](1\frac{1}{\alpha _{ij}})[(x_ix_j)x_s+(y_iy_j)y_s],}{\left(\frac{1}{2}\frac{1}{\alpha _{kl}}\right)\left[(x_k^2+y_k^2)(x_l^2+y_l^2)\right]\left(1\frac{1}{\alpha _{kl}}\right)\left[(x_kx_l)x_s+(y_ky_l)y_s\right]},$$
(31)
which is obtained by rewriting eq. 30.
## 4 Application: the surface density and time delay models for PG1115+080
### 4.1 Data
As a simple application, we model the image positions and time delays of the well-studied quadruple system PG1115+080. This system has been extensively studied ever since the first models by Young and collaborators (1981), and has received closer attention after Schechter et al.’s (1997) measurements of its time delay. All models, except those of Saha & Williams (1997), adopt elliptical/circular shapes and a few common radial profiles, with the models of Keeton & Kochanek (1997) and Impey et al. (1998) being the most comprehensive. Only one year after the discovery of the legendary double-image radio-loud quasar Q0957+561, this system was identified as a multiple-imaged system by Weymann et al. (1980) in their survey of nearby bright QSOs. It is now known to consist of four images with the names $`A_1`$, $`A_2`$, $`B`$ and $`C`$ (with flux ratios about 4 : 2.5 : 0.7 : 1, cf. the HST observations of Kristian et al. 1993) of a radio-quiet QSO at redshift $`z_s=1.722`$. The images $`A_1`$ and $`A_2`$ are within $`0.5^{\prime \prime }`$; see the inset of Fig. 1. Interestingly, the lens galaxy is also one of the bright members of a galaxy group ($`N10`$) at redshift $`z_l=0.310`$, first mapped by Young et al. (1981). The center of the group is to the south-west of the lens, roughly at $`r_g=(20^{\prime \prime }\pm 2^{\prime \prime })`$ and $`\theta _g=(117^o\pm 3^o)`$. The lens galaxy has been resolved by both HST and the 8.2-m Subaru telescope in $`0.3^{\prime \prime }`$ seeing. It appears to be an early type galaxy with a de Vaucouleurs profile and a half-light radius of $`0.55^{\prime \prime }`$. There is no sign of differential dust-extinction in the lens galaxy. While NICMOS observations by Impey et al. (1998) show no flattening for the lens, ground infrared images by Iwamuro et al. (2000) found it to be an E1 galaxy elongated towards $`\theta _p65^o`$. Both observations reveal a $`1^{\prime \prime }`$ infrared Einstein ring connecting the four images, which is thought to be the infrared image of the QSO host galaxy. PG1115+080 is also one of the two quadruple systems where the time delay between images has been measured, the other one being the radio-loud quasar B1608+656 from the CLASS survey (cf. Fassnacht et al. 1999). Although two different sets of values are quoted in the literature (Schechter et al. 1997, Barkana 1997), the leading image is the furtherest image (the image C), and the innermost image (image B) arrives last. The time delay ratio $`r_{ABC}=t_{AC}/t_{BA}`$, for the delay between image $`C`$ and $`A_1+A_2`$ vs. image $`B`$ and $`A_1+A_2`$, provides an extra discriminator of the models; the images $`A_1`$ and $`A_2`$ are within $`0.5^{\prime \prime }`$ of each other, and the small relative delay is undetected. Schechter et al. first reported $`r_{ABC}=0.7\pm 0.3`$ from their photometric monitoring program in 1995-1996. Later analysis by Barkana (1997) found $`r_{ABC}=1.13\pm 0.18`$, after taking into account correlations of errors amongst the time delays. The delay between images B and C, $`t_{BC}=25.0\pm 1.7`$ days.
Here we illustrate the application of our models to the most recent data from Impey et al. (1998) of PG1115+080. We do not attempt to model B1608+656 because the lens appears to be a merging pair of galaxies, and the morphology is too complex for our model. We denote with the index $`i=1,2,3,4`$ the four images $`A_1`$, $`A_2`$, $`B`$ and $`C`$. All results are quoted for a standard flat universe without a cosmological constant $`(\mathrm{\Omega },\mathrm{\Lambda })=(1,0)`$. Table 1 gives the relevant quantities to calibrate our results to other universes. The angular distance from redshift $`z_1`$ to $`z_2`$ in a universe of a matter and vacuum density $`(\mathrm{\Omega },\mathrm{\Lambda })`$ times the closure density is generally given by
$$D(z_1,z_2)=\frac{c}{H_0(1+z_2)\sqrt{\mathrm{\Omega }_c}}\mathrm{sinh}\left\{\sqrt{\mathrm{\Omega }_c}_{z_1}^{z_2}𝑑z\left[\mathrm{\Lambda }+\mathrm{\Omega }_k(1+z)^2+(1+z)^3\mathrm{\Omega }\right]^{\frac{1}{2}}\right\},\mathrm{\Omega }_k=1\mathrm{\Lambda }\mathrm{\Omega }.$$
(32)
The predicted Hubble constant should be reduced from the value achieved with the standard $`(\mathrm{\Omega },\mathrm{\Lambda })=(1,0)`$ universe by 3% in the presently favored $`\mathrm{\Lambda }`$-dominated universe with $`(\mathrm{\Omega },\mathrm{\Lambda })=(0.3,0.7)`$ from surveys of distant supernovae.
### 4.2 The source position, the lens shape and mass inside images
First we solve for the lens shape and source position from the linear equations 19 and 23. The solutions for PG1115+080 are given in Table 2, sorted according to the value of $`\beta `$ or $`\theta _p`$. The resulting potential model with $`\beta =\frac{1}{8}`$ or $`\theta _p=60.5^{}`$ has a flattening of between E0 and E1, and interestingly the lens principal axis points towards the location of the galaxy group in the lens plane. Models with other values of $`\beta `$ or $`\theta _p`$ will be discussed in section 5.
To proceed with determining the lens mass at each image position, we substitute the now known flattening parameters $`(\delta ,\theta _p)`$ and the source positions $`(r_s,\theta _s)`$ in eq. 26, to predict four independent data points $`(\omega _i,m_i)`$ with $`i=1,2,3,4`$ in the radius vs. the enclosed mass plane. Fig. 3 shows the predicted lens mass enclosed at the four image positions. Note that the mass rises faster than the light, implying a growing dark mass component at large radius; the light distribution is modeled as an observed de Vaucouleurs $`r^{\frac{1}{4}}`$-law with a half-light radius of $`0.55^{\prime \prime }`$.
### 4.3 Piecewise power-law model
So far we did not enforce any strict parameterization of the radial profile. We only restrict the profile to be of the form of eq. 11; in practice, this is an insignificant restriction. In the following sections we show several ways of modeling the radial profile assuming an isolated lens. None of the models is completely satisfactory.
First we use a minimal model, which assumes that the mass-radius relation is a piecewise power-law, that is, we connect a straight line through two images in the $`\mathrm{log}m`$ vs. $`\mathrm{log}\omega `$ plane. The piecewise values for the power-law slope and axis ratios are given in Table 2.
The Hubble constant $`H_0=100h`$ can be estimated by normalizing the time delay to the observed value $`t_{AC}`$ (Schechter et al. 1997),
$$h=\frac{\tau _{100}}{2t_{AC}}\left[P_{14}+\left(1\frac{1}{\alpha _{14}}\right)S_{14}\right],$$
(33)
where
$$P_{ij}=\left[(x_i^2+y_i^2)(x_j^2+y_j^2)\right]$$
(34)
depends on the image positions $`(x_i,y_i)`$ alone and
$$S_{ij}=2P_{ij}2\left[(x_ix_j)x_s+(y_iy_j)y_s\right]$$
(35)
depends on the source position $`(x_s,y_s)`$ as well. The model yields a Hubble constant $`H_030`$ km/s/Mpc, much lower than determined by other authors (e.g., Impey et al. 1998). The low $`H_0`$ is a result of the high value for the power-law slope $`\alpha _{AC}=1.6`$. Other values are given in Table 2 for various image pairs and observed time delay. The time delay ratio can also be predicted with (cf. eq. 31)
$$r_{ABC}\frac{t_{14}}{t_{32}}=\frac{P_{14}+\left(1\frac{1}{\alpha _{14}}\right)S_{14}}{P_{32}+\left(1\frac{1}{\alpha _{32}}\right)S_{32}}.$$
(36)
Substituting the slopes $`\alpha _{AC}`$ and $`\alpha _{BA}`$ to eq. 31 we find the ratio between images $`r_{ABC}=t_{AC}/t_{BA}0.65`$.
Note that the time delay predictions here are robust and independent of details of the density, since the discontinuity in the density is completely smoothed out in the lensing potential.
### 4.4 Single power-law model
The piecewise-power-law model above necessarily implies a discontinuous density profile. This could be cured if we enforce a single power-law model, that is, we fit a straight line to the four points in the $`\mathrm{log}m`$-$`\mathrm{log}\omega `$ plane. This would give us a power-law slope $`\alpha =1.38`$. Fig. 1 shows our model surface density contours. Similar to the non-parametric models of Saha & Williams (1997) we find a peanut-shaped lens.
We can estimate the goodness of the fit by recomputing the image positions from a given mass model. The general procedure of simulating images of our theoretical lens model is as follows: First combine eq. 26 with the power-law radial profile to get
$$r^2rr_s\mathrm{cos}(\theta \theta _s)=m(\omega )=b_0r_0\left(\frac{\omega }{r_0}\right)^\alpha ,\omega =r\left|1\delta \mathrm{cos}(2\theta 2\theta _p)\right|^\beta .$$
(37)
Then upon substitution of eq. 18 to eliminate $`r`$, we obtain a one-dimensional non-linear equation for the image position angle $`\theta `$, which can be solved easily numerically. As a comparison, one would be dealing with a minimum-finding or a root-finding numerical problem in a two-dimensional plane $`(r,\theta )`$ in the general case without the separation of the angular vs. radial part. For our model of PG1115+080, the image solutions are shown in Fig. 1 as dashed circles, together with the input observed image positions (solid circles). The predicted four images are off by about 60 milli arcsecs, a residual which is inconsistent with the $`3`$ milli arcsecs accuracy of Impey et al. positions, and is marginally consistent with earlier data by Kristian et al. with a quoted error of 50 milli arcsec for the lens galaxy and 5 milli arcsec for the images. The images $`A_1`$ and $`A_2`$ are at two sides of the critical curve (the line of infinite amplification in the source plane), hence are highly amplified with opposite parity. The time delay ratio can be estimated with eq. 31, assuming a constant power-law slope $`\alpha =\alpha _{BC}=1.38`$. This yields $`r_{ABC}1.3`$, close to the Barkana (1997) value of $`1.13\pm 0.18`$. The large difference here is due to the large residual in terms of fitting the images $`A_1`$ and $`A_2`$ with a straight power-law.
We also compute the amplification patterns by taking double derivatives of the time delay surface. The circles in Fig. 1 show the observed (solid circles) and predicted (dashed circles) fluxes of each image and the source (half-closed circle), with the area of each circle in proportion to the flux. The $`B`$ to $`C`$ ratio is well-reproduced and the predictions for images $`A_1`$ and $`A_2`$ are also consistent with observations at about the 0.3 magnitude level.
Finally, if we put a host galaxy around the QSO, the model predicts that the image of the host galaxy will be stretched into an arc. We see a nearly closed ring. This agrees very well with the diffuse ring that Impey et al. discovered in their NICMOS images. The ring maps back to the source plane as a disk with an area of $`0.03`$ square arcsec.
### 4.5 Double power-law model
Alternatively we can fit a smooth, five-parameter lens model with a lens potential
$$\psi (\omega )=c_0\left(\frac{\omega }{a_0}\right)^{\alpha _{in}}\left[1+\left(\frac{\omega }{a_0}\right)^n\right]^{\frac{\alpha _{out}\alpha _{in}}{n}},\omega =r\left|1\delta \mathrm{cos}(2\theta 2\theta _p)\right|^\beta .$$
(38)
The corresponding mass profile is given in Appendix B. This model assumes the lens potential (as well as the lens mass) increases like a double-power-law with an inner slope $`\alpha _{in}`$ and outer slope $`\alpha _{out}`$. The transition is defined by the normalization constant $`c_0`$, the radius $`r=a_0`$ and the sharpness parameter $`n`$; bigger $`n`$ corresponds to sharper transition. When fitting the mass model to the four points in Fig. 3, it turns out that the inner slope $`\alpha _{in}`$ is fully unconstrained. The other four free parameters $`(c_0,a_0,\alpha _{out},n)`$ are determined from fitting the four data points; the procedure is explained in Appendix B. The values of $`(c_0,a_0,\alpha _{out})(0.6,1.1,1.63)`$, nearly independent of the value for the inner slope $`\alpha _{in}`$. Note $`\alpha _{out}\alpha _{AC}`$, consistent with the fact that the mass at $`A_1`$, $`A_2`$ and $`C`$ follows nearly a power-law (cf. Fig. 3). The value of $`n`$ increases from $`15`$ to $`39`$ for $`\alpha _{in}=02`$. That $`n2`$ means that the density changes sharply at the transition. All fits have zero residual and predict nearly identical mass profiles at radii between images $`B`$ and $`C`$. They differ only in the mass profile inside image $`B`$, where we have no direct constraint on the dark matter profile. For all purposes it is sufficient to set $`n=20`$ so that $`\alpha _{in}2`$, in which case the model has a finite core at small radii. From the potential model we can compute the dimensionless surface density contours (cf. Fig. 4). Near the position of the images, the density has a flattening of E2-E3, flatter than the potential, as expected. The potential model can also be substituted in eq. 6 to predict the time delay contours, shown in Fig. 4. The observed images $`C`$, $`A_2`$, $`A_1`$, $`B`$ (the solid circles) are exactly the valley, peak, saddle, valley points of the delay contour, where the theoretical images should lie according to Fermat’s principle. The images $`A_1`$ and $`A_2`$ have nearly the same arrival time. Substituting the lens potential in eq. 28 we can predict the time delay ratio for the double-power-law model. The result, $`r_{ABC}0.65`$, is nearly the same as the piece-wise power-law model. These predicted ratios are in good agreement with the earlier measurement of $`0.7\pm 0.3`$ of Schechter et al. (1997).
We can recompute the image positions from the double-power-law model by solving
$$r^2rr_s\mathrm{cos}(\theta \theta _s)=m(\omega )=c_0\left(\frac{\omega }{a_0}\right)^{\alpha _{in}}\left[1+\left(\frac{\omega }{a_0}\right)^n\right]^{\frac{\alpha _{out}\alpha _{in}}{n}},\omega =r\left|1\delta \mathrm{cos}(2\theta 2\theta _p)\right|^\beta ,$$
(39)
and eq. 18.
Unlike the single-power-law model it predicts five images (cf. Fig. 4). Four predicted images fall exactly on the observed positions. But the model predicts an extra image near image B. This can be identified with the usually highly demagnified fifth image that is sometimes observed in lensed systems. The extra image the model predicts is, however, too bright to be consistent with observations. This is a direct consequence of the fact that $`\kappa `$ at its radius is nearly constant and $`1`$. The amplification can be calculated from $`\kappa `$ as $`(1\kappa )^225`$.
All images lie on the curve defined by eq. 18. The extra image arises because the model density profile (cf. Fig. 5) is non-monotonic near $`1^{\prime \prime }`$ radius, the transition radius $`a_0`$ of the double-power-law potential model. A wiggle in density happens when the transition parameter $`n>2`$, i.e., a sharp transition of the potential.
## 5 Discussion
### 5.1 Power-law slope of PG1115+080 vs. $`H_0`$ and time delay ratio
The previous sections have illustrated the procedure for fitting known quadruple image systems with our models, and have also revealed some puzzling problems with PG1115+080, for example, the low and non-unique $`H_0`$ and the peculiar and non-unique lens density. Fig. 6 shows the values for $`\alpha `$ predicted from either the mass-radius relation (cf. eq. 27) or from the time delay (cf. eq. 29). Two comments are in order: (1) There is a large spread with the power-law slope $`\alpha `$ depending on how it is predicted. To fit the observed time delay within $`2\sigma `$ with a reasonable $`H_0`$ (between 50 and 100 km/s/Mpc), we obtain only a loose constraint $`0.3<\alpha <1.6`$. (2) The rise of the power-law slope $`\alpha `$ with radius in the piecewise-power-law model is somewhat unphysical; realistic models typically have a monotonically steeper density profile (smaller $`\alpha `$) with increasing radius. These results are consistent with the finding of a wiggle near the transition radius ($`1^{\prime \prime }`$) in the surface density of the smooth double-power-law fits. Nevertheless the density is positive everywhere.
A detailed treatment of these problems should include the effect of the neighboring group. Such models are presented elsewhere, since they are beyond the interest of this paper on methods for quadruple systems in general. Nevertheless some further study of the model parameter space is clearly needed. Here it is sufficient to comment on the range of the predicted model parameters, in particular, the power-law slope and the value of $`H_0`$ when we vary the parameter $`\beta `$ or $`\theta _p`$.
Fig. 7 shows the parameter space of the isolated lens model by varying the orientation of the lens principal axis $`\theta _p`$ (cf. eq. 12 and 13). It turns out a one-parameter family of models fits the observations; we effectively vary $`\beta `$. In fact, $`\theta _p61^o+4\beta `$, and $`\mathrm{\Delta }=2\beta \delta \frac{\beta }{2}`$ to a fair accuracy. We find that models with $`\beta <1/3`$ will generally result in negative density models, models with $`\beta 0`$ have a radially-increasing density, and models with $`|\beta |0`$ also produce a nearly constant axisymmetric density, hence implying an unphysically small Hubble constant. Only in the range $`\frac{1}{4}\beta <0`$ do we find physical models with a reasonable flattening. The source positions for different models are shown in Fig. 8. Table 2 compares the predicted parameters of the $`\beta =\frac{1}{4}`$ or $`\theta _p=60^{}`$ model with the $`\beta =\frac{1}{8}`$ or $`\theta _p=60.5^{}`$ model. Fig. 3 shows the predicted radial profile for both models. Both models would have a significant residual if we fit the image positions with a straight power-law, or strange wiggles and extra images if we fit a double-power-law. We find this is generally the case with all isolated lens models without the shear from the neighbouring group.
Fig. 9 shows that there is significant degeneracy with the value for $`H_0`$, depending on the adopted model and the adopted time delay. The most reliable estimate is from the delay between the innermost image B and the outermost image C using
$$h=\frac{\tau _{100}}{2t_{BC}}\left[P_{14}+\left(1\frac{1}{\alpha _{14}}\right)S_{14}P_{32}+\left(1\frac{1}{\alpha _{32}}\right)S_{32}\right],$$
(40)
where we neglected the time delay between the images $`A_1`$ and $`A_2`$. The $`H_0`$ thus predicted scales approximately as $`60(2\alpha )`$ km/s/Mpc (cf. Fig. 9), where $`\alpha =\alpha _{BC}`$ is the average power-law slope. The factor $`2\alpha `$ is the average power-law slope of the model surface density, and is not well-constrained at the radius of the images. Models with $`\beta \frac{1}{4}`$ or $`\theta _p=60^{}`$ predict a steeper density profile, and closer to isothermal than models with $`\beta \frac{1}{8}`$ or $`\theta _p=60.5^{}`$, hence yield a more plausible $`H_0`$ around $`60`$ km/s/Mpc. Models with very shallow profiles ($`\alpha >1`$) are unfavored by the consensus value for $`H_0`$.
In all our isolated lens models, the delay ratio is closer to the Schechter value, while the Barkana value seems to be the widely accepted one. This explains why we see in Fig. 9 a tighter prediction of $`H_0`$ using $`t_{AC}`$ and $`t_{BA}`$ of Schechter values than the Barkana values. Interestingly in the limit that the isothermal model is applicable, $`\alpha =1`$, the time delay ratio can be estimated from the image positions directly (i.e. without computing the source position) with the following equation:
$$r_{ABC}=\frac{P_{14}}{P_{32}}=\frac{r_1^2r_4^2}{r_3^2r_2^2}0.65,\text{for isothermal lens},$$
(41)
where $`r_i`$ is the observed radius of the image $`i`$ from the lens center. The ratio increases when we introduce the shear from the nearby group.
### 5.2 Rotation curve and mass-to-light ratio
Fig. 10 compares the predicted velocity dispersion for the lens with the observed value. The dispersion is predicted with the following formula,
$$\sigma _i^2=V^2\frac{m_i}{2r_i},V^2\frac{D_{os}}{D_{ls}}\frac{(3\times 10^5\mathrm{km}/\mathrm{s})^2}{2\pi \mathrm{radian}}\frac{1\mathrm{radian}}{206265\mathrm{arcsec}},$$
(42)
where the predictions are made from each image $`i=1,2,3,4`$ (strictly speaking the formula is valid only for a singular isothermal lens). Applying this for each value of the lens position angle $`\theta _p`$ we get the full range of possible values for the dispersion. The difference of the dispersion among the four images shows the deviation from an isothermal model. On average the predicted lens galaxy velocity dispersion, 200-300 km/s, suggests that the lens is close to a massive $`L_{}`$ galaxy. The observed dispersion $`\sigma =(281\pm 50)`$ km/s ($`95\%`$ confidence, shaded region) from the Keck spectrum of Tonry (1998). This value is comparable to the theoretical models, but on the high side. A similar problem has been noted in previous models (Schechter et al. 1997). Tonry noted that the discrepancy might be due to a steep radial fall-off of the velocity dispersion or a simple over-estimation of the observed dispersion; the spectrum of the faint lens was made from subtracting two Keck spectra taken with a $`1^{\prime \prime }`$-slit, one passing the lens and the image B, one passing the images $`A_1`$ and $`A_2`$.
Fig. 11 compares the increase of the lens mass from the innermost image $`B`$ to the outermost image $`C`$ vs the increase of the lens light. The mass ratio is predicted using
$$\frac{M_C}{M_B}=\frac{m_4}{m_3}=\left(\frac{r_4}{r_3}\right)^{\alpha _{BC}}.$$
(43)
where the power-law slope $`\alpha _{BC}`$ depends the principal axis $`\theta _p`$. The light ratio
$$\frac{L_C}{L_B}=\frac{_0^{r_4}𝑑r2\pi r\mathrm{exp}\left[7.67\left(\frac{r}{r_e}\right)^{\frac{1}{4}}\right]}{_0^{r_3}𝑑r2\pi r\mathrm{exp}\left[7.67\left(\frac{r}{r_e}\right)^{\frac{1}{4}}\right]},r_e=0.55^{\prime \prime },$$
(44)
where the observed de Vaucouleurs law is used for the projected light. The fact that $`\frac{M_C}{M_B}>\frac{L_C}{L_B}>1`$ implies that the mass grows faster than the light at these radii for all these models, consistent with a large amount of dark matter between 0.8 and 1.5 arcsec.
## 6 Conclusion
Here we itemize our main results.
(i) We found a very general class of lens models that allow for non-elliptical and non-scale-free lenses. These models include previous isothermal elliptical models as special case. We can derive the radial mass distribution of the lens in a non-parametric way. We can study the deviation from the usually assumed straight power-law profile. We also give simple formulas for computing the time delay ratios and for estimating $`H_0`$.
(ii) The models are very easy to compute and can be used for efficient exploration of a large parameter space because the lens equations can be reduced to a set of linear equations.
(iii) We have applied the models to PG1115+080, and have explored a large parameter space of isolated lens models. All models using piece-wise-power-law or double-power-law fit the positions of the images exactly. The models can also reproduce the flux ratios between the images and the stellar velocity dispersion of the lens approximately.
(iv) Our models are consistent with a dark halo up to a radius of three times the half-light-radius of the lens. The enclosed mass increases much faster than the enclosed light as we move radially from the innermost image to the outermost image.
(v) We reconfirm earlier results by (e.g. Schechter et al. 1997) that the principal axis of the lens potential points, within a few degrees, to the external group, and is consistent with the observed value (Iwamuro et al. 2000).
(vi) Our models do not yield a unique prediction of $`H_0`$, because it is sensitive to the lens power-law slope $`\alpha `$ (cf. Fig. 9), which appears to have a large spread, depending on how it is predicted. The power-law slope is also sensitive to small uncertainties of the position angle of the lens $`\theta _p`$. A change of $`\theta _p`$ by a few degrees from the observed value can change $`\alpha `$ from $`0`$ to $`2`$ and $`H_0`$ from 100 km/s/Mpc to 0.
(vii) Our models predict consistently a low time delay ratio $`r_{ABC}`$, which fits only the Schechter value. This and the peculiar oscillation in the predicted density profiles, we believe, are due to the neglect of the external group. This will be dealt with in detail in a follow-up paper.
In summary, our semi-analytical lens models allow us to explore a large range of physical lens density distributions. We find that there is a large systematic uncertainty in the lens models from fitting image positions and time delays of PG1115+080 and isolated lens models are always unsatisfactory for this quadruple system.
The authors thank Tim de Zeeuw for a careful reading of the manuscript and the referee for a constructive report. HSZ thanks Paul Schechter for discussions and encouragement.
## Appendix A Power-law models
Within a certain range of radii, realistic lens mass distributions can be approximated by the following bi-symmetric power-law of the radius $`\omega `$
$$m(\omega )\omega \frac{d\psi }{d\omega }=b_0r_0\left(\frac{\omega }{r_0}\right)^\alpha ,\omega =rf(\theta )=r\left|1\delta \mathrm{cos}(2\theta 2\theta _p)\right|^\beta ,r_0=1^{\prime \prime },0<\alpha <2,$$
(45)
where $`b_0`$ is the characteristic deflection strength at the characteristic length scale $`r_0`$ (one arcsec) of the lens system, and $`\theta _p`$ defines the principle axis of the lens. The corresponding lens potential is related to the enclosed mass $`m(\omega )`$ by
$$\psi (\omega )=\mathrm{const}+\frac{m(\omega )}{\alpha }.$$
(46)
The axis ratio (or its inverse) of the bi-symmetric potential
$$q_\psi \left(\frac{1\delta }{1+\delta }\right)^\beta 1\mathrm{\Delta },\mathrm{\Delta }=2\beta \delta ,$$
(47)
and the axis ratio of the density
$$q_\kappa \left(\frac{1\delta }{1+\delta }\right)^{\frac{1\alpha \beta }{2\alpha }}\left[\frac{1(\frac{4\beta }{\alpha }1)\delta }{1+(\frac{4\beta }{\alpha }1)\delta }\right]^{\frac{1}{2\alpha }}1\mathrm{\Delta }(1+2\alpha ^1),$$
(48)
where the approximations are valid if the flattening is small. The surface density $`\kappa `$ of the model is given by
$$\kappa =\frac{1}{2}br^{\alpha 2}g^{\alpha \beta }(\alpha g^2+(\beta ^2\alpha \beta )g^2+\beta g^{\prime \prime }g),$$
(49)
where $`g1\delta \mathrm{cos}(2\theta 2\theta _p),g^{}2\delta \mathrm{sin}(2\theta 2\theta _p),g^{\prime \prime }4\delta ^2\mathrm{cos}(2\theta 2\theta _p)`$.
## Appendix B Double-power-law models
For the double power-law lens potential
$$\psi (\omega )=c_0\left(\frac{\omega }{a_0}\right)^{\alpha _{in}}\left[1+\left(\frac{\omega }{a_0}\right)^n\right]^{\frac{\alpha _{out}\alpha _{in}}{n}},$$
(50)
the corresponding mass profile is given by
$$m(\omega )=\omega \frac{d\psi (\omega )}{d\omega }=\alpha (\omega )\psi (\omega ),$$
(51)
where
$$\alpha (\omega )=\frac{d\mathrm{log}\psi (\omega )}{d\mathrm{log}\omega }=\frac{\alpha _{in}+\alpha _{out}\left(\frac{\omega }{a_0}\right)^n}{1+\left(\frac{\omega }{a_0}\right)^n}.$$
(52)
Generally speaking, a double-power-law fit (cf. eq. 38) involves searching for solutions of the following equations in a five-parameter space $`(n,\xi _0,\alpha _{out},\gamma _0,a_0)`$:
$$\mathrm{log}m_i=\left[\xi _0+(1+\gamma _0)\alpha _{out}\mathrm{log}\omega _i\right]\left(1+\frac{\gamma _0\alpha _{out}}{n}\right)\mathrm{log}\left[1+\left(\frac{\omega }{a_0}\right)^n\right]+\mathrm{log}\left[1+\gamma _0+\left(\frac{\omega }{a_0}\right)^n\right],i=1,2,3,4,$$
(53)
where
$$\gamma _0\frac{\alpha _{in}}{\alpha _{out}}1,\xi _0\mathrm{log}c_0+\mathrm{log}\alpha _{out}(1+\gamma _0)\alpha _{out}\mathrm{log}a_0,$$
(54)
and $`(\omega _i,m_i)`$ for $`i=1,2,3,4`$ are the four data points in the mass-radius diagram (cf. Fig. 3). Note that the equations 53 are linear in $`(\xi _0,\alpha _{out})`$, so these two variables can be eliminated by Gaussian substitution. If we fix $`n`$, say $`n=20`$, then we are left with two equations for two variables $`(\gamma _0,a_0)`$, which can be solved with a two-dimensional iterative root finding routine. The solution converges rapidly from an initial guess of $`\gamma _0=0`$ and $`a_0=1^{\prime \prime }`$.
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# Phonon-mediated drag at 𝜈=1/2: A test of the Chern-Simons composite fermion theory
## I Introduction
Comparisons between measured transport properties and theoretical predictions have been fruitful in providing important information about elementary excitations in electronic systems. A recent and rather unusual example may be found in studies of the frictional drag between separately contacted, nearby, two-dimensional electron layers. Because it is a measure of the interlayer scattering rate, drag resistance is sensitive to Coulomb and phonon-mediated interactions between the layers and is related to correlated density fluctuations in the coupled bilayer system. Since these depend crucially on underlying electronic degrees of freedom, frictional drag can serve as an important test of specific theoretical predictions concerning the fluctuation spectrum of a given electronic system. At zero field, these studies have verified many qualitative features predicted by the random phase approximation (RPA) theory of correlations in moderate density two-dimensional electron systems. On the other hand, they have also demonstrated the importance, at a quantitative level, of various corrections. These studies have a greater potential importance in the fractional quantum Hall regime, where simple perturbative RPA-like theories fail and our understanding of correlations is much less complete. In the fractional quantum Hall regime, all electrons share the same quantized kinetic-energy and low-energy properties are entirely determined by electron-electron interactions. Electrons in the fractional quantum Hall regime are in the limit of extreme correlations. It is the failure of perturbation theory which admits all the peculiarities and surprises which have been discovered in this field over the past two decades.
One example of a peculiar strong-correlation consequence is presented by the Fermi-liquid-like properties of the compressible states which occur when the lowest Landau level is half-filled. Experimental evidence then supports the existence of gapless fermionic excitations which form a Fermi surface. This is particularly surprising because the kinetic energy of the electrons in the lowest Landau level is quenched by the large magnetic fields, so this Fermi surface must result purely from interaction effects. On the theoretical front, the Chern-Simons composite fermion theory has been spectacularly successful in explaining various experimental findings at the phenomenological level . It starts from a singular gauge transformation which attaches two flux quanta to each electron. These new ‘composite particles’, like electrons, obey Fermi statistics and are therefore called composite fermions. At $`\nu =1/2`$, the mean-field flux-density from the attached flux quanta exactly cancels that coming from the external magnetic field. Neglecting fluctuations, composite fermions at $`\nu =1/2`$ therefore experience zero magnetic field and have a Fermi surface. In the random phase approximation (RPA), the density-density correlation function for composite fermions is equal to the density-density correlation function for completely spin-polarized electrons at zero magnetic field. However, the density-density correlation function of the physical electrons contains an additional local field term associated with the attached flux quanta. We will refer to this additional term as the magnetic local field. In RPA theory, the magnetic-field dependence of the correlation function is completely embedded in this quantity.
Chern-Simons theory has the advantage of being technically and conceptually simple. It suffers, however, from various problems, most associated with the composite fermion mass scale $`m^{}`$. On physical grounds, it is known that $`m^{}`$ is determined by the interaction energy scale in the lowest Landau level. However, the theory does not have a proper lowest Landau level projection so that $`m^{}`$ has to be fixed phenomenologically. Unfortunately, this naive approach leads to a violation of the f-sum rule and Kohn’s theorem in the sense that certain response functions should contain the bare electron mass $`m_b`$. A scheme to correct for this deficiency was developed by Simon and Halperin who introduced the Modified Random Phase Approximation (MRPA). The MRPA includes a Fermi Liquid like correction, and allows for an effective mass $`m^{}`$ different from the band mass $`m_b`$ while still satisfying the f-sum rule and Kohn’s Theorem. Even though this prescription turns out to be surprisingly successful, a more microscopic understanding of the role of the lowest Landau level projection in the theory is desirable.
Recently several attempts have been made to construct a truly lowest Landau level theory of composite fermions. Inspired by $`\nu =1/2`$ variational wavefunction, it has been realized that the position of an electron and a nearby zero of the wavefunction (called a ‘vortex’) are displaced by $`k_il_B^2`$ for each electron labeled by $`i=1,\mathrm{},N`$. The lowest Landau level constraint and strong correlations tend to bind vortices to electrons. The composite object that consists of an electron and two vortices satisfies fermionic statistics and, as a result, the $`k_i`$’s have to be chosen differently for each composite object. The ground state of these ‘neutral composite fermions’ is given by a Fermi sea in the space of $`k_i`$. One can also show that $`k_i`$ acts as the guiding-center-translation generator in the lowest Landau level. Within a Hartree-Fock theory of these composite objects, the effective kinetic energy has the form $`k_i^2/2m^{}`$ with an effective mass $`m^{}`$ given correctly by the interaction energy scale . A self-consistent theory of neutral dipolar composite fermions has been constructed and it has been shown that, in the long-wavelength low-energy limit, this theory (taking the bare electron mass to infinity to reflect the absence of a kinetic energy term in the lowest Landau level Hamiltonian) gives the same physical response functions as those of the Chern-Simons composite fermion theory . Coulomb mediated frictional drag between two $`\nu =1/2`$ systems depends mostly on the long wavelength and low energy limits of density-density correlation in each layer. As a result, the two theories mentioned above should make identical predictions for Coulomb drag. This quantity has been measured and the observations are consistent at relatively high temperatures with the theoretically predicted $`T^{4/3}`$ power law. The data does, however, exhibit a low temperature anomaly that has not yet been fully understood.
When two 2D electron systems are widely separated, the contribution to drag from Coulomb scattering between the layers is suppressed and phonon mediated interaction dominates. Unlike the Coulomb drag case, phonon mediated drag is sensitive to density-density correlations in each layer at relatively high frequencies and at wavevectors comparable to the inverse interparticle separation. Therefore studies of phonon mediated drag should provide useful information about excitations at finite frequencies and at large wavevectors, leading to a test of the underlying theory at many different energy and length scales, not just at the long wavelength and low energy limits. Since the equivalence between Chern-Simons theory and the neutral dipolar composite fermion theories has been established only in the long wavelength and low energy limits, studies of phonon drag may reveal important differences between these two theories. As the first step to investigate the validity of these theories, in this paper we report on a comparison between recent measurements of frictional drag between widely separated two-dimensional electron layers at $`\nu =1/2`$ and Chern-Simons composite fermion theory predictions. We find some qualitative differences between observations and trends predicted by theory and speculate on their implications for the theory of the half-filled Landau level.
In the following section we review the theory of frictional drag and discuss some of the features that are important in analyzing drag at $`\nu =1/2`$. The theoretical framework we use to describe frictional drag is based on earlier $`B=0`$ work involving two of us. While there are some issues in the theoretical literature on frictional drag which are not completely settled, these do not impact on the conclusions we draw here and we do not discuss them further. In Sec. III we present numerical results for drag, calculated from (M)RPA composite-fermion theory, and compare these with recent experimental results obtained by Zelakiewicz et al.. Our conclusions are presented in Sec. IV.
## II Frictional drag
When two or more electronic systems are placed in close proximity, their transport properties are interdependent. In particular, currents in one subsystem can induce (or ’drag’) currents in other subsystems. The first drag experiments were on coupled two-dimensional and three-dimensional electron system. Drag between two-dimensional layers has been measured in electron-electron, electron-hole, and recently in hole-hole systems. In this paper we consider frictional drag between two two-dimensional electron gases embedded in a GaAs-AlGaAs double quantum well system. The quantum wells widths $`L`$ are $`200`$ Å and are separated by a center-to-center distance $`d`$ which is large enough to ensure that interlayer tunneling can be neglected ($`d300`$ Å). We imagine a current density $`J_1`$ being drawn in the first layer while the second layer is an open circuit. To counter balance the stochastic drag force between the layers, an electric field $`E_2`$ will build up in the second layer. The transresistivity, $`\rho _{21}`$, is defined as
$$\rho _{21}E_2/J_1,$$
(1)
$`\rho _{21}`$ depends on the strength of interlayer interactions and on the phase space available for interlayer scattering events.
Theoretical expressions for $`\rho _{21}`$ can be derived by conventional linear response theory at various levels of sophistication. Under conditions which are met in current experiments, the transresistivity is given by
$`\rho _{21}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{8\pi ^2e^2n_1n_2k_\mathrm{B}T}}{\displaystyle _0^{\mathrm{}}}𝑑qq^3`$ (3)
$`{\displaystyle _0^{\mathrm{}}}𝑑\omega \left|{\displaystyle \frac{W_{21}(q,\omega )}{\epsilon (q,\omega )}}\right|^2{\displaystyle \frac{\mathrm{Im}\mathrm{\Pi }_1(q,\omega )\mathrm{Im}\mathrm{\Pi }_2(q,\omega )}{\mathrm{sinh}^2(\mathrm{}\omega /2k_\mathrm{B}T)}}.`$
Here $`n_i`$ and $`\mathrm{\Pi }_i(q,\omega )`$ are the two-dimensional electron density and the polarization function, respectively, of layer $`i`$. In Eq. (3) $`W_{ij}(q,\omega )`$ is the reciprocal space interaction between layers $`i`$ and $`j`$, and $`\epsilon (q,\omega )`$ is the interlayer screening function, given by
$$\epsilon =[1+W_{11}\mathrm{\Pi }_1][1+W_{22}\mathrm{\Pi }_2]W_{21}^2\mathrm{\Pi }_1\mathrm{\Pi }_2.$$
(4)
Given the interlayer interaction, this theory of frictional drag depends only on the polarization functions of the 2D electron layers.
Both interlayer Coulomb interactions and phonon-mediated interactions contribute to interlayer friction:
$$W_{ij}(q,\omega )=U_{ij}(q)+𝒟_{ij}(q,\omega ).$$
(5)
Taking account of the finite quantum well widths, the Coulomb contribution is
$$U_{ij}(q)=\frac{e^2}{2\kappa q}B_{ij}(qd,qL/2)$$
(6)
where $`\kappa `$ is the dielectric constant of the semiconductor and $`B_{ij}`$ is a width-dependent form factor. The phonon mediated interaction is given by
$$𝒟_{ij}(q,\omega )=\underset{\lambda }{}\frac{dQ_z}{2\pi \mathrm{}}|M_\lambda (𝐐)|^2F_i(Q_z)F_j(Q_z)D_\lambda (𝐐,\omega )$$
(7)
where $`F_i(Q_z)=𝑑z|\phi _i(z)|^2e^{iQ_Zz}`$ is the Fourier transform of the electron density in the direction perpendicular to the layers ($`\phi _i`$ is the subband wavefunction in layer $`i`$), and $`D_\lambda (𝐐,\omega )`$ is the phonon propagator. We use the notation $`𝐐=(𝐪,Q_z)`$. The index $`\lambda `$ denotes longitudinal and transverse phonon modes, $`\lambda =l,t`$. Both acoustic and optical phonons can contribute to drag, but in GaAs/AlGaAs systems acoustic phonons dominate under the conditions studied in this paper. The electron-acoustic phonon coupling matrix elements are given by
$$|M_l(𝐐)|^2=\frac{\mathrm{}Q}{2\varrho c_l}\left[D^2+(eh_{14})^2\frac{9q^4Q_z^2}{2Q^8}\right]$$
(8)
$$|M_t(𝐐)|^2=\frac{\mathrm{}}{2\varrho c_t}(eh_{14})^2\frac{8q^2Q_z^4+q^6}{4Q^7}$$
(9)
where the $`c_\lambda `$ are the phonon velocities for longitudinal and transverse phonons, $`\varrho `$ is the mass density of GaAs/AlGaAs, $`D`$ is the deformation potential, and $`eh_{14}`$ is the piezoelectric constant. In our numerical evaluations we use $`c_l=5140`$ m/s, $`c_t=3040`$ m/s, $`\varrho =5300`$ kg/m<sup>3</sup>, $`D=13.0`$ eV, and $`eh_{14}=1.2\times 10^9`$ eV/m.
Because their contributions come dominantly from different regions of wavevector and frequency, it is possible, both experimentally and theoretically, to separate Coulomb and phonon-mediated contributions to drag. Coulomb drag comes from $`qd^1`$, whereas phonon-mediated drag is dominated by contributions from $`q2k_F`$ and frequencies at the peak in $`𝒟_{ij}`$ near $`\omega =c_\lambda q`$. Work at zero magnetic field has established that Coulomb drag dominates at small interlayer separations($`d300`$ Å) but falls off rapidly with increasing $`d`$. Phonon mediated drag, on the other hand, depends only weakly on $`d`$ and has been observed for layers separated by distances approaching $`1\mu `$m.
### A Frictional drag at $`B=0`$
Frictional drag in the absence of a field is now well understood and the extensive literature has recently been reviewed. Eq. (3), combined with the Random Phase Approximation for the 2D electron system polarization function, successfully predicts a low-temperature drag $`\rho _{21}T^2`$, a substantial enhancement of $`\rho _{21}`$ due to the excitation of plasmon modes at around half the Fermi temperature, and a crossover from Coulomb drag to phonon-mediated drag for layers separated by more than $`50\mathrm{n}\mathrm{m}`$. Phonon-mediated drag is readily identified by its signature dependences on the ratio of electron densities in the two layers and on temperature. For phonon-mediated drag $`\rho _{21}T^6`$ at low temperatures ($`T0.1`$ K) and varies approximately linearly at higher temperatures ($`T5`$ K). As a consequence the scaled transresistivity $`\rho _{21}/T^2`$, which is more nearly constant in the Coulomb case, will have a well defined peak. The origin of this peak in $`\rho _{21}/T^2`$ is easily understood by looking at Fig. 1 which shows the phonon resonance and the particle-hole continuum in phase space. The integrand in Eq. (3) is cut off exponentially with $`\omega `$ at $`\mathrm{}\omega k_\mathrm{B}T`$. The full phonon resonance is thus exploited when $`T=T_{\mathrm{peak}}\mathrm{}c_\lambda 2k_F/k_\mathrm{B}`$, and a further increase in the temperature will only lead to a relatively smaller increase in $`\rho _{21}`$.
The occurrence of a peak depends only on the phonon resonance and the existence of a Fermi surface, which leads to a sharp edge of the particle-hole continuum as shown in Fig. 1. The precise location of the peak, however, will depend on detailed properties of the polarization function at large $`q`$ as well as on many other parameters of the experiment. At $`B=0`$, good agreement with experiment for the overall shape of scaled transresistivity curves and for the location of their peaks corroborates the theory. Disagreements in magnitude can in part be attributed to uncertainties in the electron-phonon interaction parameters ($`\rho _{21}`$ varies approximately as the fourth power of the deformation potential constant $`D`$ and this parameter is not accurately known). Corrections to the RPA polarization function at large wavevector may also play a role. Moreover, the magnitude of $`\rho _{21}`$ depends on the phonon mean free path $`\mathrm{}_{\mathrm{ph}}`$ which appears in the phonon propagator. In Ref. it was shown that for $`\mathrm{}_{\mathrm{ph}}`$ larger than a critical value $`\mathrm{}_{\mathrm{crit}}`$, the transresistivity is enhanced by to a collective mode (i.e. a zero in the real part of the screening function near $`q=2k_F`$). At zero magnetic field $`\mathrm{}_{\mathrm{crit}}200\mu `$m for $`n=1.5\times 10^{15}\mathrm{m}^2`$.
The view taken in this paper is that frictional drag is sufficiently well understood that it can be used to probe, at least at a qualitative level, the wavevector and frequency dependent polarization functions of the coupled 2D electron layers. Just such a probe is urgently needed to test composite fermion theory of the half-filled Landau level.
### B Coulomb drag at $`\nu =1/2`$
Sakhi,, Ussishkin and Stern, and Kim and Millis, using Chern-Simons RPA theory, independently predicted that at $`\nu =1/2`$ the low-temperature transresistivity $`\rho _{21}`$ would be greatly enhanced compared to its $`B=0`$ value. In these theories $`\rho _{21}(T/d)^{4/3}`$ for $`T0`$. This property follows from the form of the Chern-Simons RPA polarization function in the limits $`qk_F`$ and $`\omega v_Fq`$ (see Appendix A). Subsequent experiments by Lilly et al., did demonstrate increased drag which vanished less rapidly than $`T^2`$ with declining temperature. The experimental results are consistent with a $`T^{4/3}`$ law, except below $`0.5\mathrm{K}`$ where the temperature dependence slows further and the drag value becomes sample specific.
In Chern-Simons RPA theory, the Coulomb drag enhancement is due to the diffusive character of composite-fermion density fluctuations at long wavelength and low frequency. In Fig. 2 we plot the Coulomb contribution to $`\rho _{21}`$ (evaluated by numerical integration of Eq. (3)) as a function of $`d`$ for different temperatures. The plot shows that for $`T0.1`$ K, the transresistivity falls of approximately as $`d^2`$, significantly slower than the $`d^4`$ behavior at $`B=0`$, but still fast enough to guarantee that phonon-mediated drag will dominate for widely separated layers. The plots in Fig. 2 were calculated using the Chern-Simons RPA with an effective mass $`m^{}=m_\mathrm{e}`$. In Fig. 3 we illustrate the dependence of Coulomb drag on the choice of the effective mass, and the difference between RPA and MRPA predictions. Phonon mediated drag which is dominated by short wavelengths is not enhanced by the diffusion-like pole in the polarization function. It will, however, as we discuss below, be enhanced over the corresponding zero-field result. A numerical comparison shows that, generally, phonon-mediated drag will dominate over Coulomb drag for $`d200`$ nm.
## III Phonon mediated drag at $`\nu =1/2`$
Two qualitative aspects of the comparison between phonon-mediated drag at $`\nu =1/2`$ and at $`B=0`$ can be anticipated in advance of any detailed calculation. Firstly, since the polarization function is proportional to the effective mass, one would expect the phonon-mediated transresistivity to be larger at $`\nu =1/2`$ by a factor of approximately $`(m^{}/m_b)^2100`$. Notice that this enhancement is weaker than for Coulomb drag which is enhanced by a factor of approximately 1000 (independent of mass for $`T0`$). Secondly, the value of $`T_{\mathrm{peak}}`$ should be larger than at $`B=0`$ for several reasons. The Fermi wavevector is increased by a factor of $`\sqrt{2}`$ if, as we assume here, the electron system is completely spin polarized. The smaller Fermi energy and Fermi velocity, which go along with the larger effective mass, should also increase $`T_{\mathrm{peak}}`$. For larger $`m^{}`$, the ratio of the phonon velocity to the Fermi velocity approaches unity, and the most important drag contributions will then come from energies near the Fermi energy. Correspondingly, within the particle-hole continuum, the largest possible momentum transfer (which is given by $`q=2k_F(1+c_\lambda /v_F)`$) is larger at $`\nu =1/2`$ than for $`B=0`$ ($`c_\lambda v_F`$ for $`B=0`$). Furthermore, since $`T_{\mathrm{peak}}`$ is close to the Fermi temperature $`T_F`$, it is necessary to evaluate the polarization functions at finite temperature. (At finite $`T`$ the Fermi gas polarization function can no longer be evaluated analytically. See Appendix A for a detailed discussion of its numerical evaluation.) When evaluated at finite temperature, the polarization function has finite contributions from outside the $`T=0`$ particle-hole continuum (see Fig. 1). This additional effect suggests that a further increase of $`T_{\mathrm{peak}}(\nu =1/2)/T_{\mathrm{peak}}(B=0)`$ which should therefore be larger than $`\sqrt{2}(1+c_\lambda /v_F)`$. Indeed our numerical calculations confirm this expectation which is not, however, in agreement with experiment.
Zelakiewicz et al. have measured the transresistivity in samples with $`d=2600`$ Å and $`d=5200`$ Å where they find no evidence of a contribution from the Coulomb interaction. The transresistivity was indeed found to have the same qualitative features as at $`B=0`$, except that the transresistivity is enhanced by a factor of $`200`$. The scaled transresistivity $`\rho _{21}/T^2`$ does have a peak as a function of temperature whose location scales with electron density as $`\sqrt{n}`$, just as in the $`B=0`$ case. However, the value of $`T_{\mathrm{peak}}`$ was found to be slightly smaller at $`\nu =1/2`$ than at $`B=0`$. The detailed calculations we present below thus tend to indicate that Chern-Simons composite fermion theory overestimates electron polarizability at wavevectors $`2k_F`$, as anticipated by the authors of Ref. . Our calculations are similar to ones performed earlier at $`B=0`$; the most troublesome complication is the requirement that the free-particle polarization function be evaluated at finite temperatures. This aspect of the calculation is explained in detail in Appendix A.
Figs. 4, 5, and 6 show numerical evaluations of $`\rho _{21}/T^2`$ as a function of $`T`$ for $`m^{}=0.067m_e`$, $`m^{}=0.4m_e`$ and $`m^{}=m_e`$ respectively. Calculations at intermediate values of $`m^{}`$ interpolate smoothly between the trends illustrated in these three figures. The striking differences we find in the temperature and density dependences for the different values of $`m^{}`$ are due primarily to the different values of the Fermi temperature compared to the energy of an acoustic phonon with a wavelength comparable to the distance between electrons. We concentrate first on the property which is most prominent in the experimental data, as conventionally presented, $`T_{\mathrm{peak}}`$. We discuss its dependence on density and the phenomenological mass parameter $`m^{}`$. At zero magnetic field $`T_{\mathrm{peak}}=2.6\mathrm{K}`$ in a calculation which uses the same parameters as in Fig. 4. Fig. 4 shows results obtained when the band mass is used for $`m^{}`$ so that differences from the $`B=0`$ case are due only to spin-polarization and the introduction of the magnetic local field. The local field increases the relative importance of small angle scattering and decreases $`T_{\mathrm{peak}}`$ below the naively expected value $`\sqrt{2}\times 2.6`$ K. This tendency of the local field is closely connected to the modification which changes the temperature dependence from $`T^2`$ to $`T^{4/3}`$ in the Coulomb drag case. Even with this effect, however, $`T_{\mathrm{peak}}`$ is safely above its zero field value, in disagreement with the experimental findings of Ref. . For this mass, $`T_{\mathrm{peak}}`$ increases with density approximately as $`\sqrt{n}`$ as in the $`B=0`$ case. The origin of this behavior is simply that the highest momentum transfer within the particle-hole continuum is proportional to $`\sqrt{n}`$ and that the Fermi temperature for these values of $`m^{}`$ is sufficiently high compared to $`T_{\mathrm{peak}}`$ that Im$`\mathrm{\Pi }(q,\omega )`$ has negligible weight outside the particle-hole continuum.
In Figs. 5 and 6 we see that $`T_{\mathrm{peak}}`$ is still larger at the larger values of $`m^{}`$ required to fit other data in the quantum Hall regime. This is in accord with the considerations outlined above. With increasing $`m^{}`$ the Fermi temperature decreases and momentum transfers greater than $`2k_F`$ receive more weight. For higher values of the mass, it is interesting to see that the value of $`T_{\mathrm{peak}}`$ has a more complex density dependence, reflecting a competition between two effects. As the density increases, the maximum possible momentum transfer increases, which tends to increase $`T_{\mathrm{peak}}`$. At the same time, however, the Fermi temperature increases which leads to a smaller contribution from outside the particle-hole continuum. This tends to decrease $`T_{\mathrm{peak}}`$. Because of this competition, the density-dependence of $`T_{\mathrm{peak}}(n)`$ changes from increasing at $`m^{}=0.067m_e`$, to decreasing at $`m^{}=0.4m_e`$, to approximately constant for $`m^{}>0.4m_e`$. As for the magnitude of the transresistivity, we see that for a phonon mean free path $`\mathrm{}_{\mathrm{ph}}=100\mu \mathrm{m}`$, we would have to choose an effective mass in the neighborhood of $`m^{}=0.3m_\mathrm{e}`$ in order to match the experimental result that $`\rho _{21}/T_{\mathrm{peak}}^2200\mathrm{m}\mathrm{\Omega }/\mathrm{K}^2`$. (For this choice of parameters , the numerical results need to be multiplied by a factor of 5.7 to match the experimental results at zero magnetic field).
Phonon mediated drag might not, however, be a good way of determining the effective mass of Composite Fermions. The critical phonon mean free path $`\mathrm{}_{\mathrm{crit}}`$ itself depends on $`m^{}`$ ($`\mathrm{}_{\mathrm{crit}}1/(m^{})^2`$) and will thus be different for $`B=0`$ and $`\nu =1/2`$. In Fig. 7 we plot the transresistivity as a function of $`\mathrm{}_{\mathrm{ph}}`$ for different values of $`m^{}`$. At $`\mathrm{}_{\mathrm{ph}}\mathrm{}_{\mathrm{crit}}`$ a long-lived collective mode develops, leading to an enhancement of $`\rho _{21}`$ before it saturates at $`\mathrm{}_{\mathrm{ph}}\mathrm{}_{\mathrm{crit}}`$. The qualitative distance dependence is rather similar to the zero magnetic field case (see Ref. ): For $`\mathrm{}_{\mathrm{ph}}\mathrm{}_{\mathrm{crit}}`$, the transresistivity depends logarithmically on $`d/\mathrm{}_{\mathrm{ph}}`$, i.e. $`\rho _{21}\mathrm{ln}(Ad/\mathrm{}_{\mathrm{ph}})`$ for $`d\mathrm{}_{\mathrm{ph}}/A`$ where $`A`$ is a constant of order $`2k_FL`$. For $`d\mathrm{}_{\mathrm{ph}}/A`$ the transresistivity falls off exponentially, $`\rho _{21}\mathrm{exp}(Ad/\mathrm{}_{\mathrm{ph}})/d`$. If, on the other hand, the phonon mean free path is larger than the critical value $`\mathrm{}_{\mathrm{crit}}`$, the distance dependence is more complicated.
## IV Discussion
In this paper we describe the dependence of phonon mediated drag at $`\nu =1/2`$ on system parameters and composite fermion mass predicted by the Chern-Simons MRPA theory. We have focused on the characteristic temperatures $`T_{\mathrm{peak}}`$ at which the scaled drag resistivity $`\rho _{21}(T)/T^2`$ reaches its peak. Crudely we expect $`k_BT_{\mathrm{peak}}`$ to equal the phonon energy at the largest wavevector for which the two-dimensional electron systems have substantial charge fluctuations. Because of the sharp particle-hole continuum of Chern-Simons composite fermion theory, the maximum momentum transfer is close to $`2k_F^{\mathrm{cf}}`$, where $`k_F^{\mathrm{cf}}`$ is the Fermi wavevector of the composite fermion Fermi sea. These considerations lead to $`T_{\mathrm{peak}}\mathrm{}c_\lambda 2k_F^{\mathrm{cf}}/k_B`$. Since $`k_F^{cf}=\sqrt{2}k_F`$ for completely spin polarized composite fermions, $`T_{\mathrm{peak}}`$ at $`\nu =1/2`$ should be larger than $`T_{\mathrm{peak}}`$ at $`B=0`$, in disagreement with experiment. Our detailed calculations confirm that, at a qualitative level, this is indeed the prediction of Chern-Simons composite fermion theory.
At a quantitative level, the change of Fermi radius is not the only difference which arises in comparing phonon-mediated drag at $`B=0`$ and at $`\nu =1/2`$. One source of complication is the strong dependence of the drag resistivity on the effective mass $`m^{}`$ of the composite fermions. The large effective mass of composite fermions leads to a substantially smaller Fermi energy and Fermi velocity. This in turn leads to larger maximum momentum transfer, $`q=2k_F^{\mathrm{cf}}(1+c_\lambda /v_F)`$, since $`c_\lambda `$ is then comparable to $`v_F`$. In addition, contributions to drag from outside the $`T=0`$ particle-hole continuum increase. On the other hand, the magnetic local field correction increases the relative importance of small angle scattering and tends to decrease $`T_{\mathrm{peak}}`$. Because of these mutually competing effects, the position of the maximum is sometimes larger than value expected on the basis of naive considerations and quoted above. For example, when $`m^{}`$ is comparable to $`m_b`$, the Fermi temperature effect is relatively small and the magnetic local field effect leads to a $`T_{\mathrm{peak}}`$ smaller than the naive value. Nevertheless, the position of $`T_{\mathrm{peak}}`$ at $`\nu =1/2`$ is always well above that of $`B=0`$ case.
We now present a detailed summary of the comparison between our theoretical results and experiment . The following aspects of the experimental data are consistent with theoretical results.
1. The magnitude of the phonon mediated drag at $`\nu =1/2`$ is about 200 times larger than at $`B=0`$. In the theory, the factor of a few hundred comes from the enhanced density of states, $`m^{}/2\pi \mathrm{}^2`$, of composite fermions compared to the bare density of states, $`m_b/2\pi \mathrm{}^2`$, of electrons at $`B=0`$. $`m^{}/m_b10`$ according to numerical and theoretical estimations . The presence of two layers leads to $`(m^{}/m_b)^2100`$ fold enhancement in the drag rate.
2. When the phonon mean free path of $`\mathrm{}_{ph}=100\mu m`$ is chosen consistently for $`B=0`$ and $`\nu =1/2`$, the experimental data can be fit by choosing $`m^{}`$ around $`0.3m_e`$ which is consistent with previous numerical and theoretical estimations .
3. There is a well defined maximum of $`\rho _{21}/T^2`$, suggesting that $`\nu =1/2`$ systems have a fairly sharp wavevector cutoff for low-energy charge fluctuation, as predicted by Chern-Simons composite fermion theory.
On the other hand, the following experimental results do not agree with theoretical predictions.
1. Experimentally, $`T_{\mathrm{peak}}`$ at $`\nu =1/2`$ is always smaller than that of $`B=0`$. Theoretically it is always larger, even when unrealistically small values are used for the composite fermion effective mass.
2. Experimentally, $`T_{\mathrm{peak}}`$ is proportional to $`\sqrt{n}`$. Theoretically, this simple behavior is found only when unrealistically small values are used for the composite fermion effective mass.
We now discuss some possible sources of these discrepancies.
Incomplete spin polarization: In our calculations, we assumed complete spin polarization. If the spins are only partially polarized, the Fermi sea of the majority spin would be smaller than in the fully polarized case, leading to a smaller maximum momentum transfer and a smaller $`T_{\mathrm{peak}}`$. Recently the spin polarization of $`\nu =1/2`$ state was measured by NMR . In the experiment by S. Melinte et al. , the electron density, $`n=1.4\times 10^{11}cm^2`$, of their sample M242 is similar to the density, $`n=1.5\times 10^{11}cm^2`$, of the sample studied in the drag experiment of S. Zelakiewicz et al. . In the NMR experiment , $`\nu =1/2`$ is reached when the external magnetic field is $`B=11.4`$ T, compared to the field strength $`B=12.82`$ T for $`\nu =1/2`$ in the drag experiment . This NMR experiment can therefore be used to get a reasonable estimate of the $`\nu =1/2`$ spin polarization in the drag experiment. We conclude that the spin polarization in the $`\nu =1/2`$ drag experiment at $`12`$ K was about $`70\%`$ of the full polarization. This would imply that majority-spin Fermi wavevector was $`\sqrt{1.7}k_F1.3k_F`$ and that of the minority spin Fermi wavevector is $`\sqrt{0.3}k_F0.55k_F`$. This reduction in the majority spin-wavevector ($`1.3k_F`$ compared to $`\sqrt{2}k_F`$) is not sufficient to explain the low $`T_{\mathrm{peak}}`$ values seen in the experiments. In our theory, we would not expect a large contribution to drag from the minority spins, and in any event the Fermi radius $`0.55k_F`$ is too small for it to be associated with the observed value of $`T_{\mathrm{peak}}`$ slightly below that of $`B=0`$ case. We conclude that incomplete spin polarization alone cannot explain the discrepancy between theory and experiment.
Breakdown of (M)RPA at large wavevectors: This experimental result may signal that the (M)RPA is not adequate to describe the large wavevector response of a $`\nu =1/2`$ system. It is generally accepted that Ward identities related to conservation laws, limit beyond (M)RPA contributions to density-density and current-current response functions in the long wavevectors and low energy limit. However, it has been also shown that response functions at large wavevectors, in particular near $`2k_F`$, may be modified at low energies by singular vertex corrections . Even though these singular corrections appear at low energies at $`2k_F`$, it is certainly possible that these singular corrections persist to higher energy scales comparable to $`T_{\mathrm{peak}}`$ or $`\mathrm{}c_\lambda 2k_F^{\mathrm{cf}}/k_B`$. If this is true, it may be necessary to get beyond the (M)RPA to obtain an accurate description of phonon-mediated drag.
Dipolar Composite Fermions: As discussed in the introduction, there exist at present two descriptions of $`\nu =1/2`$ composite fermions. It has been established that these two descriptions are equivalent in the low energy and long wavelength limits. However, these two descriptions may lead to quite different predictions at large wavevectors because dipolar composite fermions have a finite size which can be comparable to $`k_F^1`$. Indeed, it has been observed that the equivalence of two approaches may break down at higher energies even in the long wavelength limit . The singular behavior of the response functions near $`2k_F`$ mentioned above in the Chern-Simons theory approach may be a signature of the breakdown of the theory at large wavevector scales and a proper description of the system at large wavevectors may require the fully lowest Landau level dipolar composite-fermion theory. Calculations of the relevant response functions and of the drag resistivity in the dipolar composite fermion approach are in progress .
In conclusion, the Chern-Simons theory of composite fermions overestimates charge fluctuations at the Fermi wavelength scale and a proper description of the system at large wavevectors may have to take into account the extended nature of dipolar composite fermions.
## V Acknowledgements
This work was initiated at the Institute for Theoretical Physics in University of California at Santa Barbara (NFS grant PHY9407194), and was further supported by the NSF under grants DMR-9714055 (M.C.B and A.H.M) and DMR-9983783 (Y.B.K.), by the Danish Research Academy (M.C.B), and by the A. P. Sloan Foundation (Y.B.K.). The authors are grateful T. J. Gramila, B. Y.-K. Hu, and S. Zelakiewicz, for informative and stimulating interactions. Y.B.K. would like to thank also Isaac Newton Institute at the University of Cambridge for the hospitality.
## A Polarization functions
In this appendix we summarize the techniques we use to evaluate the free-fermion polarization functions which appear in Eq. (3) at finite temperatures and, for completeness, present explicit expressions for the (Modified) Random Phase Approximation polarization functions.
The electron polarization function $`\mathrm{\Pi }^\mathrm{e}(q,\omega )`$ expresses the linear response of an electron system to the total electromagnetic field. For a two-dimensional system (in the $`x`$-$`y`$ plane) and a magnetic induction restricted to the $`z`$-direction, $`𝐁=B\widehat{z}`$, it is sufficient to consider the number density $`n`$ and the $`y`$-component of the number current density $`j_y`$ ($`j_x`$ is fixed by particle conservation). In a 2-vector notation, the polarization function is defined by
$$j^\alpha (q,\omega )=ec^2\mathrm{\Pi }^{\mathrm{e},\alpha \beta }(q,\omega )A_{\mathrm{total}}^\beta (q,\omega ),$$
(A1)
where $`𝐣=(cn,j_y)`$ and $`𝐀=(\mathrm{\Phi }/c,A_y)`$. The scalar potential $`\mathrm{\Phi }`$ and vector potential $`𝐀`$ are related to the electric field $`𝐄`$ and the magnetic induction $`𝐁`$ is given by $`𝐄=\mathrm{\Phi }`$ and $`𝐁=\times 𝐀`$ with $`𝐀=0`$. This choice is known as the Coulomb gauge. The total field is the sum of the external field and the field induced by the current and charge response of the system.
The tensor $`\mathrm{\Pi }^\mathrm{e}(q,\omega )`$ has four components $`\mathrm{\Pi }_{00}^\mathrm{e}`$, $`\mathrm{\Pi }_{01}^\mathrm{e}`$, $`\mathrm{\Pi }_{10}^\mathrm{e}`$, and $`\mathrm{\Pi }_{11}^\mathrm{e}`$, which are related by linear response theory to density-density, density-current, current-density, and current-current correlation functions, respectively. For the drag calculations we are only interested in the density-density response $`\mathrm{\Pi }_{00}^\mathrm{e}`$ which is called $`\mathrm{\Pi }(q,\omega )`$ in the main text.
$`\mathrm{\Pi }^\mathrm{e}(q,\omega )`$ is known only approximately, even at $`B=0`$. In the Random Phase Approximation, the polarization function correlation functions are approximated by their non-interacting fermion forms. At $`B=0`$ we thus use $`\mathrm{\Pi }^\mathrm{e}=2\mathrm{\Pi }_{00}^{(0)}`$ with
$$\mathrm{\Pi }_{00}^{(0)}=\frac{1}{A}\underset{𝐤}{}\frac{n_F(\xi _{𝐤+𝐪})n_F(\xi _𝐤)}{\mathrm{}\omega \xi _{𝐤+𝐪}+\xi _𝐤+i\eta }.$$
(A2)
At zero temperature the wavevector integral can be evaluated analytically:
$`\mathrm{Re}\mathrm{\Pi }_{00}^{(0)}`$ $`=`$ $`{\displaystyle \frac{g_0}{2z}}(2zC_{}[(zu)^21]^{1/2}`$ (A4)
$`C_+[(z+u)^21]^{1/2})`$
$`\mathrm{Im}\mathrm{\Pi }_{00}^{(0)}`$ $`=`$ $`{\displaystyle \frac{g_0}{2z}}(D_{}[1(zu)^2]^{1/2}`$ (A6)
$`D_+[1(z+u)^2]^{1/2})`$
where $`g_0=m_b/2\pi \mathrm{}^2`$, $`z=q/2k_F`$, $`u=\omega /qv_F`$, $`m_b`$ is the band mass, $`v_F`$ is the Fermi velocity, and
$$C_\pm =\{\begin{array}{cc}\frac{z\pm u}{|z\pm u|}\hfill & ,|z\pm u|>1\hfill \\ 0\hfill & ,|z\pm u|<1\hfill \end{array}$$
(A7)
$$D_\pm =\{\begin{array}{cc}0\hfill & ,|z\pm u|>1\hfill \\ 1\hfill & ,|z\pm u|<1\hfill \end{array}$$
(A8)
In the fractional quantum Hall regime, electron-electron interactions cannot be ignored. In Chern-Simons composite fermion theory interactions enter only through the local field produced by attached flux quanta. The composite fermion polarization function is defined by
$$j^\alpha (q,\omega )=ec^2\mathrm{\Pi }_{\mathrm{CF}}^{\alpha \beta }(q,\omega )A_{\mathrm{CF},\mathrm{total}}^\beta (q,\omega ),$$
(A9)
where $`𝐀_{\mathrm{CF},\mathrm{total}}`$ is the total electromagnetic field seen by the Composite Fermions which includes the Chern-Simons field. In the RPA (MRPA) one approximates $`\mathrm{\Pi }_{\mathrm{CF}}`$ by $`\mathrm{\Pi }^{(0)}`$ and relates $`\mathrm{\Pi }^\mathrm{e}`$ to $`\mathrm{\Pi }_{\mathrm{CF}}`$ by
$$[\mathrm{\Pi }^\mathrm{e}]^1=[\mathrm{\Pi }_{\mathrm{CF}}]^1+𝐂+𝐅$$
(A10)
with
$$𝐂=\frac{ec\mathrm{\Phi }_0\stackrel{~}{\varphi }}{q}\left[\begin{array}{cc}0& i\\ i& 0\end{array}\right]$$
(A11)
and
$$𝐅=\frac{\mathrm{\Delta }m}{n_0}\left[\begin{array}{cc}(\frac{\omega }{q})^2& 0\\ 0& c^2\end{array}\right]$$
(A12)
Here $`\mathrm{\Delta }m=m^{}m_b`$ where $`m^{}`$ it the effective mass, and $`n_0`$ is the equilibrium electron number density, $`\mathrm{\Phi }_0=2\pi \mathrm{}/e`$, and $`\stackrel{~}{\varphi }=2`$ at half filling. The difference between RPA and MRPA lies in the matrix $`𝐅`$ which is set to zero for the RPA. Notice that MRPA reduces to RPA if the effective mass is chosen to be equal to the band mass. Solving for $`\mathrm{\Pi }_{00}^\mathrm{e}`$ we get (since $`\mathrm{\Pi }_{01}^{(0)}=\mathrm{\Pi }_{10}^{(0)}=0`$)
$$\mathrm{\Pi }_{00}^\mathrm{e}=\frac{\mathrm{\Pi }_{00}^{(0)}}{1+\frac{\mathrm{\Delta }m}{n_0}\left(\frac{\omega }{q}\right)^2\mathrm{\Pi }_{00}^{(0)}\left(\frac{ec\mathrm{\Phi }_0\stackrel{~}{\varphi }}{q}\right)^2\mathrm{\Pi }_{00}^{(0)}\mathrm{\Pi }_{11}^{(0)}\frac{n_0}{n_0+c^2\mathrm{\Delta }m\mathrm{\Pi }_{11}^{(0)}}}$$
(A13)
In Eq. (A13) the band mass in the expression for $`\mathrm{\Pi }_{00}^{(0)}`$ is replaced by the effective mass. $`\mathrm{\Pi }_{11}^{(0)}`$ is related to the current-current correlation function and is given by
$$\mathrm{\Pi }_{11}^{(0)}=\frac{n_0}{m^{}c^2}+\frac{1}{A}\underset{𝐤}{}\left(\frac{\mathrm{}k_y}{m^{}c}\right)^2\frac{n_F(\xi _{𝐤+𝐪})n_F(\xi _𝐤)}{\mathrm{}\omega \xi _{𝐤+𝐪}+\xi _𝐤+i\eta }$$
(A14)
which at zero temperature is
$`\mathrm{Re}\mathrm{\Pi }_{11}^{(0)}`$ $`=`$ $`{\displaystyle \frac{n_0}{m^{}c^2}}+({\displaystyle \frac{\mathrm{}k_F}{m^{}c}})^2{\displaystyle \frac{g_0^{}}{6z}}(3z+C_+[(z+u)^21]^{3/2}`$ (A16)
$`+C_{}[(zu)^21]^{3/2}(z+u)^3(zu)^3)`$
$`\mathrm{Im}\mathrm{\Pi }_{11}^{(0)}`$ $`=`$ $`({\displaystyle \frac{\mathrm{}k_F}{m^{}c}})^2{\displaystyle \frac{g_0^{}}{6z}}(D_{}[1(zu)^2]^{3/2}`$ (A19)
$`D_+[1(z+u)^2]^{3/2})`$
with $`g_0^{}=m^{}/2\pi \mathrm{}^2`$.
At finite temperature, the expressions for $`\mathrm{\Pi }_{00}^{(0)}`$ and $`\mathrm{\Pi }_{11}^{(0)}`$ must be evaluated numerically. We have done this by evaluating the angular integral analytically and expressing the functions in terms of Fermi-Dirac Integrals which are defined according to
$$_j(x,b)=\frac{1}{\mathrm{\Gamma }(j+1)}_b^{\mathrm{}}\frac{t^j}{1+\mathrm{exp}(tx)}𝑑t.$$
(A20)
We find
$`(g^{})^1\mathrm{Re}\mathrm{\Pi }_{00}^{(0)}`$ $`=`$ $`1{\displaystyle \frac{\sqrt{\pi }\stackrel{~}{T}^{1/2}}{4z}}(_{1/2}({\displaystyle \frac{\mathrm{\Omega }_+^2\stackrel{~}{\mu }}{\stackrel{~}{T}}},0)`$ (A22)
$`_{1/2}({\displaystyle \frac{\mathrm{\Omega }_+^2\stackrel{~}{\mu }}{\stackrel{~}{T}}},{\displaystyle \frac{\mathrm{\Omega }_+^2}{\stackrel{~}{T}}}))`$
$``$ $`\mathrm{sgn}(\mathrm{\Omega }_{}){\displaystyle \frac{\sqrt{\pi }\stackrel{~}{T}^{1/2}}{4z}}(_{1/2}({\displaystyle \frac{\mathrm{\Omega }_{}^2\stackrel{~}{\mu }}{\stackrel{~}{T}}},0)`$ (A24)
$`_{1/2}({\displaystyle \frac{\mathrm{\Omega }_{}^2\stackrel{~}{\mu }}{\stackrel{~}{T}}},{\displaystyle \frac{\mathrm{\Omega }_{}^2}{\stackrel{~}{T}}}))`$
$`(g^{})^1\mathrm{Im}\mathrm{\Pi }_{00}^{(0)}`$ $`=`$ $`{\displaystyle \frac{\sqrt{\pi }\stackrel{~}{T}^{1/2}}{4z}}(_{1/2}({\displaystyle \frac{\stackrel{~}{\mu }\mathrm{\Omega }_{}^2}{\stackrel{~}{T}}},0)`$ (A26)
$`_{1/2}({\displaystyle \frac{\stackrel{~}{\mu }\mathrm{\Omega }_+^2}{\stackrel{~}{T}}},0))`$
$`({\displaystyle \frac{m^{}c}{\mathrm{}k_F}})^2(g_0^{})^1\mathrm{Re}\mathrm{\Pi }_{11}^{(0)}`$ $`=`$ $`u^2z^2/3`$ (A27)
$`+`$ $`{\displaystyle \frac{\sqrt{\pi }\stackrel{~}{T}^{3/2}}{8z}}(_{1/2}({\displaystyle \frac{\mathrm{\Omega }_+^2\stackrel{~}{\mu }}{\stackrel{~}{T}}},0)`$ (A29)
$`_{1/2}({\displaystyle \frac{\mathrm{\Omega }_+^2\stackrel{~}{\mu }}{\stackrel{~}{T}}},{\displaystyle \frac{\mathrm{\Omega }_+^2}{\stackrel{~}{T}}}))`$
$`+`$ $`\mathrm{sgn}(\mathrm{\Omega }_{}){\displaystyle \frac{\sqrt{\pi }\stackrel{~}{T}^{3/2}}{8z}}(_{1/2}({\displaystyle \frac{\mathrm{\Omega }_{}^2\stackrel{~}{\mu }}{\stackrel{~}{T}}},0)`$ (A31)
$`_{1/2}({\displaystyle \frac{\mathrm{\Omega }_{}^2\stackrel{~}{\mu }}{\stackrel{~}{T}}},{\displaystyle \frac{\mathrm{\Omega }_{}^2}{\stackrel{~}{T}}}))`$
$`({\displaystyle \frac{m^{}c}{\mathrm{}k_F}})^2(g_0^{})^1\mathrm{Im}\mathrm{\Pi }_{11}^{(0)}`$ $`=`$ $`{\displaystyle \frac{\sqrt{\pi }\stackrel{~}{T}^{3/2}}{8z}}(_{1/2}({\displaystyle \frac{\stackrel{~}{\mu }\mathrm{\Omega }_{}^2}{\stackrel{~}{T}}},0)`$ (A33)
$`_{1/2}({\displaystyle \frac{\stackrel{~}{\mu }\mathrm{\Omega }_+^2}{\stackrel{~}{T}}},0))`$
Here $`\mathrm{\Omega }_\pm =z\pm u`$, and we have defined the dimensionless quantities $`\stackrel{~}{\mu }=\mu /E_F`$ and $`\stackrel{~}{T}=k_BT/E_F`$ where $`\mu `$ is the chemical potential and $`E_F`$ is the Fermi energy. For parabolic bands $`\stackrel{~}{\mu }/\stackrel{~}{T}=\mathrm{ln}(e^{1/\stackrel{~}{T}}1)`$. For evaluating the Fermi-Dirac Integrals we have used efficient algorithms that are publicly available.
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# 1 Introduction
## 1 Introduction
The coupling of light vector mesons \[$`\rho `$(770) and $`\omega `$(782)\] to low-lying baryon resonances is still to a large extent unknown. This lack of information is a particularly important source of uncertainties in the theoretical description of the propagation of vector mesons in a nuclear medium, where resonance-hole states are expected to contribute largely to the dynamics.
The $`\pi ^{}p\rho ^0n`$ and $`\pi ^{}p\omega n`$ processes have been described recently in the framework of a relativistic coupled-channel model . They are particular processes included in a broader scheme aiming at reproducing data on pion-nucleon elastic scattering and pion-induced production reactions involving the $`\pi \mathrm{\Delta }`$, $`\rho `$N, $`\omega `$N, K$`\mathrm{\Lambda }`$, K$`\mathrm{\Sigma }`$ and $`\eta `$N channels. The model is restricted to s-wave scattering in the $`\rho `$N and $`\omega `$N channels. The corresponding s- and d-wave resonances in the $`\pi `$N channel are generated dynamically. The meson-baryon coupling strengths are determined from the fit to the available data on the channels included in the calculation.
The $`\pi ^{}p\rho ^0n`$ and $`\pi ^{}p\omega n`$ amplitudes are very sensitive to the presence of the s- and d-wave pion-nucleon resonances lying below the vector meson production threshold ($`1.3<\sqrt{s}<1.7`$ GeV). This point is discussed and illustrated in Section 2. Data that directly reflect these amplitudes would provide very useful constraints on the underlying dynamics. The $`\pi ^{}pe^+e^{}n`$ reaction appears as a particularly relevant process to study the $`\pi ^{}p\rho ^0n`$ and $`\pi ^{}p\omega n`$ amplitudes. This reaction offers the possibility to test experimentally the $`\rho ^0`$ and $`\omega `$ strengths below threshold and the quantum interference in the $`e^+e^{}`$ decays of the $`\rho ^0`$\- and $`\omega `$-mesons is very sensitive to the magnitudes and the relative phase of the production amplitudes. In Section 3 we present briefly the formalism and preliminary numerical results. The perspectives of this work are discussed in Section 4.
## 2 The $`\pi ^{}p\rho ^0n`$ and $`\pi ^{}p\omega n`$ amplitudes close to the vector meson production threshold
The $`\pi ^{}p\rho ^0n`$ and $`\pi ^{}p\omega n`$ amplitudes of Ref. entering our calculation of the $`\pi ^{}pe^+e^{}n`$ reaction are displayed in Fig. 1. We shall restrict our discussion to $`e^+e^{}`$ pairs of invariant masses ranging from 0.5 to 0.8 GeV. The exclusive measurement of the $`e^+e^{}n`$ outgoing channel ensures that the $`e^+e^{}`$ pairs come from vector meson decays (pseudoscalar mesons decay into an $`e^+e^{}`$ pair and an additional photon). We recall however that only s- and d-wave pion-nucleon resonances are at present included in the model of Ref. . To be complete, the description of the $`\pi ^{}pe^+e^{}n`$ reaction in the energy range discussed in this work ($`1.2<\sqrt{s}<1.8`$ GeV) should include also the effect of other partial waves. We will return to this question in Section 4.
The $`\pi ^{}p\rho ^0n`$ and $`\pi ^{}p\omega n`$ scattering amplitudes of Fig. 1 illustrate the importance of baryon resonances in vector meson production below threshold. These resonances induce a rich structure in both the real and imaginary parts of the amplitudes. In particular, the presence of the d-wave N\*(1520) resonance is clearly reflected in the J=3/2 amplitudes for $`\rho ^0`$ and $`\omega `$ production. This is an immediate consequence of the strong coupling of the N\*(1520) to both the $`\rho ^0n`$ and $`\omega n`$ channels . The strong couplings imply that there is considerable vector-meson strength in the N\*-hole modes in the nuclear medium.
An experimental test of the N\*N$`\rho ^0`$ and N\*N$`\omega `$ vertices through the $`\pi ^{}pe^+e^{}n`$ reaction below the vector meson production threshold would be a most valuable constraint on the in-medium propagation of $`\rho ^0`$\- and $`\omega `$-mesons.
## 3 The $`\pi ^{}pe^+e^{}n`$ reaction
The $`\pi ^{}p\rho ^0n`$ and $`\pi ^{}p\omega n`$ amplitudes are simply related to the $`\pi ^{}pe^+e^{}n`$ amplitudes through the Vector Dominance assumption . In this picture, the $`e^+e^{}`$ decay of vector mesons is described by their conversion into time-like photons which subsequently materialize into $`e^+e^{}`$ pairs. The magnitude of the coupling constants $`f_\rho `$ and $`f_\omega `$, which characterize the conversion of $`\rho `$\- and $`\omega `$-mesons into photons, is determined from the measured partial widths of $`\rho ^0`$\- and $`\omega `$-mesons into $`e^+e^{}`$ pairs . The relative phase of the $`\rho `$ and $`\omega `$ amplitudes is not determined by hadronic observables. We determine this phase in each channel by comparing with the photon-decay helicity amplitudes of the corresponding resonance , assuming Vector Meson Dominance. We use $`f_\rho `$=0.036 GeV<sup>2</sup> and $`f_\omega `$=0.011 GeV<sup>2</sup> .
The squared amplitude for the $`\pi ^{}pe^+e^{}n`$ reaction with intermediate $`\rho ^0`$\- and $`\omega `$-mesons in the Vector Dominance Model is illustrated in Fig. 2. Schematically, this quantity can be written as
$`\left|<ne^+e^{}|\pi ^{}p>\right|^2={\displaystyle \frac{|<e^+e^{}|\gamma >|^2}{m^4}}{\displaystyle \frac{f_\rho _{\pi ^{}p\rho ^0n}}{m^2m_\rho ^2+im_\rho \mathrm{\Gamma }_\rho (m)}}`$
$`+{\displaystyle \frac{f_\omega _{\pi ^{}p\omega n}}{m^2m_\omega ^2+im_\omega \mathrm{\Gamma }_\omega (m)}}^2,`$ (1)
where the first term of the right-hand side describes the propagation of the time-like photon and its decay into an $`e^+e^{}`$ pair of invariant mass m and the second term contains the vector meson production dynamics. The vector mesons are characterized by their mass m<sub>V</sub> and energy-dependent width $`\mathrm{\Gamma }_V(m)`$. The interference of the complex $`_{\pi ^{}p\rho ^0n}`$ and $`_{\pi ^{}p\omega n}`$ amplitudes (Fig. 1) in the $`\pi ^{}pe^+e^{}n`$ cross section is sensitive to their relative phase. The importance of measuring such a phase in the $`e^+e^{}`$ or $`\pi ^+\pi ^{}`$ decays of $`\rho ^0`$\- and $`\omega `$-mesons has been evidenced by the contribution of such data to the understanding of other processes, like the photoproduction of $`\rho ^0`$\- and $`\omega `$-mesons in the diffractive regime ($`\gamma Bee^+e^{}Be`$) and the $`e^+e^{}\pi ^+\pi ^{}`$ reaction .
We indicate the magnitude of the $`\rho ^0\omega `$ interference in the $`\pi ^{}pe^+e^{}n`$ reaction as function of the total center of mass energy in Fig. 3. We have selected $`e^+e^{}`$ pairs of invariant mass m=0.55 GeV. This figure illustrates the role of baryon resonances with masses in the range of 1.5 to 1.6 GeV in generating strong interference effects.
Above the vector meson threshold, the $`\rho ^0\omega `$ interference in the $`\pi ^{}pe^+e^{}n`$ cross section is particularly interesting for $`e^+e^{}`$ pair invariant masses close to the $`\omega `$ mass. This effect is manifested in the invariant mass spectrum displayed in Fig. 4 ($`\sqrt{s}`$=1.8 GeV). The model of Ref. for the $`_{\pi ^{}p\rho ^0n}`$ and $`_{\pi ^{}p\omega n}`$ amplitudes predicts a constructive interference at this energy. This feature appears to be a very sensitive test of the model.
A detailed discussion of these interference patterns will be presented in a forthcoming publication .
## 4 Perspectives
The study of the $`\pi ^{}pe^+e^{}n`$ reaction provides a particularly stringent test of the $`\pi ^{}p\rho ^0n`$ and $`\pi ^{}p\omega n`$ amplitudes close and below the vector meson production threshold ($`1.2<\sqrt{s}<1.8`$ GeV).
We have computed the cross section of the $`\pi ^{}pe^+e^{}n`$ reaction using the model of Ref. for the vector meson production amplitude and indicated its main features as function of the total center of mass energy.
A natural extension of the present work would be to include the p-wave pion-nucleon resonances in the coupled channel scheme of Ref. , thereby increasing the expected domain of validity of the $`\pi ^{}p\rho ^0n`$ and $`\pi ^{}p\omega n`$ amplitudes. Projecting the coupled-channel amplitudes on specific s- and t-channel exchanges could be a useful step in providing a simple interpretation of our numerical results.
We note that the study of the quantum interference of $`\rho ^0`$\- and $`\omega `$-mesons produced in the $`\pi ^{}p\rho ^0n`$ and $`\pi ^{}p\omega n`$ reactions in other channels than the $`e^+e^{}`$ decay ($`\pi ^0\gamma `$ for example) may also be of interest.
Data on the $`\pi ^{}pe^+e^{}n`$ cross section in the energy range considered in this work are at present not available. Such measurements would provide an important test of the dynamics in a reaction which is crucial for the understanding of the in-medium propagation of vector mesons.
## 5 Acknowledgements
One of us (M. S.) acknowledges the generous hospitality of the Theory Group of GSI, where much of this work was done.
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# Disentangling Intertwined Embedded-States and Spin Effects in Light-Front Quantization
## I Introduction
With the advent of Light-Front field theory, one can be quite hopeful to develop a connection between Quantum Chromodynamics (QCD) and the relativistic constituent quark model which is used in various electroweak form factor calculations. QCD provides a fundamental description of hadronic and nuclear structure in terms of elementary quark and gluon degrees of freedom. It is very succesful in the perturbative regime, for instance when it is used in the explanation of the evolution of distribution functions in deep-inelastic scattering. There, the basic hard interaction is described by perturbative QCD and the Electro-weak interaction. The soft parts, distribution and fragmentation functions, are non-perturbative ingredients, which are less well understood. Lattice calculations are becoming increasingly accurate, but there is still much room and need for other non-perturbative approaches.
It is part of the nature of the description of hadronic systems in terms of quarks and gluons that the characteristic momenta are of the same order or even very much larger than the masses of the particles involved. Therefore a relativistic treatment is called for. A very promising technique is Light-Front Dynamics (LFD), which treats relativistic many-body effects in a consistent way . In LFD a Fock-space expansion of bound states is made. The wave function $`\psi _n(x_i,k_i^{},\lambda _i)`$ describes the component with $`n`$ constituents, with longitudinal momentum fraction $`x_i`$, perpendicular momentum $`k_i^{}`$ and helicity $`\lambda _i`$, $`i=1,\mathrm{},n`$. It is the aim of LFD to determine those wave functions and use them in conjunction with hard scattering amplitudes to describe the properties of hadrons and their response to electroweak probes.
Recently, important steps were taken towards a realization of this goal. In the work of Brodsky,Hiller and McCartor , it is demonstrated how to solve the problem of renormalizing light-front Hamiltonian theories while maintaining Lorentz symmetry and other symmetries. (The genesis of the work presented in may be found in and additional examples including the use of LFD methods to solve the bound-state problems in field theory can be found in the review ). These results are indicative of the great potential of LFD for a fundamental description of non-perturbative effects in QCD. However, at present there are no realistic results available for wave functions of hadrons based on QCD alone. In order to calculate the response of hadrons to external probes, one might resort to the use of model wave functions. This way to estimate matrix elements was used by Ji et al.. The same reasons that make LFD so attractive to solve bound-state problems in field theory make it also useful for a relativistic description of nuclear systems. Presently, it is realized that a parametrization of nuclear reactions in terms of non-relativistic wave functions must fail. LF methods have the advantage that they are formally similar to time-ordered many-body theories, yet provide relativistically invariant observables.
Until now we have sketched a rather rosy picture for the application of LFD to hadron physics. However, not all is well and this is just the reason for the present investigation. Since the 1980’s it has been assumed that the observables computed in the framework of LFD using the methods of perturbation theory are invariants, just as in covariant perturbation theory. Many authors have shown that LFD has this feature in particular cases and some years ago some general statements to the same effect could be made . A case in point is the calculation of a current matrix element in quantum field theory. A typical amplitude is given by the triangle diagram. One encounters this diagram e.g. when computing the pion form factor (see Fig. 1).
The vertices denoted by $`\mathrm{\Psi }`$ are coupling constants in covariant perturbation theory. The hard scattering process is the absorption of a photon of momentum $`q`$ by a (anti-)quark. In the LFD approach the covariant amplitude is replaced by a series of LF time-ordered diagrams. In the case of the triangle diagram they are depicted in Fig. 2.
The first of these two diagrams is easily interpreted in terms of the LF wave functions $`\mathrm{\Psi }`$. However, the other diagram has a vertex that can again be written in the same way as before, but it contains also another vertex, denoted by $`\mathrm{\Psi }^{}`$, that cannot be written as a LF wave function. It is a new element in LFD. We call this vertex the non-wave-function vertex. In order to compute the form factors in the time-like region, the contributions from these vertices must be included. Semileptonic meson decay processes also require the contributions from these vertices. One may try to avoid using them by choosing special kinematic conditions. It is known however that this will not be a simple task .
In the present work, we investigate the contributions from the non-wave-function vertices. We construct both the wave-function and non-wave-function vertices using pointlike covariant ones. The model used here is essentially an extension of Mankiewicz and Sawicki’s $`(1+1)`$-dimensional quantum field theory model , which was later reinvestigated by several others . The starting model wave function is the solution of the covariant Bethe-Salpeter equation in the ladder approximation with a relativistic version of the contact interaction . The covariant model wave function is a product of two free single particle propagators, the overall momentum-conserving Dirac delta function, and a constant vertex function. Consequently, all our form factor calculations are various ways of evaluating the Feynman triangle diagram in quantum field theory.
The importance of the contributions of the non-wave-function vertices was investigated in two cases: the electromagnetic form factors of a scalar and a pseudoscalar meson with spin-1/2 constituents. We also calculated the same form factor of a scalar meson with spin-0 constituents to see the spin effects. In 3+1 dimensions both the covariant and the LF calculations are divergent and the model without any smeared vertex for the fermion loop is not well defined. This is in dramatic contrast to the case of spin-0 (boson) constituents, where regularization is not needed at all. In order to disentangle the issue of the non-wave-function vertices from the need of regularization, we performed our calculations in 1+1 dimensions, where at least the covariant calculations for spin-1/2 constituents give finite results.
It is commonly believed and widely used that the LF energy integration of the covariant Feynman amplitude generates the corresponding equivalent amplitude in the LFD. As we will show in this work, however, the equivalence between the LFD and the covariant Feynman calculation is not always guaranteed. The bad component of the current, $`J^{}`$, with spin-1/2 constituents exhibits a persistent end-point singularity in the contribution from the non-wave-function vertex. Unless the divergence in this contribution is properly subtracted, the singular behavior leads to an infinitely different result from that obtained by the covariant Feynman calculation. Ensuring the equivalence to the Feynman amplitude, we have identified the divergent term that needs to be removed from $`J^{}`$. Only after the identified term is subtracted, the result is covariant and satisfies the current conservation.
In the next Section (Section II), we present both the covariant Feynman calculations and the LF calculations using the LF energy integration for the electromagnetic form factors of a pseudoscalar and a scalar meson with spin-1/2 constituents as well as the same form factor of a scalar meson with spin-0 constituents. Section III contains the numerical estimates of both the wave-function and non-wave-function vertices to the electromagnetic form factors of each case presented in Section II. The conclusion and discussion follow in Section IV.
## II Calculations
The electromagnetic form factors can be extracted from the matrix elements of the current $`J^\mu `$
$$p^{}|J^\mu |p=ie_m(p^\mu +p^\mu )F(q^2),$$
(1)
where $`e_m`$ is the charge of the meson and $`q^2=(p^{}p)^2`$ is the square of the four momentum transfer. If one uses the plus-component, $`J^+=(J^0+J^3)/\sqrt{2}`$, the LF calculation gives two finite contributions, the wave-function part and the non-wave-function part, that add up to the covariant result, as expected. The importance of the non-wave-function contribution varies strongly with the momentum transfer and depends sensitively on the binding energy of the meson. For small values of $`q^2`$ and small binding energy, the wave-function part is dominant, but elsewhere the non-wave-function is essential for agreement between the LF calculation and the covariant results.
The form factor can also be extracted from the minus-component of the current, $`J^{}=(J^0J^3)/\sqrt{2}`$. Covariance guarantees that it makes no difference whether the form factor is determined using the plus or the minus current matrix element. As LFD is not manifestly covariant, it may happen that $`J^{}`$ leads to a form factor different from the one determined using $`J^+`$. As we show in this Section, the matrix element of $`J^{}`$ diverges in LFD. Unless one regulates $`J^{}`$, the current cannot be conserved. To assure the current conservation, it is crucial to identify the term that causes the divergence. We have identified this term exactly and found that it is an infinite function of the momentum transfer. If this infinite term is subtracted, the two LF contributions become finite as it must be in the conserved current. Moreover, their sum equals again the covaraint result as expected. However, the regularized LF contributions are different from the two parts of the form factor extracted from the plus current. The differences grow with increasing binding energy.
### A Pseudoscalar Meson with the Fermion Loop
The covariant fermion triangle-loop (Fig. 1) for the pseudoscalar meson leads to the amplitude given by
$$p^{}|J^\mu |p=4N\frac{d^2k}{(2\pi )^2}\frac{(m^2k^2+pp^{})k^\mu +(k^2m^2kp^{})p^\mu +(k^2m^2kp)p^\mu }{(k^2m^2+iϵ)((kp)^2m^2+iϵ)((kp^{})^2m^2+iϵ)},$$
(2)
where $`m`$ is the fermion mass and $`N`$ modulo the obvious charge factor $`e_m`$ is the normalization constant fixed by the unity of the form factor at zero momentum transfer. Even though we will present the unequal constituent mass case such as the kaon in the next Section of our numerical analysis, for the clarity of presentation we will focus in this Section on the equal mass case, such as the pion, only.
The usual Feynman parametrization and the covariant integration yields
$$p^{}|J^\mu |p=i(p^\mu +p^\mu )\frac{N}{\pi }_{0}^{}{}_{}{}^{1}𝑑x_{0}^{}{}_{}{}^{1x}𝑑y\frac{2m^2(1x)(m^2+(x+y1)^2M^2xyq^2)x}{((x+y)(x+y1)M^2+m^2xyq^2)^2},$$
(3)
where $`M`$ is the meson mass. For $`q^2=0`$, the integration leads to fix the normalization as
$$1/N=\frac{4m^2}{\pi M(4m^2M^2)}\left[\frac{M}{4m^2}+\frac{1}{\sqrt{4m^2M^2}}\mathrm{arctan}\left(\frac{M}{\sqrt{4m^2M^2}}\right)\right].$$
(4)
In LFD, the form factor $`F(q^2)`$ can be obtained by calculating either $`p^{}|J^+|p`$ or $`p^{}|J^{}|p`$. In principle, the result must be identical to the above covariant Feynman result regardless of which component of the current is used. However, this is not necessarily the case as we demonstrate in the following.
First, the calculation of $`p^{}|J^+|p`$ integrating out the LF energy $`k^{}`$ in Eq. (2) yields $`F(q^2)`$ given by
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{N}{\pi (2+\alpha )}}[{\displaystyle _{0}^{}{}_{}{}^{1}}dx{\displaystyle \frac{(1+\alpha )^2m^2}{[m^2x(1x)M^2][(1+\alpha )^2m^2(1x)(\alpha +x)M^2]}}`$ (5)
$`+`$ $`{\displaystyle _{0}^{}{}_{}{}^{\alpha }}dx{\displaystyle \frac{(1+\alpha )(\alpha x)[x(\alpha x)M^2(1+\alpha )^2m^2]}{\alpha [(\alpha x)(1+x)M^2(1+\alpha )^2m^2][x(\alpha x)M^2+(1+\alpha )m^2]}}],`$ (6)
where $`\alpha `$ is given by $`q^2=\frac{\alpha ^2M^2}{1+\alpha }`$. In Eq. (6) the first and second terms correspond to the contributions from the wave-function and non-wave-function vertices depicted in the first and second diagrams in Fig. 2, respectively. We have verified that the contribution from the non-wave-function part vanishes at $`\alpha =0`$, i.e. at $`q^2=0`$, indicating the absence of a zero-mode contribution in the good component of the current $`J^+`$. Adding both contributions in Eq. (6), we obtain
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{2N(1+\alpha )m^2}{\pi \alpha M[(2+\alpha )^2m^2(1+\alpha )M^2]}}`$ (9)
$`[{\displaystyle \frac{2+2\alpha +\alpha ^2}{\sqrt{4(1+\alpha )m^2+\alpha ^2M^2}}}\mathrm{Artanh}\left({\displaystyle \frac{\alpha M}{\sqrt{4(1+\alpha )m^2+\alpha ^2M^2}}}\right)`$
$`\text{ }+{\displaystyle \frac{2\alpha }{\sqrt{4m^2M^2}}}\mathrm{arctan}\left({\displaystyle \frac{M}{\sqrt{4m^2M^2}}}\right)].`$
This result is identical to the form factor obtained by the covariant Feynman calculation given by Eq. (3) and F(0)=1 gives the same normalization given in Eq. (4).
On the other hand, the calculation of $`p^{}|J^{}|p`$ yields $`F(q^2)`$ given by
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{N}{\pi (2+\alpha )}}[(1+\alpha )M^2{\displaystyle _{0}^{}{}_{}{}^{1}}dx{\displaystyle \frac{(1x)^2}{[m^2x(1x)M^2][(1+\alpha )^2m^2(1x)(\alpha +x)M^2]}}`$ (10)
$``$ $`{\displaystyle \frac{(1+\alpha )^2m^2}{\alpha M^2}}{\displaystyle _{0}^{}{}_{}{}^{\alpha }}{\displaystyle \frac{dx}{\alpha x}}{\displaystyle \frac{(1+\alpha )^2m^2+\{(1+\alpha )^2(1+x)\}(\alpha x)M^2}{[(\alpha x)(1+x)M^2(1+\alpha )^2m^2][x(\alpha x)M^2+(1+\alpha )m^2]}}],`$ (11)
where again it is apparent that the first and second terms correspond to the contributions from the wave-function and non-wave-function vertices, respectively. However, the non-wave-function part shows the end-point singularity coming from $`\frac{1}{\alpha x}`$. Without subtracting the end-point singularity, the result is infinitely different from that obtained in the $`p^{}|J^+|p`$ calculation. This is an astonishing result that deviates from the common belief in the equivalence of the LFD and the covariant Feynman calculation. Neither covariance nor current conservation is satisfied without a certain adjustment. In order to identify the term that must be subtracted, we rewrite the above equation as follows:
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{N}{\pi (2+\alpha )}}[(1+\alpha )M^2{\displaystyle _{0}^{}{}_{}{}^{1}}dx{\displaystyle \frac{(1x)^2}{[m^2x(1x)M^2][(1+\alpha )^2m^2(1x)(\alpha +x)M^2]}}`$ (13)
$`\text{ }+{\displaystyle _{0}^{}{}_{}{}^{\alpha }}dx{\displaystyle \frac{R(x,\alpha )}{\alpha x}}],`$
where $`R(x,\alpha )`$ is defined by
$$R(x,\alpha )=\frac{(1+\alpha )^2m^2[(1+\alpha )^2m^2+\{(1+\alpha )^2(1+x)\}(\alpha x)M^2}{\alpha M^2[(\alpha x)(1+x)M^2(1+\alpha )^2m^2][x(\alpha x)M^2+(1+\alpha )m^2]}.$$
(14)
In order to obtain the identical result to Eq. (9) from the $`p^{}|J^+|p`$ calculation, we find that $`R(\alpha ,\alpha )`$ must be subtracted from the numerator of the non-wave-function part integrand, i.e.,
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{N}{\pi (2+\alpha )}}[(1+\alpha )M^2{\displaystyle _{0}^{}{}_{}{}^{1}}dx{\displaystyle \frac{(1x)^2}{[m^2x(1x)M^2][(1+\alpha )^2m^2(1x)(\alpha +x)M^2]}}`$ (16)
$`\text{ }+{\displaystyle _{0}^{}{}_{}{}^{\alpha }}dx{\displaystyle \frac{R(x,\alpha )R(\alpha ,\alpha )}{\alpha x}}].`$
The subtracted term $`R(\alpha ,\alpha )=\frac{1+\alpha }{\alpha M^2}`$ depends on the momentum transfer and never vanishes. While the subtracted result Eq. (16) is identical to Eq. (9), it is interesting to note that the zero-mode contribution does not vanish in the $`p^{}|J^{}|p`$ calculation as the non-wave-function part in Eq. (16) still survives even at $`q^2=0`$. However, the subtracted result Eq. (16) with the zero-mode contribution assures covariance and satisfies current conservation.
### B Scalar Meson with the Fermion Loop
The Feynman parametrization and the covariant integration of the fermion triangle-loop for the scalar meson gives the following amplitude:
$`p^{}|J^\mu |p`$ $`=`$ $`4N{\displaystyle \frac{d^2k}{(2\pi )^2}\frac{(3m^2+k^2pp^{})k^\mu (k^2+m^2kp^{})p^\mu (k^2+m^2kp)p^\mu }{(k^2m^2+iϵ)((kp)^2m^2+iϵ)((kp^{})^2m^2+iϵ)}}`$ (17)
$`=`$ $`i(p^\mu +p^\mu ){\displaystyle \frac{N}{\pi }}{\displaystyle _{0}^{}{}_{}{}^{1}}𝑑x{\displaystyle _{0}^{}{}_{}{}^{1x}}𝑑y{\displaystyle \frac{x((1xy)^2M^2m^2xyq^2)}{((x+y)(x+y1)M^2+m^2xyq^2)^2}},`$ (18)
where the normalization $`N`$ is again fixed by $`F(0)=1`$ and given by
$$1/N=\frac{4m^2}{\pi M^3}\left[\frac{M}{4m^2}\frac{1}{\sqrt{4m^2M^2}}\mathrm{arctan}\left(\frac{M}{\sqrt{4m^2M^2}}\right)\right].$$
(19)
In LFD, the calculation of $`p^{}|J^+|p`$ leads to $`F(q^2)`$ given by
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{N}{\pi (2+\alpha )}}[{\displaystyle _{0}^{}{}_{}{}^{1}}dx{\displaystyle \frac{(1+\alpha )\{2(1x)(2x+\alpha )(1+\alpha )\}m^2}{[m^2x(1x)M^2][(1+\alpha )^2m^2(1x)(\alpha +x)M^2]}}`$ (20)
$`+`$ $`{\displaystyle _{0}^{}{}_{}{}^{\alpha }}dx{\displaystyle \frac{(1+\alpha )(\alpha x)[(1+\alpha )(1+4x\alpha )m^2x(\alpha x)M^2]}{\alpha [(\alpha x)(1+x)M^2(1+\alpha )^2m^2][x(\alpha x)M^2+(1+\alpha )m^2]}}],`$ (21)
where the contribution from the non-wave-function part again vanishes at $`q^2=0`$, indicating the absence of zero-mode contribution in the good component of the current, $`J^+`$. Adding both contributions, we find
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{2N(1+\alpha )m^2}{\pi \alpha M^3[(2+\alpha )^2m^2(1+\alpha )M^2]}}`$ (24)
$`[{\displaystyle \frac{8(1+\alpha )m^2(2+2\alpha \alpha ^2)M^2}{\sqrt{4(1+\alpha )m^2+\alpha ^2M^2}}}\mathrm{Artanh}\left({\displaystyle \frac{\alpha M}{\sqrt{4(1+\alpha )m^2+\alpha ^2M^2}}}\right)`$
$`\text{ }{\displaystyle \frac{2\alpha }{\sqrt{4m^2M^2}}}\mathrm{arctan}\left({\displaystyle \frac{M}{\sqrt{4m^2M^2}}}\right)].`$
This result is identical to the form factor obtained by the covariant Feynman calculation given by Eq. (18) and F(0)=1 gives the same normalization presented in Eq. (18).
However, the calculation of $`p^{}|J^{}|p`$ generates an end-point point singularity similar to the one observed in the pseudoscalar case. Defining the function
$$S(x,\alpha )=\frac{(1+\alpha )^2m^2[(1+\alpha )(1+4x3\alpha )m^2+\{\alpha (2\alpha )+(2\alpha 1)x\}(\alpha x)M^2]}{\alpha M^2[(\alpha x)(1+x)M^2(1+\alpha )^2m^2][x(\alpha x)M^2+(1+\alpha )m^2]},$$
(25)
we find $`F(q^2)`$ given by
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{N}{\pi (2+\alpha )}}[{\displaystyle \frac{(1+\alpha )}{M^2}}{\displaystyle _{0}^{}{}_{}{}^{1}}dx{\displaystyle \frac{4(1+\alpha )m^42(2+\alpha )(1x)m^2M^2+(1x)^2M^4}{[m^2x(1x)M^2][(1+\alpha )^2m^2(1x)(\alpha +x)M^2]}}`$ (27)
$`\text{ }+{\displaystyle _{0}^{}{}_{}{}^{\alpha }}dx{\displaystyle \frac{S(x,\alpha )}{\alpha x}}],`$
where the non-wave-function part again shows the end-point singularity coming from $`\frac{1}{\alpha x}`$. In order to obtain the identical result to Eq. (24) from the $`p^{}|J^+|p`$ calculation, we find that $`S(\alpha ,\alpha )`$ must be subtracted from the numerator of the non-wave-function part integrand, i.e.,
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{N}{\pi (2+\alpha )}}[{\displaystyle \frac{(1+\alpha )}{M^2}}{\displaystyle _{0}^{}{}_{}{}^{1}}dx{\displaystyle \frac{4(1+\alpha )m^42(2+\alpha )(1x)m^2M^2+(1x)^2M^4}{[m^2x(1x)M^2][(1+\alpha )^2m^2(1x)(\alpha +x)M^2]}}`$ (29)
$`\text{ }+{\displaystyle _{0}^{}{}_{}{}^{\alpha }}dx{\displaystyle \frac{S(x,\alpha )S(\alpha ,\alpha )}{\alpha x}}].`$
The subtracted term $`S(\alpha ,\alpha )=\frac{1+\alpha }{\alpha M^2}`$ again cannot vanish. The zero-mode contribution is also visible since the second term in Eq. (29) doesn’t vanish even if $`q^2=0`$. The subtracted result Eq. (29), however, assures the covariance and satisfies the current conservation as in the previous calculation of the pseudoscalar meson.
### C Scalar Meson with the Boson Loop
For a comparison with the scalar constituents neglecting spin effects, we present in this subsection the calculation of $`F(q^2)`$ for a scalar meson with a boson loop. The Feynman parametrization and the covariant integration of the boson triangle-loop for the scalar meson gives the following amplitude:
$`p^{}|J^\mu |p`$ $`=`$ $`N{\displaystyle \frac{d^2k}{(2\pi )^2}\frac{p^\mu +p^\mu 2k^\mu }{(k^2m^2+iϵ)((kp)^2m^2+iϵ)((kp^{})^2m^2+iϵ)}}`$ (30)
$`=`$ $`i(p^\mu +p^\mu ){\displaystyle \frac{N}{4\pi }}{\displaystyle _{0}^{}{}_{}{}^{1}}𝑑x{\displaystyle _{0}^{}{}_{}{}^{1x}}𝑑y{\displaystyle \frac{2x1}{((x+y)(x+y1)M^2+m^2xyq^2)^2}},`$ (31)
where the normalization $`N`$ is given by
$$1/N=\frac{1}{2\pi M^2(4m^2M^2)}\left[1+\frac{2(2m^2M^2)}{M\sqrt{4m^2M^2}}\mathrm{arctan}\left(\frac{M}{\sqrt{4m^2M^2}}\right)\right].$$
(32)
In LFD, the calculation of $`p^{}|J^+|p`$ leads to $`F(q^2)`$ given by
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{N}{4\pi (2+\alpha )}}[1{\displaystyle _{0}^{}{}_{}{}^{1}}dx{\displaystyle \frac{(1+\alpha )(2x+\alpha )(1x)}{[m^2x(1x)M^2][(1+\alpha )^2m^2(1x)(\alpha +x)M^2]}}`$ (34)
$`+{\displaystyle _{0}^{}{}_{}{}^{\alpha }}dx{\displaystyle \frac{(1+\alpha )^2(\alpha 2x)(\alpha x)}{\alpha [(\alpha x)(1+x)M^2(1+\alpha )^2m^2][x(\alpha x)M^2+(1+\alpha )m^2]}}],`$
where the contribution from the non-wave-function part again vanishes at $`q^2=0`$, indicating the absence of a zero-mode contribution in the good component of the current $`J^+`$ as in the fermion loop cases. Adding both contributions, we find
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{N(1+\alpha )}{2\pi \alpha M^3[(2+\alpha )^2m^2(1+\alpha )M^2]}}`$ (37)
$`[\sqrt{4(1+\alpha )m^2+\alpha ^2M^2}\mathrm{Artanh}\left({\displaystyle \frac{\alpha M}{\sqrt{4(1+\alpha )m^2+\alpha ^2M^2}}}\right)`$
$`\text{ }+{\displaystyle \frac{2\alpha (2m^2M^2)}{\sqrt{4m^2M^2}}}\mathrm{arctan}\left({\displaystyle \frac{M}{\sqrt{4m^2M^2}}}\right)].`$
This result is identical to the form factor obtained by the covariant Feynman calculation given by Eq. (31) and F(0)=1 gives the same normalization presented in Eq. (32).
Similarly, the calculation of $`p^{}|J^{}|p`$ generates $`F(q^2)`$ given by
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{N}{4\pi (2+\alpha )}}[{\displaystyle \frac{(1+\alpha )}{M^2}}{\displaystyle _{0}^{}{}_{}{}^{1}}dx{\displaystyle \frac{2(1+\alpha )m^2(2+\alpha )(1x)M^2}{[m^2x(1x)M^2][(1+\alpha )^2m^2(1x)(\alpha +x)M^2]}}`$ (39)
$`+{\displaystyle \frac{(1+\alpha )^2}{\alpha M^2}}{\displaystyle _{0}^{}{}_{}{}^{\alpha }}dx{\displaystyle \frac{2(1+\alpha )m^2+\alpha (\alpha x)M^2}{[(\alpha x)(1+x)M^2(1+\alpha )^2m^2][x(\alpha x)M^2+(1+\alpha )m^2]}}].`$
Unlike the cases of fermion constituents, however, the non-wave-function part here does not exhibit the end-point singularity. Without any adjustment, we find that Eq. (39) is identical to Eq. (37) obtained from $`p^{}|J^+|p`$. Thus, the result is automatically covariant and satisfies current conservation. However, the zero-mode contribution is still present in $`J^{}`$ current as one can easily see that the second term in Eq. (39) doesn’t vanish when $`\alpha `$ goes to zero. In the next section, we numerically estimate the importance of the non-wave-function part in all three cases that we presented in this section.
## III Numerical Results
We have estimated both the wave-function and non-wave-function vertices to the electromagnetic form factors of each case presented in Section II. For the three cases (A,B,C) presented in Section II, we show the numerical results of the form factor calculated via the minus component as well as the plus component of the current. We denote the contributions from the wave-function and non-wave-function vertices as $`F_{\mathrm{val}}^{+()}`$ and $`F_{\mathrm{nv}}^{+()}`$, respectively when the plus(minus) component of the current is used. The sum of the two contributions is denoted by $`F_{\mathrm{tot}}^{+()}`$,i.e. $`F_{\mathrm{tot}}^{+()}=F_{\mathrm{val}}^{+()}+F_{\mathrm{nv}}^{+()}`$.
For the numerical computation, we take for the experimental meson masses of the pion and the kaon $`m_\pi =0.140`$ GeV and $`m_K=0.494`$ GeV resp. and vary the quark masses to investigate the binding-energy dependence of the meson form factors. We call the pseudoscalar meson with the equal quark masses and mass $`m_\pi =0.140`$ GeV the ”pion”. Likewise, the pseudoscalar meson with the unequal quark masses and the meson mass $`m_K=0.494`$ GeV is called the ”kaon”.
The ”pion” form factor with the quark mass $`m_q=0.250`$ GeV is shown in Fig. 3. When the plus current is used, the valence contribution $`F_{\mathrm{val}}^+`$ diminishes very quickly as $`Q^2`$ gets larger even though the normalization at $`Q^2=0`$ is entirely coming from the valence part as we pointed out in Section II. The crossing between $`F_{\mathrm{val}}^+`$ and $`F_{\mathrm{nv}}^+`$ appears at $`Q^2`$ below $`0.05`$ GeV<sup>2</sup>. When the minus current is used, however, the valence contribution $`F_{\mathrm{val}}^{}`$ is negligible even at $`Q^2=0`$ and the entire result is essentially given by $`F_{\mathrm{nv}}^{}`$. The value of $`F_{\mathrm{nv}}^{}(0)`$ corresponds to the zero-mode contribution in the minus current $`J^{}`$ and it is interesting to note that more than $`90\%`$ of the form factor at $`Q^2=0`$ is contributed by the zero-mode. Since $`F_{\mathrm{tot}}^+=F_{\mathrm{tot}}^{}`$ exactly coincide with the covariant result obtained by Eq.(3) for all $`Q^2`$ as they must, only a single solid line is depicted in Fig. 3. The same applies to all of the other figures presented in this work.
In Fig. 4, we present the results for the ”pion” by changing the quark masses in the following way: $`m_q=0.140`$, 0.077 and 0.0707 GeV, respectively. The closer $`m_q`$ is to $`m_\pi /2=0.07`$ GeV, the smaller the binding energy gets, and the slope of $`F_{\mathrm{tot}}^{+/}`$ at $`Q^2=0`$ (or the charge radius) increases with decreasing quark mass, as expected. We find that the crossing between $`F_{\mathrm{val}}^+`$ and $`F_{\mathrm{nv}}^+`$ occurs at a larger value of $`Q^2`$ and $`F_{\mathrm{val}}^{}`$ becomes larger near $`Q^2=0`$ as the binding gets weaker. This may be explained by the reduction of the probability to generate the non-wave-function vertex (or the higher Fock state) compared to the valence state as the interaction between the constituents gets weaker. Thus, in the weaker binding, $`F_{\mathrm{val}}^+`$ dominates over $`F_{\mathrm{nv}}^+`$. Similarly, $`F_{\mathrm{val}}^{}`$ becomes the main contribution near $`Q^2=0`$ which is the only region where the form factor exists in the weak binding limit. Consequently, the zero-mode $`F_{\mathrm{nv}}^{}(0)`$ gets substantially diminished as shown in Fig. 4.
The ”kaon” form factor with $`m_q=0.25`$ GeV and $`m_s=0.37`$ GeV is shown in Fig. 5. Compared to the ”pion” case, the dominance of $`F_{\mathrm{val}}^+`$ extends to the larger $`Q^2`$ region and the crossover between $`F_{\mathrm{val}}^+`$ and $`F_{\mathrm{nv}}^+`$ is postponed beyond the range of $`Q^2=1`$ GeV<sup>2</sup>. The zero-mode $`F_{\mathrm{nv}}^{}(0)`$ is also much smaller than in the ”pion” case even though $`F_{\mathrm{nv}}^{}`$ rises very quickly as $`Q^2`$ gets away from the zero range. As one can see in Figs.5(a) and (b), the contribution from the heavier quark struck by the photon is larger than that from the lighter quark struck by the photon. We have indeed confirmed that as $`m_s`$ gets larger only the contribution from the heavy quark struck by the photon dominates as expected.
In Fig. 6, the form factors of the scalar partner to the ”pion”, which we call ”s-pion” in the following, with $`m_q=0.25`$, 0.14, 0.077, and 0.0707 GeV are presented for comparison with the ”pion” case. The basic features of $`F_{\mathrm{val}}^+`$ and $`F_{\mathrm{nv}}^+`$ near $`Q^2=0`$ are same as in the ”pion” case because $`F_{\mathrm{val}}^+(0)=1`$ must hold for any meson. However, as the binding gets weaker, we find that the ”s-pion” form factors $`F_{\mathrm{tot}}^{+/}`$ change sign at smaller $`Q^2`$-values. This indicates that electron scattering off the ”s-pion” not only has zero cross section at a certain electron energy, but also that the electron energy that yields zero scattering gets smaller as the binding of the ”s-pion” is weaker. Another dramatic difference from the pseudoscalar meson is the astonishing cancellation between $`F_{\mathrm{val}}^{}`$ and $`F_{\mathrm{nv}}^{}`$. Especially in the strong binding case, both $`F_{\mathrm{val}}^{}`$ and $`F_{\mathrm{nv}}^{}`$ are huge but they cancel in a very remarkable way to yield exactly the same result as $`F_{\mathrm{tot}}^+`$.
In Fig. 7, the form factor of the scalar partner to the ”kaon”, i.e. ”s-kaon”, is plotted. The basic feature is similar to the ”s-pion”. In Fig. 8, we show the results for the ”s-pion” when the spinor quark is replaced by a bosonic quark. As we extensively discussed in Section II, the subtraction of the end-point singularity is not required in the scalar quark case in contrast to the spinor quark case. In the scalar quark case, it is interesting to note that $`F_{\mathrm{val}}^{}`$ and $`F_{\mathrm{nv}}^{}`$ reveal a large difference compare to the spinor quark case, while $`F_{\mathrm{val}}^+`$ and $`F_{\mathrm{nv}}^+`$ are very similar to the spinor quark case. We find that the huge cancellation between $`F_{\mathrm{val}}^{}`$ and $`F_{\mathrm{nv}}^{}`$ observed in the spinor quark case does occur in the scalar quark case only for very strong binding, and the most of $`F_{\mathrm{tot}}^{}`$ is saturated by $`F_{\mathrm{nv}}^{}`$. However, the tiny contribution from $`F_{\mathrm{val}}^{}`$ near $`Q^2=0`$ grows as the binding gets weaker and we have demonstrated the dominance of the valence part in the small binding limit regardless of the spin content, as discussed above. Fig. 9 shows the corresponding results for the ”s-kaon” when the spinor quark is replaced by a bosonic quark. While the general features are similar to the spinor quark case, the bosonic quarks are more tightly bound together than the spinor quarks so that the charge radius of the meson is smaller than the case of spinor quarks as one might expect from the Pauli’s exclusion principle.
## IV Conclusion and Discussion
In this paper, we have analyzed both the plus and minus components of the current quantized on the light-front to compute the electromagnetic form factors of pseudoscalar and scalar mesons. We considered spin-1/2 consituents as well as spin-0 constituents and found dramatic differences between the two cases. Comparing with the covariant Feynman calculations, we notice that the common belief of equivalence between the manifestly covariant calculation and the LF calculation linked by the LF energy integration of the Feynman amplitude is not always realized. The minus component of the LF current generated by the fermion loop has a persistent end-point singularity that must be removed to assure covariance and current conservation. A similar singularity was observed in the calculation of the fermion self-energy in and . The plus component of the LF current, however, is immune to this disorder and provides a form factor identical to the one obtained doing the covariant Feynman calculation. This phenomenon is also associated with the spin-effect of the constituents because the calculation with the scalar(spin-0) constituents does not have the same symptom. Decomposing the LF amplitude into the wave-function and non-wave-function parts, it is interesting to note that the end-point singularity exists only in the non-wave-function vertex contribution.
Even after the singularity is removed, the minus component of the current sustains the zero-mode contribution while the plus component is free from the zero-mode. We have numerically estimated the importance of the non-wave-function vertices in all three cases that we discussed in Section II. We considered also the unequal constituent mass cases such as the kaon form factor. We find that the behaviors of $`F_{val}^{}`$ and $`F_{nv}^{}`$ are trimendously different between pseudoscalar and scalar meson cases, while $`F_{val}^+`$ and $`F_{nv}^+`$ have very similar features in both cases. The huge but remarkably exact cancellation between $`F_{val}^{}`$ and $`F_{nv}^{}`$ shown in the scalar meson case persists even if the spinor quark is replaced by the bosonic quark. In the bosonic quark case, however, the binding between the constituents is stronger than the spinor quark case. We also notice that the zero-mode $`F_{nv}^{}(0)`$ diminishes as the binding gets weaker. Our results are quite consistent to the earlier observation exhibiting the smaller zero-mode contribution in the heavier quark systems. In all of these cases, our results show that if the meson is weakly bound then the contributions from the wave-function and the non-wave-function vertices to the plus current are separately almost the same as those for the minus current. Of course, their sums add up to the same number as the covariant Feynman result in both the plus and minus cases.
The calculations carried out so far are semi-realistic as the model was 1+1-dimensional and only a point-vertex was considered. It is clear from a formal analysis of the 3+1-dimensional case, however, that a singularity of the same form will occur in the matrix element of $`J^{}`$ calculated in LFD regardless of dimensionality. A recent analysis of the Burkhardt-Cottingham sum-rule seems to reveal a similar divergence in the polarized spin-1/2 structure functions . While the additional regularization may be provided by smearing the point-vertex with a realistic wave-function in the 3+1-dimensional covariant treatment of the current, the identification of the singular term as we achieved in this work would still be necessary for the smeared vertex cases. The importance of the non-wave-function parts may nevertheless differ numerically from the 1+1-dimensional case. This point is presently under investigation.
###### Acknowledgements.
This work was supported in part by a grant from the US Department of Energy and the Netherlands Organisation for Scientific Research (NWO).
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# Realizing 3D Spectral Imaging in the Far-Infrared: FIFI LS
## 1 INTRODUCTION
Astronomers will soon have access to unprecedented spatial resolution and sensitivity in the far-infrared with the Stratospheric Observatory For Infrared Astronomy (SOFIA). Far-infrared astronomical observations, which are impossible from the ground due to water absorption, are necessary to understand fully a number of important astronomical problems and issues. Many astrophysical conditions require far-infrared probes because the areas of interest are mostly inaccessible at other wavelengths due to severe extinction from interstellar dust or the physics of interest is only manifest at far-infrared wavelengths. In particular, far-infrared spectroscopy, pioneered and developed on SOFIA’s predecessor the Kuiper Airborne Observatory (KAO) and greatly extended with the Infrared Space Observatory (ISO), will provide an injection of important data into astrophysical issues.
Building upon the success of our previous imaging Fabry-Perot far-infrared spectrometer (FIFI) that was designed for the KAO, we are developing a successor instrument for SOFIA: the Field-Imaging Far-Infrared Line Spectrometer (FIFI LS). FIFI LS will utilize integral field spectral imaging in two wavelength bands: 42 to 110 $`\mu `$m and 110 to 210 $`\mu `$m. This will allow the instrument to simultaneously obtain dual wavelength band, spectral imaging of a 5 $`\times `$ 5 pixel field, without scanning a Fabry-Perot or multiple pointings with a long-slit spectrometer. Thus, FIFI LS is designed as a “major-step” forward and will take advantage of the unique benefits that SOFIA offers.
With the increased sensitivity and resolution provided by SOFIA, the main scientific targets for FIFI LS will include the detailed morphological studies of: (1) the heating and cooling of galaxies, (2) star formation and the interstellar matter under low-metalicity conditions, as found in dwarf galaxies, (3) active galactic nuclei and their environment, (4) merging and interacting galaxies, and (5) large surveys of nearby galaxies. To reach our scientific goals, very high observing sensitivities and efficiencies are essential, requiring a compromise of spectral resolution; however, for the science objectives listed above, a comparably low spectral resolution (R $``$2000) is more than sufficient. Overall, FIFI LS on SOFIA will be more sensitive than the ISO Long-Wavelength Spectrometer and have much higher spatial resolution and mapping capabilities. As a future option, an extension of the instrument to the 25-42 $`\mu `$m range is planned upon availability of the Si:Sb detector arrays developed for SIRTF.
## 2 Instrument Design
### 2.1 Spectrometer Concept
FIFI LS achieves 2-dimensional spatial mapping and simultaneous spectral multiplexing by optically slicing, or re-arranging, the 2D field of view onto a single slit, which is then dispersed via a standard long-slit spectrometer. This type of “optical slicer” was originally devised for laboratory spectroscopy and later successfully implemented for near-infrared astronomy, but this is the first time that this technique has been applied in the far-infrared. Specifically, FIFI LS has a 5 $`\times `$ 5 pixel field of view which is sliced into five individual slitlets, re-arranged into a continuous 25 $`\times `$ 1 pixel slit, and finally, fed into a grating spectrometer and dispersed onto a 2D 16 $`\times `$ 25 pixel detector array.
One of the other unique aspects of FIFI LS is its dual-channel nature. By using a dichroic beamsplitter in the entrance optics and two optical slicers, FIFI LS covers the diffraction-limited field of view simultaneously in two wavelength bands: a short wavelength band 42 to 110 $`\mu `$m and a long wavelength band 110 to 210 $`\mu `$m. Fig. 1 shows how the dual channel integral field system is realized in more detail as a projection of the focal plane onto the two detector arrays with diffraction-limited optics. On the right side of the figure is the long wavelength channel of the spectrometer which has a pixel scale of 14<sup>′′</sup> per pixel. On the left side of the figure is the short wavelength channel which has a pixel scale of 7<sup>′′</sup> per pixel. This of course implies that the shorter wavelength channel has a smaller field of view, as demonstrated in Fig. 1. Overall, this scheme ensures that for all spatial elements in the field, spectra are observed simultaneously in the two bands, thereby increasing observing efficency.
### 2.2 Optical Design
The geometrical layout, optimization, and analysis of the entire spectrometer optics was carried out with the optical design software ZEEMAX-EE Ver. 8.0. The two spectrometers were designed to yield diffraction-limited image quality: the geometrical spot diameters are small compared to the diffraction disk diameter of a point source imaged by the telescope. A 3-D solid model of the complete FIFI LS optics is shown in Fig. 2, viewed from the top; all optical components are cooled to cryogenic temperatures.
A rotating K-mirror assembly (hidden on the left in Fig. 2) at liquid nitrogen temperature compensates for rotation of the field of view during long integration times. Two chopped calibrator sources with temperatures close to liquid nitrogen can be switched into the beam for internal calibration and flat-fielding at irradiation levels close to that of the SOFIA telescope background. The calibrator sources are located near a pupil image so that the light paths are equivalent for internal calibration and observation. All optics after the calibration system are at liquid helium temperature. After the calibration optics, the beam enters the spectrograph through a cold Lyot stop (not shown in Fig. 2) to suppress light diffracted by the entrance aperture. A far-infrared dichroic beam splitter feeds both the long and short wavelength spectrometer. In the short wavelength branch, re-imaging optics double the image size on the image slicer.
In Figure 3, the optical slicer is shown in more detail. The light enters into the slicer system from the flat mirror (upper-left in Fig. 3), then continues down toward the slicer mirrors (the cube-like object near the left bottom of Fig. 3). These 5 mirrors in the mirror stack form the image slicer, which acts as a field mirror creating a pupil image on one of the 5 re-imaging capture mirrors (the five mirrors near the top-middle). The capture mirrors re-arrange the images of the 5 slices along one, slightly curved line on the slit-mirror, near the center of the image. Working in combination, these three mirrors, the slicer, capture, and slit mirrors, perform the slicing and re-arranging of the field of view onto the entrance slit for a grating spectrograph. In addition, the three mirror system re-aligns the pupil of each slit, so that the virtual pupils of each slice coincide. In this case, we do not use simple flat mirrors for the slicing process, as is usually done in near-infrared instruments, because of the larger A$`\times \mathrm{\Omega }`$ product (area times the beam solid angle); curved mirrors allow for a much more compact slicer assembly design.
The FIFI LS spectrometers use slightly off-axis (0.5) Littrow mounted gratings (the long rectangular optics in Fig. 2); a truly Littrow mounted grating spectrometer would have the same entrance and exit optical paths, but in the FIFI LS system the entrance and exit paths are slightly separated so that the outgoing beam can easily be re-imaged onto the detectors. This arrangement allows for the dual use of the two anamorphic collimators (the two large mirrors right-most and the two large mirrors on top and bottom in Fig. 2). From the image slicer, the beam continues in each channel via the anamorphic optics, which expand the beam to an elliptical shaped cross-section that illuminates the grating over its 300 mm length, attaining the proposed spectral resolution of about 1000-2000 ($`\mathrm{\Delta }\nu `$ 100 - 250 km/s) In each return path, the anamorphic re-imaging is used to match the spatial and spectral resolution of the system to the square pixels of the detector.
In order to cover the wide wavelength range necessary, the Littrow-mounted gratings are operated in the 1st order (long wavelengths) or 1st and 2nd orders (short wavelengths). To select first or second diffraction order, exchangeable filters are utilized in the short-wavelength branch of the spectrometer. The observing wavelength of each spectrograph is then tuned by tilting the grating, and exchangeable filters are used to select the appropriate working order of the grating.
Since the grating efficiency strongly depends upon wavelength and diffraction order, the groove profile and separation for each grating were optimized separately. Rigorous vectorial diffraction analysis of the grating efficiency in the actual Littrow configuration were performed using the PCGrate-1E Ver. 3.0 software. The calculated grating efficiencies as a function of wavelength and diffraction order are shown in Fig. 4. Our goal was to maximize the efficiencies across the two bands. For the long wavelengths, this was done by using a symmetric grating profile that provided a high efficiency across the band, and for the shorter wavelengths one grating order was insufficient, and the best result was obtained by maximizing the efficency over two orders by using an asymmetric grating profile. In general, the optical parameters of FIFI LS are summarized in Table 1.
### 2.3 Grating Mechanical Layout
To cover the specified wavelength range of 42 to 110 $`\mu `$m and 110 to 210 $`\mu `$m in first and second order, both gratings have to be tilted by an angle of $`\pm `$20. However, to reach a spectral accuracy of $``$100 - 250 km/s, the grating has to be moved and controlled with a precision of less than 4<sup>′′</sup> while maintaining mechanical stability of the optical surface. In addition, this precision must be reached at liquid helium temperatures. In order to minimize deformation to the grating surface from vibration, the grating structure design was extensively tested with finite element modeling. The final grating structure, as shown in Fig. 5, has $`<`$ 1 $`\mu `$m of deflection with 10 N of applied force at the outer edge of the structure.
To obtain a large tilting range and precise positioning, the grating is actuated by a two stage tilting mechanism. The first stage, for coarse positioning, consists of a support structure connected to the bottom of the grating. The support structure is driven by a roller-screw lever arm mechanism (a so-called sine-bar mechanism) that tilts the grating via a push-pull movement. The roller screw is driven by a stepper motor at liquid nitrogen temperature that is thermally isolated from the grating by a magnetic feed-thru. The second stage, for fine positioning, utilizes a PZT, which drives the grating with respect to the support structure via a directly attached lever arm. The PZT movement is controlled in a software loop with closed-loop bandwidth up to 1kHz. Additionally, an eddy current damping system is mounted to the grating to minimize in-flight vibrations of the airplane, which are a perpetual source of error on airborne experiments, especially affecting moving parts where resonant modes can raise error motion to an unacceptable value.
### 2.4 Grating Position Read-out
For highly reliable measurement of the angular position of the gratings, we directly attached an INDUCTOSYN<sup>TM</sup> position transducer. An INDUCTOSYN<sup>TM</sup> is effectively a transformer with the primary and secondary windings placed on a rotor and stator, respectively. The winding pattern on the rotor is excited by a 10 kHz signal, while on the stator, there are two periodic patterns that are 90 out of phase with each other. In operation, the rotor and stator windings inductively couple such that the two output signals from the stator have amplitudes which vary as the sine and cosine functions based on the relative position in each winding cycle. By comparing the two amplitudes, a high-resolution difference positional accuracy is obtained.
The position transducers work with a fairly high operating current of $``$ 0.25 Amperes. To provide thermal insulation of the 4K worksurface, only thin signal wires are tolerated within the cryogenic regions. Thin wires, on the other hand, raise the power dissipation, leading again to higher heat input. To overcome this problem, we use a superconducting transformer which steps down the signal voltage by a factor of 50 after entering the 4K region. In this way, a high-voltage, but low- current excitation signal, can be conveyed via thin wires, reducing dissipation losses.
The relative amplitudes of the two transducer outputs are a measure of the actual position of the grating, with proper initial calibration. Both signals are amplified separately in a low-noise amplifier stage and translated into 14 bit digital position data by a Resolver-to-Digital Converter. Since the transducer output produces a roll over of the position data every 1.4, an additional loop counter keeps track on the position data over the full 40 tilting range of the grating. First tests with a prototype setup at liquid helium temperature, showed that an angular resolution of less than 0.5<sup>′′</sup> can be obtained easily.
### 2.5 Detectors
As mentioned above, FIFI LS uses two 16 $`\times `$ 25 detector arrays to cover the 42 - 110 $`\mu `$m and 110 - 210 $`\mu `$m wavelength bands. We chose to use Gallium-doped Germanium photoconductor detectors since they are proven to be very sensitive in the wavelength range 40 - 120 $`\mu `$m, and, with the application of $``$ 600 N mm<sup>-2</sup> of stress, their wavelength sensitivity shifts to 100 - 220 $`\mu `$m. Thus, FIFI LS uses two Ge:Ga detector arrays, one stressed and one unstressed; our design and initial testing is discussed in detail in this volume.
## 3 CRYOSTAT
The FIFI LS instrument is directly attached to the science instrument mounting flange of the SOFIA telescope. In Fig. 6, a 3D rendered model of the FIFI LS instrument is shown. The cryostat is enclosed by a vacuum vessel which also provides the mechanical interface for mounting. Mounted to the bottom of the vacuum vessel are the dichroic beam splitter and the field optics for the focal plane guiding camera. The dichroic filter separates the telescope’s light (entering Fig. 6 on the bottom right) into two beams: the optical, directed downward to the telescope guiding camera, and the infrared, directed upward into the cryostat. The infrared beam from the telescope enters the vacuum vessel through a polyethylene window, which also serves as a pressure barrier between stratospheric pressure and the vacuum inside the instrument.
There are three cryogenic containers in FIFI LS for liquid nitrogen, liquid helium, and superfluid helium. The manner in which the three vessels are mounted together is shown in Fig. 7, a center cut-thru the cryogenic vessels. The liquid nitrogen container has a capacity of 25 liters, which provides cooling for the outer radiation shields, the liquid nitrogen worksurface and the entrance optics (i.e. the K-mirror assembly and the re-imaging optics). The cut-thru of the liquid nitrogen vessel in Fig. 7 shows two of its inner support ribs. The liquid nitrogen worksurface is suspended from the warm vacuum vessel by G-10 fiberglass stand-offs which provide high mechanical stiffness and low thermal conductivity. The liquid nitrogen container is designed for a cryogen holding time of about 28 hours.
The 35 liter main liquid helium reservoir provides cooling for the inner radiation shields and the liquid helium optical bench. The liquid helium optical bench, suspended from the liquid nitrogen worksurface by G-10 tabs, mechanically supports and cools all of the optical components after the calibration system (excluding the detectors). The expected cryogen holding time for the main liquid helium reservoir is up to 50 hours.
Since the detector arrays require operating temperatures below 4K, they are mounted to a small 2.8 superfluid helium tank, which is again suspended by G-10 tabs from the liquid helium optical bench. This tank is pumped in order to reach a temperature of $``$ 2K. The expected maximum holding time for the liquid helium tank is about 18 hours. To ensure safe operation during flight, all cryogen vessels are provided with coaxial neck tubes and warm pressure relief valves.
### 3.1 Cryostat Analysis
In order to comply with Federal Aviation Administration requirements, we are designing the Cryostat to withstand 3 times the operational pressure— about 3.5 bars. Fig. 8 shows finite element analysis results for the three large vessels at 3.5 bars of internal pressure– overemphasizing the deformation for display purposes The maximum deformation in the three vessels are 0.91 mm for the vacuum vessel, 0.20 mm for the LN<sub>2</sub> vessel, and 0.24 mm for the LHe vessel. In addition, we used FEA for determining the best design for stability of the internal optical work surfaces and to minimize the amount of stress on the welding joints of the vessels. In the latter case, for example, we find that the best approach is to create a nearby zone that can elastically deform, thereby relocating areas of high stress away from the welding seam.
## 4 Scientific Capabilities
FIFI LS employs two fixed pixel sizes of 7<sup>′′</sup> (short wavelength spectrometer) and 14<sup>′′</sup> (long wavelength spectrometer), respectively, determined by the image slicer. The 5 $`\times `$ 5 pixel fields of view are observed simultaneously with two Ge:Ga photoconductor arrays. Observing wavelength are adjusted by tilting the Littrow mounted grating in each channel. Spectral coverage of $``$ 1500 km/s around a selected far-infrared line is obtained simultaneously for all 25 spatial pixels. A summary of important instrument properties is shown in Table 2.
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# Thermal/Electronic Transport Properties and Two-Phase Mixtures in La5/8-xPrxCa3/8MnO3
\[
## Abstract
We measured thermal conductivity $`\kappa `$, thermoelectric power S, and dc electric conductivity $`\sigma `$ of La<sub>5/8-x</sub>Pr<sub>x</sub>Ca<sub>3/8</sub>MnO<sub>3</sub>, showing an intricate interplay between metallic ferromagnetism (FM) and charge ordering (CO) instability. The change of $`\kappa `$, S and $`\sigma `$ with temperature (T) and x agrees well with the effective medium theories for binary metal-insulator mixtures. This agreement clearly demonstrates that with the variation of T as well as x, the relative volumes of FM and CO phases drastically change and percolative metal-insulator transition occurs in the mixture of FM and CO domains.
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Immense resurgent activities on mixed-valent manganites reveal the importance of Jahn-Teller-type electron/lattice coupling in addition to the double exchange mechanism. Another important aspect of manganites is the existence of various-scale, real-space variation of physical properties or parameters. For example, even though the orthorhombicity of the average structure drastically decreases with the replacement of La by divalent ions in LaMnO<sub>3</sub>, the large variation of local Mn-O bond lengths still remains intact. In high divalent-ion doping ranges, charge ordering results in sheet-like arrangements of Mn<sup>3+</sup> and Mn<sup>4+</sup> ions with 10-30 Å length-scales . Furthermore, ferromagnetic resonance experiments have shown the presence of two types of signals in ferromagnetic manganites, which was interpreted as evidence of electronic phase separation . Various experiments suggest the existence of magnetic polarons or mobile ferromagnetic clusters at high temperature (T), which could be viewed as resulting from dynamic phase separation . Recently, an electron diffraction study on low Curie temperature (T$`_\text{C}`$) manganites has revealed that there coexist ferromagnetic (FM)-metallic and CE-type charge-ordered (CO), insulating domains . It has been emphasized that this particular type of static phase separation is responsible for colossal magnetoresistance in low T$`_\text{C}`$ manganites. Various theoretical models for mixed-valent manganites also reveal the general tendency of static or dynamic electronic phase separation .
The transport properties of metal-insulator (M-I) mixtures have been perennial topics for both theoretical and experimental condensed matter physics . Most of the experiments were performed on films with deposited M-I mixture or artificial bulk M-I composites prepared under pressure . The total electric and thermal conductivity and thermoelectric power of binary M-I mixtures were successfully explained by the effective medium theories .
In this letter, we report the absolute values of the magnetization M, thermal conductivity $`\kappa `$, thermoelectric power S, and dc electric conductivity $`\sigma `$ of La<sub>5/8-x</sub>Pr<sub>x</sub>Ca<sub>3/8</sub>MnO<sub>3</sub> with various x and T. Various aspects of our results are consistent with the coexistence of FM and CO phases, whose relative volumes change with both T and x, and the percolative M-I transition in FM-CO mixtures. Furthermore, the T and x dependence of $`\sigma `$, $`\kappa `$, and S agrees well with the (general) effective medium theories for M-I mixtures.
High-quality polycrystalline specimens of La<sub>5/8-x</sub>Pr<sub>x</sub>Ca<sub>3/8</sub>MnO<sub>3</sub> with x=0.0, 0.1, 0.2, 0.25, 0.3, 0.35, 0.375, 0.40, 0.42, and 0.625 have been prepared with the standard solid state reaction. We fixed the Ca concentration at 3/8 because our previous studies showed that T$`_\text{C}`$ is optimized at this particular Ca doping level . $`\sigma `$ of all specimens with accurate geometry was measured with the standard four probe method, and M was measured with a SQUID magnetometer. Both $`\kappa `$ and $`S`$ of the representative samples (x=0.0, 0.1, 0.25, 0.35, 0.375, 0.42, and 0.625) have been measured from 8 to 310 K with the steady state method. A radiation shield was used to obtain absolute $`\kappa `$ values.
The systematic T-dependent and M/H curves are shown in Fig. 1. M/H curves were measured in H = 2 kOe, which was carefully chosen to align FM domains without influencing the CO insulating phase. These results indicate that La<sub>5/8-x</sub>Pr<sub>x</sub>Ca<sub>3/8</sub>MnO<sub>3</sub> (x=0.0) is, basically, FM-metallic below 275 K, and that the ground state of Pr<sub>5/8</sub>Ca<sub>3/8</sub>MnO<sub>3</sub> (x=0.625) is CO-insulating below 225 K. The behaviors of $`\sigma `$ and M/H for other x compositions are systematically in-between those for x=0.0 and 0.625. Open circles (Fig. 1 (a)) represent the M-I transition points where $`\sigma `$ becomes (local) minimum . Open circles in Fig. 1 (b) show the M/H values at the M-I transition points determined from the curves. The average of those M/H values (dotted line) is 17$`\pm `$2 % of 8.1 emu/mol, the saturated M/H value of x=0.0. Thus, with changing T, the M-I transition occurs when M of each sample becomes about 17 % of that of x=0.0, independent from x value. If we assume that the $`T`$-dependent volume fraction f($`T`$) of the FM domain is proportional to M($`T`$) in H=2 kOe, the M-I transition with changing $`T`$ occurs when f reaches $``$0.17, close to the three-dimensional percolation threshold (f$`_\text{c}`$). This behavior can be also seen in the variation of $`\sigma `$ with x for fixed $`T`$. In Fig. 2 (a), $`\sigma `$ vs. M<sub>x</sub>($`T`$)/M$`_{\text{0.0}}`$($`T`$) plot is shown at 10 and 100 K for various x. Interestingly, the samples with M<sub>x</sub>($`T`$)/M$`_{\text{0.0}}`$($`T`$) $`<`$ 0.15 are in the insulating regions of the curves in Fig. 1, while those with M<sub>x</sub>($`T`$)/M$`_{\text{0.0}}`$($`T`$) $`>`$ 0.17 are in the metallic regimes. Therefore, with variation of x, the M-I transition occurs when M<sub>x</sub>/M$`{}_{0.0}{}^{}`$0.15-0.17. These observation clearly suggest that the M-I transition with both $`T`$ and x takes place when f($`T`$,x) ($``$M<sub>x</sub>($`T`$)/M$`_{\text{0.0}}`$($`T`$)) becomes close to f$`_\text{c}`$. It is important to note that M<sub>0.0</sub>($`T`$) contains the natural $`T`$-dependence of ferromagnetic moment so that f($`T`$,x) equals to M<sub>x</sub>($`T`$)/M$`_{\text{0.0}}`$($`T`$) (not M<sub>x</sub>($`T`$)/M$`_{\text{0.0}}`$($`T`$=0)).
To gain further insights into the nature of the M-I transition, we measured $`T`$-dependent $`\kappa `$ and $`S`$ as shown in Fig. 3. First, we note that the estimated $`\kappa _e`$ from $`\sigma `$, using the Wiedemann-Franz law, is by two orders of magnitude smaller than the measured $`\kappa `$($`T`$) for all x, indicating the dominant phonon contribution. Furthermore, at high $`T`$ above $`T__\text{C}`$ or $`T_{_{\text{CO}}}`$, $`\kappa `$ always increases when $`T`$ is raised, and the magnitude of $`\kappa `$ is in the range of 0.5-2 W/mK, comparable to that of amorphous solids . This behavior has been attributed to local anharmonic lattice distortions associated with small polarons . Related to this, $`S`$($`T`$) of all the samples at high $`T`$ follows the form S<sub>0</sub>+E<sub>g</sub>/k$`{}_{B}{}^{}T`$ with the gap energy E<sub>g</sub> systematically increasing from 4 meV (x=0.0) to 12 meV (x=0.625). These E<sub>g</sub> values are significantly smaller than the activation energies (125 meV: x=0.0 to 175 meV: x=0.625) associated with $`\sigma `$($`T`$). This difference can result from the small polaronic transport .
As shown in Fig. 3 (a), with increasing x, $`\kappa `$($`T`$) smoothly evolves from that of x=0.0 to that of x=0.625, and the $`\kappa `$ increase at $`T_\text{C}`$ becomes smaller. Consistently, $`\kappa `$ vs. M<sub>x</sub>($`T`$)/M$`_{\text{0.0}}`$($`T`$) at 10 and 100 K (Fig. 2 (b)) shows that $`\kappa `$ varies monotonically from the maximum ($`\kappa `$ of x=0.0) to the minimum ($`\kappa `$ of x=0.625). For x=0.0, $`\kappa `$ at $`T_\text{C}`$ increases suddenly, probably due to the suppression of local lattice distortions associated with small polarons. For x=0.625, $`\kappa `$($`T`$) behaves as in amorphous solids in the entire $`T`$ range, except for a slight increase at $`T_{\text{co}}`$.
In comparison with $`\kappa `$ and $`\sigma `$, $`S`$ exhibits seemingly different behaviors with x. $`S`$ is very close to the metallic value, $`S`$ of x=0.0, even near f$`_\text{c}`$ where $`\sigma `$ (or $`\kappa `$) is still significantly smaller than $`\sigma `$ (or $`\kappa `$) of x=0.0. $`S`$ vs. M<sub>x</sub>/M<sub>0.0</sub> at 100 K (Fig. 2 (c)) clearly demonstrates this tendency; $`S`$ is close to zero (slightly negative), and is insensitive to M<sub>x</sub>/M$`_{\text{0.0}}`$ as long as M<sub>x</sub>/M$`_{\text{0.0}}`$ $``$10 %. In contrast, $`\sigma `$ for M<sub>x</sub>/M$`_{\text{0.0}}`$ = 10 % at 100 K is more than three orders of magnitude smaller than that of x=0.0. In fact, $`T`$-dependence of $`S`$ near $`T_\text{C}`$ is also consistent with this metallic $`S`$ behavior near f$`_\text{c}`$. With decreasing $`T`$ near $`T_\text{C}`$, $`S`$ starts to decrease, i. e., becomes metallic at $`T`$ higher than those for $`\kappa `$ and $`\sigma `$ changes. For example, in the heating curves of x=0.35, $`S`$ starts to decrease around 130 K, significantly higher than $`T`$ (100 K) for abrupt $`\kappa `$ increase or $`T`$ ($``$110 K) for $`\sigma `$ minimum.
When thermal/electronic transport properties of our systems are viewed as those of M-I mixtures, the above peculiar $`S`$ behavior is, in fact, consistent with the theoretical prediction of effective thermoelectric power $`S_\text{E}`$ by Bergmann and Levy . For an isotropic binary mixture, they showed that in terms of $`\sigma `$, $`\kappa `$, and $`S`$ of each component,$`S_\text{E}`$ is given by
$$S_\text{E}=S_\text{M}+(S_\text{I}S_\text{M})\left(\frac{\kappa _\text{E}/\kappa _\text{M}}{\sigma _\text{E}/\sigma _\text{M}}1\right)/\left(\frac{\kappa _\text{I}/\kappa _\text{M}}{\sigma _\text{I}/\sigma _\text{M}}1\right)\text{ },$$
(1)
where the subscripts M and I refer to metallic and insulating components, respectively. $`\kappa _\text{E}`$ and $`\sigma _\text{E}`$ refer to effective thermal and electric conductivity, respectively, of the binary mixture. This equation has been successfully applied to explain $`S`$ behaviors of binary Al-Ge films . When $`\sigma _\text{I}`$/$`\sigma _\text{M}`$ $`<`$$`<`$ $`\kappa _\text{I}`$ /$`\kappa _\text{M}<`$1 (which applies to our system) and for f = f$`_\text{c}`$ , the above equation leads to$`S_\text{E}S_\text{M}`$, which explains our experimental results noted above.
To quantitatively compare our results with Eq. (1), we calculated$`S_\text{E}`$ at every $`T`$ . In this comparison, experimental $`\sigma `$ and $`\kappa `$, shown in Figs. 1 and 3, were used for $`\sigma _\text{E}`$ and $`\kappa _\text{E}`$ , and ($`\sigma _\text{I}`$, $`\kappa _\text{I}`$, and $`S_\text{I}`$) and ($`\sigma _\text{M}`$, $`\kappa _\text{M}`$, and$`S_\text{M}`$) are assumed to be identical with those of x=0.0 and 0.625, respectively. (For large x, it was difficult to measure $`S`$ at low $`T`$ due to high resistivity, so we assumed that $`S_\text{I}`$ changes as 1/$`T`$, and that $`\sigma _\text{I}`$ does exponentially.) The solid lines in Fig. 3 (b) depict the calculated$`S_\text{E}`$, using Eq. (1), for heating curves of x=0.25, 0.35, and 0.42. The calculated curves match with our experimental $`S`$ surprisingly well at all $`T`$ below $`T_\text{C}`$ or $`T_{\text{co}}`$.
For $`\sigma _\text{E}`$ (or $`\kappa _\text{E}`$) of a binary M-I mixture, Mclachlan proposed the general effective medium (GEM) equation,
$$(1f)\left(\frac{\sigma _I^{1/t}\sigma _E^{1/t}}{\sigma _I^{1/t}+A\sigma _E^{1/t}}\right)+f\left(\frac{\sigma _M^{1/t}\sigma _E^{1/t}}{\sigma _M^{1/t}+A\sigma _E^{1/t}}\right)=0\text{ },$$
(2)
where $`A`$ =(1-$`f_\text{c}`$)/$`f_\text{c}`$. The same equation also works for $`\kappa `$. The critical exponent $`t`$ is close to 2 in three dimension. This equation has been successfully applied to isotropic inhomogeneous media in wide f regions including percolation regime .
To apply the GEM equation to $`\sigma `$($`T`$), we assumed that f($`T`$,x) = M<sub>x</sub>($`T`$)/M$`_{\text{0.0}}`$($`T`$), $`\sigma _\text{M}`$($`T`$) = $`\sigma `$($`T`$) of x=0.0, and $`\sigma _\text{I}`$($`T`$) = $`\sigma `$($`T`$) of x=0.625. With the parameters $`t`$=2 & f$`_\text{c}`$=0.17, the calculated$`\sigma _\text{E}`$ for various x are shown as solid lines in Fig. 4. At $`T`$ $`>`$ $``$80 K,$`\sigma _\text{E}`$($`T`$) nicely matches the experimental $`\sigma `$($`T`$) even if $`\sigma `$ changes by 6 orders of magnitude with $`T`$ and x. However, this agreement does not hold at very low $`T`$. The calculated$`\sigma _\text{E}`$($`T`$) at $`T`$ $`<`$ 80 K with the same parameters $`t`$=2 & f$`_\text{c}`$=0.17 significantly deviated from the experimental $`\sigma `$($`T`$). We found that at $`T`$ $`<`$ 80 K, the calculated$`\sigma _\text{E}`$($`T`$) matches the experimental $`\sigma `$($`T`$) better when $`t`$ is increased to $``$ 4. This is more evident in the x dependence of $`\sigma `$ at 10 and 100 K, as shown in Fig. 2 (a).$`\sigma _\text{E}`$(100 K), calculated with $`t`$=2 & f$`_\text{c}`$=0.17 (solid line), matches the experimental values (open circles) better than that with $`t`$=4 & f$`_\text{c}`$=0.17 (dotted line). However,$`\sigma _\text{E}`$(10 K), calculated with $`t`$=4 & f$`_\text{c}`$=0.15 (solid line), is closer to the experimental $`\sigma `$ (solid circles) than that with $`t`$=2 & f$`_\text{c}`$=0.15 (dashed line) . (The change of f$`_\text{c}`$ in the range of 0.15-0.17 makes little difference.) These observations demonstrate that $`t`$, normally close to the three dimensional exponent of 2, becomes $``$4 at very low $`T`$. A similar, drastic increase of $`t`$ has been noted in the case of tunneling transport for M-I mixtures , suggesting that the tunneling process between FM domains is important for $`\sigma `$ of our system at very low $`T`$ .
By using the GEM equation for $`\kappa _\text{E}`$, the $`\kappa `$($`T`$) for various x can be calculated with the assumption that $`\kappa _\text{I}`$($`T`$)=$`\kappa `$($`T`$) of x=0.625, $`\kappa _\text{M}`$($`T`$)=$`\kappa `$($`T`$) of x=0.0, $`t`$=2, and f$`_\text{c}`$=0.17. The solid lines in Fig. 3 (a) represent the estimated $`\kappa _\text{E}`$ for x=0.1, 0.25, 0.35, and 0.42. In addition, the calculated $`\kappa _\text{E}`$ as a function of M<sub>x</sub>/M$`_{\text{0.0}}`$ at 10 and 100 K is depicted as solid lines in Fig. 2 (b). Estimated $`\kappa _\text{E}`$ lines in Figs. 2 and 3 coincide with the experimental data well . In order to confirm self-consistency,$`S_\text{E}`$ at 10 and 100 K is evaluated by using the calculated$`\sigma _\text{E}`$ and $`\kappa _\text{E}`$ (solid lines of Fig. 2 (a) and (b)), and the Eq. (1). The calculated$`S_\text{E}`$ (solid lines of Fig. 2 (c)) with the variation of x is in good agreement with the experimental values.
The unambiguous agreement between the measured thermal/electronic transport properties and the calculated values based on Eqs. (1) and (2) strongly indicates that: (1) transport properties are dominated by thermal/electrical conduction in M-I mixtures, (2) the relative volume of the (FM) metallic phase is proportional to the measured M($`T`$,x), and (3) the $`T`$-dependent transport and magnetic properties of metallic and insulating phases are always that of x=0 and 5/8, respectively. Combined with the earlier electron diffraction results, this successful agreement demonstrates that all of the thermal, electronic, and magneto-transport properties of La<sub>5/8-x</sub>Pr<sub>x</sub>Ca<sub>3/8</sub>MnO<sub>3</sub> are dominated by the percolative conduction through FM-metallic domains which is statically mixed with CO insulating domains. One surprising indication from our results is that at least in low $`T_C`$ materials (x $`>`$ 0.25), the so-called Curie transition is, in fact, the M-I transition across percolation threshold, and the ordered FM moment changes smoothly near the percolative phase transition $`T`$.
We greatly thank A. J. Millis, G. Kotliar, E. Abrahams, and T. W. Noh for useful discussions. We are partially supported by the NSF-DMR-9802513. K. H. Kim and M. U. are partially supported by the KOSEF and by the JPSJ Fellowship, respectively.
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# Harmonic Splittings of Surfaces
## 1 Introduction
Let $`\mathrm{\Gamma }=\pi _1(S)`$ be the fundamental group of a closed surface $`S`$ of genus at least two. Morgan-Shalen showed \[MS2\], \[GiSh\] that every point in the Thurston compactification $`𝒫(S)`$ of the Teichmuller space $`\text{Teich}(S)`$ gives an isometric $`\mathrm{\Gamma }`$-action on an R-tree. Given a measured foliation $`𝒫(S)`$, the action is simply the $`\mathrm{\Gamma }`$-action on the leaf space of the lift of $``$ to $`\text{H}^2`$. This action is small in the sense that edge stabilizers do not contain rank two free groups. It is also minimal in the sense that it leaves no proper subtree invariant.
Shalen \[Sh\] conjectured that every minimal small action of $`\mathrm{\Gamma }`$ on an R-tree $`T`$ arises in this way. This problem has several important applications in low-dimensional geometry and topology (see \[Ot\]). Partial results were obtained by Morgan-Shalen \[MS1\] and Gillet-Shalen \[GiSh\].
The conjecture was eventually proven in two steps: Morgan-Otal \[MO\] (see also \[Ha\]) constructed the candidate foliation, with dual R-tree $`R`$, and a $`\mathrm{\Gamma }`$-equivariant morphism $`\varphi :RT`$ so that $`\varphi `$ has no “edge folds” (see below); then Skora \[Sk1, Sk2\] showed that $`\varphi `$ has no “vertex folds”, giving that $`\varphi `$ is a $`\mathrm{\Gamma }`$-equivariant isometry, completing an affirmative solution to the conjecture.
###### Theorem 1.1 (Morgan-Otal, Skora)
Let $`\mathrm{\Gamma }=\pi _1(S)`$, $`S`$ a closed surface of genus at least two. Then any small, minimal $`\mathrm{\Gamma }`$-action on an R-tree is dual to the lift of a measured foliation on $`S`$.
A complete exposition of Theorem 1.1 is given in \[Ot\].
The purpose of the present paper is to prove Theorem 1.1 from a different point of view, using harmonic maps. Harmonic maps were used by Gromov-Schoen \[GrSc\] to show that certain groups do not act nontrivially on singular spaces such as trees. Here we use harmonic maps to classify, in the special case of a surface group, all minimal, small actions on R-trees, against a background where many such actions exist (namely $`6g7`$ dimensions worth).
Our other interest in this proof is in the way it uses harmonic maps as a tool in combinatorial group theory. For example, combinatorial topology arguments become greatly simplified (via the maximum principle) when looking at a harmonic representative. Another example is the existence of a moduli space of harmonic maps (and harmonic map invariants) associated to a group action, allowing for an extra tool in proofs.
### 1.1 Outline
Here is a brief description of our approach to the proof.
Step 1 (Find a harmonic map): Given a small action of the surface group $`\mathrm{\Gamma }=\pi _1(S)`$ on an R-tree $`T`$, it is relatively straightforward to find a $`\mathrm{\Gamma }`$-equivariant harmonic map $`f:\stackrel{~}{S}T`$. Here we have endowed $`S`$ with a complex structure.
Step 2 (Associated data): The harmonic map $`f`$ automatically has associated to it the following data:
* a $`\mathrm{\Gamma }`$-equivariant holomorphic quadratic (Hopf) differential $`\stackrel{~}{\mathrm{\Phi }}`$ on the Riemann surface $`\stackrel{~}{S}`$
* a $`\mathrm{\Gamma }`$-equivariant measured foliation $`\stackrel{~}{}`$, the vertical foliation of $`\stackrel{~}{\mathrm{\Phi }}`$
* the leaf space $`R`$ of $`\stackrel{~}{}`$, with metric induced from the measure on $`\stackrel{~}{}`$, making $`R`$ into an R-tree
The map $`f`$ is projection along the leaves of $`\stackrel{~}{}`$, with the possibility of several vertical leaves being sent to the same point in $`T`$. The $`\mathrm{\Gamma }`$-action on $`S`$ induces an isometric $`\mathrm{\Gamma }`$-action on $`R`$.
Step 3 (Morphism from a geometric action to the given action): Let $`\pi :\stackrel{~}{S}R`$ be the natural projection sending each leaf of $`\stackrel{~}{F}`$ to a point. Here a leaf may have a countable number of k-pronged singularities. We then obtain a $`\mathrm{\Gamma }`$-equivariant morphism $`\varphi :RT`$ of R-trees, where $`\varphi =f\pi ^1`$. We must show that $`\varphi `$ is an isometry, which is the same as saying that $`\varphi `$ does not fold at any point.
Step 4a (No edge folds): If $`\varphi `$ folded at an edge point of $`R`$, i.e. a point whose “tangent space” has only two directions, then this would force $`f`$ to locally take the form $`z|\text{Re}z|`$ which is not harmonic. Hence there are no edge folds, nor even vertex folds at trivalent vertices. The vertex points of $`R`$ are precisely the images under $`\pi `$ of leaves of $`\stackrel{~}{}`$ passing through a singularity of $`\stackrel{~}{}`$.
Step 4b (No vertex folds): Ruling out folds at high order vertices $`vR`$ requires a global argument (see Example 3.2.1 of a local vertex fold). The smallness hypothesis implies that, if two edges adjacent to $`v`$ are folded together, then neither edge can contain a point representing the lift of a closed leaf of $``$. This basically allows us to reduce the proof to the model case (see §5.2.3) where some leaf of $``$ is dense in $`S`$.
We now exploit the fact that we have a choice of conformal structures for $`S`$. Assuming $`\varphi `$ folds at some vertex point, we can always choose a path of conformal structures $`S_t`$ on $`S=S_0`$ so that the Hubbard-Masur differential on $`\stackrel{~}{S}_t`$ (the holomorphic differential $`\stackrel{~}{\mathrm{\Psi }}_t`$ whose vertical foliation projects to $`R`$) has simple zeroes for $`t0`$, and the edges which are folded together are represented on $`\stackrel{~}{}_t`$ by domains with a common one-dimensional frontier. As the harmonic map would again take the form $`z|\text{Re}z|`$ across this frontier, we see that $`\stackrel{~}{\mathrm{\Psi }}_t\stackrel{~}{\mathrm{\Phi }}_t`$, where $`\stackrel{~}{\mathrm{\Phi }}_t`$ is the Hopf differential for the harmonic map $`f:\stackrel{~}{S}T`$.
Hence there is a family of distinct R-trees $`R_t`$ and morphisms $`\varphi _t:R_tT`$. These trees come from measured foliations on $`S`$ which themselves come from interval exchange maps. But any nontrivial continuous variation in an interval exchange map forces a nontrivial variation in the tree $`T`$. As $`T`$ is fixed this is impossible, so there can be no vertex folds.
### 1.2 Acknowledgements
We thank Mladen Bestvina for useful discussions, Howard Masur for all his help (including the idea and most of the details of the proof of Proposition 2.3), and the referee for numerous comments and corrections which greatly improved the paper. Misha Kapovich \[Ka\] also had the idea to use harmonic maps in the proof of Skora’s theorem.
## 2 Preliminaries
### 2.1 R-trees
An R-tree is a metric space $`T`$ such that any two points in $`T`$ are joined by a unique arc and every arc is isometric to a segment in R. Let $`[x,y]`$ denote the unique (geodesic) segment from $`x`$ to $`y`$ in $`T`$.
We say that $`xT`$ is an edge point (resp. vertex point) of $`T`$ if $`T\{x\}`$ has precisely two (resp. more than two) components. An edge in $`T`$ is a nontrivial, embedded segment $`[x,y]`$ in $`T`$.
A morphism of R-trees is a map $`\varphi :TT^{}`$ such that every nondegenerate segment $`[x,y]`$ has a nondegenerate subsegment $`[x,w]`$ for which $`\varphi _{[x,w]}`$ is an isometry.
The morphism $`\varphi :TT^{}`$ folds at the point $`xT`$ if there are nondegenerate segments $`[x,y_1]`$ and $`[x,y_2]`$, with $`[x,y_1][x,y_2]=\{x\}`$, such that $`\varphi `$ maps each segment $`[x,y_i]`$ isometrically onto a common segment in $`T^{}`$. It is easy to see that the morphism $`\varphi :TT^{}`$ is an isometric embedding unless $`\varphi `$ folds at some point $`xT`$.
By an action of $`\mathrm{\Gamma }`$ on $`T`$ we mean an action by isometries. The action is minimal if $`\mathrm{\Gamma }`$ leaves no proper subtree of $`T`$ invariant. For any $`\mathrm{\Gamma }`$-action on $`T`$, there is a $`\mathrm{\Gamma }`$-invariant proper subtree which is minimal (see, e.g., \[CM\]). Also, if $`A_\gamma `$ is the isometry of $`T`$ corresponding to $`\gamma \pi _1S`$ for which $`inf_{yT}d(A_\gamma y,y)>0`$, then $`A_\gamma `$ has an axis $`l_\gamma `$ in $`T`$, i.e., an isometrically embedded line in $`T`$ which is invariant under $`A_\gamma `$ and which has the property that $`xl_\gamma `$ iff $`d(A_\gamma x,x)=inf_{yT}d(A_\gamma y,y)`$. The proof is a straightforward consequence of the non-positive curvature of $`T`$.
Assumption: Henceforth we will assume, without loss of generality, that all actions are minimal.
We will need the following fact about small actions.
###### Lemma 2.1
Let $`\mathrm{\Gamma }=\pi _1(S)`$, $`S`$ a closed surface of genus at least two. If the action $`\mathrm{\Gamma }\times TT`$ is small then $`T`$ must have a vertex point.
Whenever speaking of vertex or edge points in a subtree of a tree $`T`$, we mean with respect to the space of directions in the subtree, not the ambient tree $`T`$.
Proof: If $`T`$ has no vertex points then it is isometric to R, so the action of $`\mathrm{\Gamma }`$ gives a representation $`\psi :\mathrm{\Gamma }\text{Isom}(\text{R})`$. As $`\psi (\mathrm{\Gamma })<\text{Isom}(\text{R})`$ is virtually abelian and $`\mathrm{\Gamma }`$ is not solvable, it must be that the kernel of $`\psi `$ contains two noncommuting elements. But $`S`$ is a closed surface of genus at least two, so sufficiently high powers of any two noncommuting elements of $`\mathrm{\Gamma }=\pi _1(S)`$ generate a free group. This free group lies in the kernel of the action, in particular stabilizes any nondegenerate edge of $`T`$, a contradiction. $``$
### 2.2 Holomorphic quadratic differentials
A holomorphic quadratic differential $`\mathrm{\Phi }`$ on the Riemann surface $`S`$ is a tensor given locally by an expression $`\mathrm{\Phi }=\phi (z)dz^2`$ where $`z`$ is a conformal coordinate on $`S`$ and $`\phi (z)`$ is holomorphic. Such a quadratic differential $`\mathrm{\Phi }`$ defines a measured foliation in the following way. The zeros $`\mathrm{\Phi }^1(0)`$ of $`\mathrm{\Phi }`$ are well-defined and discrete. Away from these zeros, we can choose a canonical conformal coordinate $`\zeta (z)=^z\sqrt{\mathrm{\Phi }}`$ so that $`\mathrm{\Phi }=d\zeta ^2`$. The local measured foliations ($`\{\text{Re}\zeta =\text{const}\}`$, $`|d\text{Re}\zeta |`$) then piece together to form a measured foliation known as the vertical measured foliation of $`\mathrm{\Phi }`$.
### 2.3 Actions dual to a measured foliation
Let $`(,\mu )`$ denote the vertical measured foliation of $`\mathrm{\Phi }`$. Lift it to a $`\pi _1S`$-equivariant measured foliation $`(\stackrel{~}{},\stackrel{~}{\mu })`$ on $`\stackrel{~}{S}`$. The leaf space $`R`$ of $`\stackrel{~}{}`$ is a Hausdorff topological space. Let $`\pi :\stackrel{~}{S}R`$ denote the projection. The leaf space $`R`$ of the measured foliation $`(\stackrel{~}{},\mu )`$ inherits a metric space structure from the measure $`\mu `$: a geodesic segment $`[x,y]`$ in $`R`$ is given by any path $`\gamma `$ in $`\text{H}^2`$ from a point in the leaf corresponding to $`x`$ to a point in the leaf corresponding to $`y`$, such that $`\gamma `$ is transverse to the leaves of the foliation $`\stackrel{~}{}`$. The distance $`d_R(x,y)`$ is simply $`\mu (\gamma )`$, and the metric space $`(R,d)`$ is an R-tree (see \[MS2\]). This tree is often not locally compact. For instance, when the leaves of the foliation on the surface $`S`$ are dense, we can find sequences of arcs $`C_n`$ transverse to the foliation with endpoints on singularities of $`\stackrel{~}{}`$ whose transverse measure $`\mu (C_n)`$ goes to zero, forcing the distance between the corresponding images of the (lifts of) vertices to also go to zero.
The action of $`\mathrm{\Gamma }`$ on $`\text{H}^2`$ preserves $`\mu `$, and so induces an isometric action of $`\mathrm{\Gamma }`$ on $`R`$. The stabilizers of this action are virtually cyclic, in particular are small.
The action of $`\pi _1S`$ on $`\text{H}^2`$ preserves $`\mu `$, and so induces an isometric action of $`\pi _1S`$ on $`R`$. The map $`\pi :\stackrel{~}{S}R`$ is equivariant with respect to this action.
### 2.4 The Hubbard-Masur Theorem
Holomorphic quadratic differentials on a Riemann surface $`S`$ are related to classes of measured foliations via the Hubbard-Masur Theorem. To set the notation, fix a Riemann surface $`S`$ and define a map $`HM:\text{QD}(S)(S)`$ from the space $`\text{QD}(S)`$ of holomorphic quadratic differentials on $`S`$ to the space $`(S)`$ of equivalence classes of measured foliations on $`S`$ that associates to $`\mathrm{\Phi }\text{QD}(S)`$ the class of its vertical measured foliation. A fundamental result is
###### Theorem 2.2 (Hubbard-Masur \[HM\])
HM is a surjective homeomorphism.
Remark. A proof of Theorem 2.2 in the spirit of the current work can be found in \[W2\].
An alternative phrasing will be convenient for us. Let $`Q(S)\text{Teich}(S)`$ denote the bundle of holomorphic quadratic differentials over $`\text{Teich}(S)`$: here the fiber over $`[S]\text{Teich}(S)`$ is the space $`\text{QD}(S)`$ of quadratic differentials holomorphic with respect to a complex structure $`S`$ in $`[S]`$. Let $`(,\mu )`$ denote a given measured foliation. Then the Hubbard-Masur Theorem shows that there is a well-defined section $`\mathrm{\Psi }_\mu :\text{Teich}(S)Q(S)`$ which associates to $`[S]\text{Teich}(S)`$ the holomorphic quadratic differential $`\mathrm{\Psi }_\mu (S)\text{QD}(S)`$ whose vertical measured foliation is measure equivalent to $`(,\mu )`$.
### 2.5 Moving around in the Hubbard-Masur section
In this subsection we give a basic property of the section $`\mathrm{\Psi }_\mu `$.
Let $`S`$ be a Riemann surface and let $`q`$ be a holomorphic quadratic differential with vertical measured foliation $`(,\mu )`$. Let $`p_0`$ be a singularity of $`q`$ and let $`L`$ be the maximal compact graph of singular vertical arcs through $`p_0`$ which connect $`p_0`$ to all the other singularities on the leaf through $`p_0`$. Consider a neighborhood $`𝒩`$ of L in $`S`$. We refer to the components $`\{s_i\}`$ of $`𝒩L`$ as sectors, and say that two sectors meet along a (nondegenerate) arc if their closures intersect along that arc. We observe that there is a natural correspondence of sectors near a maximal singular arc $`L`$ as above under Whitehead moves and isotopies of the foliation.
###### Proposition 2.3 (Sectors can be made adjacent)
Let $`S`$ be a Riemann surface, let $`q`$ be a holomorphic quadratic differential with vertical measured foliation $`(,\mu )`$, and let $`L`$, $`p_0`$ and $`\{s_i\}`$ be as above. Choose any pair of sectors $`s_{i_1}`$ and $`s_{i_2}`$ from the list of sectors $`\{s_i\}`$. Then there is a Riemann surface $`S^{}`$ and a holomorphic quadratic differential $`q^{}`$ on $`S^{}`$ so that the vertical foliation of $`q^{}`$ is measure equivalent to $`(,\mu )`$, and under the equivalence the sectors $`s_{i_1}^{}`$ and $`s_{i_2}^{}`$ (corresponding to $`s_{i_1}`$ and $`s_{i_2}`$, respectively) meet along an arc.
Note that both $`q`$ and $`q^{}`$ are in the image of the Hubbard-Masur section corresponding to $`(,\mu )`$. A self-contained proof of Proposition 2.3 is given in the appendix.
## 3 Harmonic maps to trees
### 3.1 Definition of harmonic map
Given a Lipschitz continuous map $`w:S(T,h)`$ from a Riemann surface $`S`$ to a locally finite metric tree $`T`$, we define the energy form to be the tensor
$$edzd\overline{z}=(w_{}_z_h^2+w_{}_{\overline{z}}_h^2)dzd\overline{z}$$
Since the map $`w`$ is Lipschitz, it is differentiable almost everywhere and bounded almost everywhere on closed balls; thus the form $`edzd\overline{z}`$ is defined almost everywhere with $`edzd\overline{z}`$ integrable over compacta. Note that even when $`T`$ is not locally finite, the image of any closed ball in $`S`$ is compact hence lies in a locally finite subtree of $`T`$, so this analysis applies.
Alternatively, for any conformal metric $`g`$ on $`S`$ with area form $`dA_g`$, the energy form may be expressed as follows. Choose an orthonormal frame $`\{v_1,v_2\}`$ at a point $`zS`$, and consider the pushforward vectors $`\{w_{}v_1,w_{}v_2\}`$. The the energy form is the 2-form $`\frac{1}{2}(w_{}v_1|_h^2+w_{}v_2_h^2)dA_g`$, or alternatively $`\frac{1}{2}tr_g(w^{}h)dA_g`$. The energy of the map $`w`$ is $`E=e𝑑zd\overline{z}`$. The map $`w`$ is a harmonic map if it is a minimum for this functional in its homotopy class of maps. We define the Hopf differential $`\mathrm{\Phi }`$ for a map $`w:ST`$ by
$$\mathrm{\Phi }=\varphi dz^2=4w_{}_z,w_{}_z_hdz^2$$
Note that $`\mathrm{\Phi }=\mathrm{\Phi }_{L^1}<2E`$.
### 3.2 Examples
In this subsection, we list some motivating examples of harmonic maps from Riemann surfaces to R-trees. Each example will illustrate a principle we will later use.
1. The map $`f(z)=\text{Re}\{z^2\}`$ as the most basic vertex fold.
The map $`f(z)=\text{Re}\{z^2\}:\text{C}\text{R}`$ can be viewed as a harmonic map from the Riemann surface C to the R-tree R. Observe that the preimage of the origin $`O\text{R}`$ is the pair of intersecting lines $`\{x=\pm y\}`$ which divides C into four sectors. The other level lines of a nonzero $`r\text{R}`$ consist of hyperbolas $`\{x^2y^2=r\}`$. The leaf space of the connected components of these level curves is the pair of coordinates axes. We conclude that the harmonic function $`f(z)`$ factors as a projection to the R-tree of the coordinates axes followed by a vertex fold of each half-axis to its negative, which results in the image R-tree R.
2. Here is an example from \[W3\]: begin with the holomorphic differential $`z^kdz^2`$ on C, whose vertical measured foliation is the set of curves $`\{Rez^{k+2}=c\}`$. When we project along this foliation, we obtain a harmonic map to a tree with $`k+2`$ prongs out of a single vertex.
3. Actions Dual to a Measured Foliation $`(,\mu )`$, as given in §2.3.
Here the harmonic map is simply the projection along the vertical foliation of the properly normalized Hubbard-Masur differential for $`(,\mu )`$. This characterization is independent of the particular Riemann surface chosen. We therefore observe the following.
###### Lemma 3.1 (When $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ agree)
If the action of $`\mathrm{\Gamma }`$ is dual to a measured foliation $`(,\mu )`$, then there is a well-defined Hopf differential section $`\mathrm{\Phi }:\text{Teich}(S)Q(S)`$ for $`\pi `$, and this section $`\mathrm{\Phi }`$ agrees with the Hubbard-Masur differential section $`\mathrm{\Psi }:\text{Teich}(S)Q(S)`$ for $``$.
Proof: The lemma is effectively the content of \[W2\], which we now summarize; for complete details, see \[W2\]. (Later on, in §4, we shall give an independent proof of the existence of a harmonic map dual to a measured foliation.) As in §2.3, a measured foliation $`(,\mu )`$ on $`S`$ lifts to an equivariant measured foliation $`(\stackrel{~}{},\stackrel{~}{\mu })`$ on $`\stackrel{~}{S}`$; we can project along the leaves to obtain an R-tree $`(R,d)`$, with this construction also yielding an equivariant map $`\pi _0:\stackrel{~}{S}(R,d)`$.
For each complex structure $`\sigma `$ on $`S`$, we can minimize the energy in the equivariant homotopy class of $`\pi _0`$ obtaining \[W2; proof of Prop. 3.1\] a map $`\pi :(S,\sigma )(R,d)`$ whose Hopf differential $`\mathrm{\Phi }_R(\sigma )`$ has vertical foliation measure equivalent to $`(,\mu )`$. (This argument is a straightforward application of Ascoli-Arzela, with a crucial use of the axes of group elements in $`R`$ to control (see \[W2; Lemma 3.4\]) the images of some points by elements of the minimizing sequence of maps.) This characterizes the differential uniquely \[HM\]; for a harmonic maps argument for this uniqueness, see \[W2; §4\]. Here the point is that both maps can be given as projections along minimal stretch foliations of Hopf differentials and the distance between the image points of the two maps can be equivariantly defined, and is a subharmonic function. \[As the pullback of a smooth convex function off of the zeroes of the Hopf differentials, this pullback of the distance function is smooth and subharmonic (i.e. submean for balls of fixed radii) away from a discrete set of singularities and continuous across them; hence it is subharmonic everywhere. (Compare Proposition 3.2)\] The maximum principle then applies, showing that the distance must be constant. Some geometry of the tree, in particular the fact that it has branches, forces that constant to vanish. Thus $`\mathrm{\Phi }_R(\sigma )=\mathrm{\Psi }_\mu (\sigma )`$. $``$
An important part of our proof of Theorem 1.1 will be a converse (Lemma 5.3) to Lemma 3.1.
4. Actions Dual to the Measured Foliation of the Hopf Differential for an Arbitrary Harmonic Map from a Surface.
Let $`f:SX`$ be a harmonic map from the Riemann surface $`S`$ to a metric space (possibly Riemannian, possibly singular). Let $`\mathrm{\Phi }`$ denote the associated Hopf differential; we will see in §3.3, that this Hopf differential is a holomorphic quadratic differential on $`S`$. Lift this differential $`\mathrm{\Phi }`$ to an (equivariant) differential $`\stackrel{~}{\mathrm{\Phi }}`$ on the universal cover $`\stackrel{~}{S}`$, and consider its vertical (corresponding to the minimal stretch directions of $`\stackrel{~}{f}`$) measured foliation $`(\stackrel{~}{},\mu )`$ and associated projection $`\pi :\stackrel{~}{S}R`$ to the leaf space (see §2.3).
Part of the content of the previous example is that the equivariant projection map $`\pi :\stackrel{~}{S}R`$ from $`\stackrel{~}{S}`$ to $`R`$ is harmonic.
Our proof of Skora’s theorem involves a study of the relationship between the harmonic map we will construct from $`\stackrel{~}{S}`$ to $`T`$ and the associated harmonic map $`\pi :\stackrel{~}{S}R`$ from $`\stackrel{~}{S}`$ to the leaf space $`R`$.
Remark. In \[W4\], we study the product harmonic map $`(\stackrel{~}{f},\pi ):\stackrel{~}{S}X\times \text{R}`$, and find that it is also conformal, after a slight homothety of R. We also find that when $`X`$ is smooth and two-dimensional, then this map is a stable minimal map.
5. A Harmonic One-Form with Integral Periods.
Project a square torus $`T^2`$ along its natural vertical foliation to the circle $`S^1`$. This map is clearly harmonic. Now, there is a genus two surface $`S`$ which is a branched cover over $`T^2`$, and the $`1`$-form $`dz`$ lifts to yield a holomorphic one form on $`S`$. One can still project along leaves of this one form to a figure 8 which is a branched cover of the original $`S^1`$. Hence by composing with the map to $`S^1`$, we see that there is an associated harmonic map $`f:SS^1`$ and, via the one-form lifted from $`dz`$, an associated holomorphic one-form with integral periods on $`S`$.
As we vary $`S`$ in Teich $`(S)`$ (say in a family $`S_t`$), the holomorphic one forms with the same $`A`$-periods (in the usual notation) varies continuously through one-forms, say, $`\omega _t`$. It is interesting to consider the topology of the foliations $`_t`$ that integrate $`\mathrm{ker}\text{Re}\omega _t`$.
The original surface $`S`$ could be described as being constructed from a pair of cylinders $`C_1`$ and $`C_2`$ bounded by circles $`\{S_{11},S_{12}\}`$ and $`\{S_{21},S_{22}\}`$, respectively. Each $`S_{ij}`$ is composed of two semicircles $`s_{ij}^\text{t}`$ and $`s_{ij}^\text{b}`$. Now, the upshot of the notation is that $`S`$ is defined by identifying $`s_{11}^\text{b}`$ to $`s_{12}^\text{t}`$, $`s_{12}^\text{b}`$ to $`s_{21}^\text{t}`$, $`s_{21}^\text{b}`$ to $`s_{22}^\text{t}`$ and $`s_{22}^\text{b}`$ to $`s_{11}^\text{t}`$ in the natural way, and the foliations parallel to the core curves of the cylinders become $`_0`$.
A natural motion in Teichmüller space is to slightly rotate one of these cylinders against the other. This has the effect in our case of preserving the topology of $`_0`$, up to a Whitehead move which splits the singularities from being locally a pair of coordinate axes as in Example 1 to a “double $`Y`$” configuration. Of course, as $`_t`$ is the foliation of a harmonic one-form, we see that this new synthetically constructed foliation inherited from the cylinders, which is actually the foliation of the Hubbard-Masur differential for $`(_0,\mu _0)`$ on $`S_t`$, is not $`_t`$. We conclude in this case that for the natural $`\pi _1S`$\- action on R defined via the one form $`\omega `$, it is not the case that the Hopf differential section $`\mathrm{\Phi }:\text{Teich}(S)Q(S)`$ agrees with the Hubbard-Masur differential section $`\mathrm{\Psi }:\text{Teich}(S)Q(S)`$.
### 3.3 Local structure
R. Schoen has emphasized (see \[Sc\]) that a map for which the energy functional is stationary under reparametrizations of the domain has a Hopf differential which is holomorphic: one uses suitable domain reparametrizations to show that the Hopf differential satisfies the Cauchy-Riemann equations weakly, and then Weyl’s Lemma forces the Hopf differential to be (strongly) holomorphic. We observe that in this argument, the range manifold may be singular.
The vertical and horizontal foliations of the Hopf differential for $`w:ST`$ integrate the directions of minimal and maximal stretch of the gradient map $`dw`$, for smooth energy minimizing maps $`w:ST`$. As the image is one-dimensional, the harmonic map $`w`$ is a projection along the minimal stretch direction. Further, if one normalizes the conformal coordinates in a domain that avoids the zeroes $`\mathrm{\Phi }^1(0)`$ of the Hopf differential $`\mathrm{\Phi }`$ so that $`\mathrm{\Phi }=dz^2`$ in that neighborhood, then one sees from the geometric definition of $`\mathrm{\Phi }`$ above that the energy-minimizing map takes maximal stretch segments of measure $`ϵ`$ to segments in $`T`$ of length $`ϵ`$.
### 3.4 Effect on convex functions
A function defined on a an R-tree is convex if its restriction to every geodesic is convex in the classical sense. Recall that a function is subharmonic if it is submean, that is it’s value at any point is less than or equal to its average in a small ball around that point. Harmonic maps between Riemannian manifolds pull back convex functions to subharmonic functions (see, e.g. \[I\]). This important property extends to the case of R-tree targets.
###### Proposition 3.2
A harmonic map from a surface $`S`$ to an R-tree pulls back (germs of) convex functions to (germs of) subharmonic functions.
Proof: We first argue that the map $`\pi :\stackrel{~}{S}R`$ to the leaf space $`R`$ pulls back germs of convex functions on $`R`$ to germs of subharmonic functions on $`\stackrel{~}{S}`$. Locally, the level set of the vertex, say $`VR`$, near a zero or pole of the Hopf differential $`\stackrel{~}{\mathrm{\Phi }}`$ divides the neighborhood of the singularity into ’sectors’, with the natural coordinate $`\zeta `$ mapping each sector conformally onto a neighborhood of zero in the upper half-plane (see \[St;§7.1\]). Under this mapping of a sector, the foliation of preimages of points in the tree $`R`$ (in a sector) is taken to the horizontal foliation of the half-plane given by curves of the form $`\{y=const\}`$.
While it is not essential for the proof at this point, let us now consider a convex function $`F`$ defined on the tree $`R`$ near the point $`VR`$. This function pulls back to a function on a collection of sectors, which is constant on the horizontal levels $`\{y=const\}`$, and convex in $`y`$. Since any sector can be taken to any other sector by an appropriate rotation, it is straightforward to see that this pullback is submean. (A more detailed argument is also given below, in the case of the tree $`T`$.)
With this in mind, let us return to the original case of the map $`f:\stackrel{~}{S}T`$. In the neighborhood of the singularity $`p`$ of the Hopf differential $`\stackrel{~}{(}\mathrm{\Phi })`$, we can regard our map as first projecting to a neighborhood of a vertex $`V`$ in R (this neighborhood of $`V`$ is metrically a k-pronged star out of $`V`$, with one prong for each sector, by construction), followed by a map of $`R`$ to $`T`$, in which several prongs of $`R`$ map to a single prong of $`T`$, this prong of T emanating out of the image $`vT`$ of the vertex $`VR`$). Here we must have each prong taken injectively to a prong, because of the form of the map $`f:\zeta Re\zeta =\xi `$ away from singularities of $`\stackrel{~}{(}\mathrm{\Phi })`$.
In order to see why the map $`f:\stackrel{~}{S}T`$ pulls back germs of convex functions on $`T`$ to germs of subharmonic functions on $`\stackrel{~}{S}`$, we make one crucial observation: we note that neighboring sectors on $`\stackrel{~}{S}`$ must be taken to different prongs out of $`f(p)T`$; this is because a small arc transverse to the common boundary leaf of the pair of sectors is projected by $`f`$ injectively into $`T`$, once again because in a neighborhood of such an arc, there are no singular points, and so the map $`f`$ is of the form $`\zeta Re\zeta =\xi `$. This implies that the pre-image under $`f`$ of any given prong in $`T`$ consists of at most half of the sectors abutting $`p`$.
Consider then a convex function $`F`$ on the tree $`T`$ near a point $`pT`$. This function pulls back to a function on a collection of sectors, which is constant on the horizontal levels $`\{y=const\}`$, and also convex in $`y`$. Suppose we have that $`F(v)=0`$, so we need the mean value of $`f^{}F`$ to be non-negative on a disk $`D`$ around $`p`$. Of course, since $`F`$ is convex on $`f(D)T`$, we know that $`F`$ can be negative on at most one prong of $`f(D)`$, and must be non-negative on the other prongs. Moreover, since $`F`$ is convex, if we average $`f^{}F`$ over a pair of sectors, one in which $`f^{}F`$ is nonpositive, and one in which $`f^{}F`$ is non-negative, we see that the sum of the averages must be non-negative. (To see this, conformally map each sector to a half-plane (say $`\{y0\}`$), and then glue the halfplanes together so that $`f^{}F`$ is convex in the coordinate $`y`$ across the foliation.) Then we simply apply the observation of the previous paragraph to conclude that since $`f^{}F`$ is non-negative on at least half of the sectors, the average of $`f^{}F`$ on the union of sectors (i.e th disk $`D`$) must be non-negative, as required. $``$
## 4 Constructing a morphism from a geometric action to the given action
Let $`\mathrm{\Gamma }=\pi _1(S)`$, $`S`$ a closed surface of genus at least 2, and let $`\mathrm{\Gamma }\times TT`$ be an action (not necessarily small) on an R-tree $`T`$. In this section we construct an action of $`\mathrm{\Gamma }`$ on an R-tree $`R`$ which is dual to a measured foliation, and a $`\mathrm{\Gamma }`$-equivariant morphism $`\varphi :RT`$.
We will think of $`S`$ as having a fixed hyperbolic structure, and so the universal cover $`\stackrel{~}{S}`$ is the hyperbolic plane $`\text{H}^2`$. Since $`T`$ is contractible, there is a $`\mathrm{\Gamma }`$-equivariant Lipschitz continuous map $`f_0:\text{H}^2T`$. To be concrete, one may lift a triangulation of $`S`$, define $`f_0`$ by equivariance on the 0-skeleton of this triangulation, then extend (by contractibility of $`T`$) equivariantly to the 1-skeleton and 2-skeleton.
### 4.1 Finding the foliation using a harmonic map
Our first goal is to find a harmonic $`f`$ map in the equivariant homotopy class of the $`\mathrm{\Gamma }`$-equivariant continuous map $`f_0:\text{H}^2T`$ constructed at the beginning of §4. The harmonic map $`f`$ will have the property that there is a measured foliation $`(,\mu )`$ on $`S`$ so that every leaf of $`\stackrel{~}{}`$ gets mapped to a point under $`f`$. While it is possible to use the general theories of Korevaar-Schoen \[KS\] and Jost \[J1, J2\] on harmonic maps to nonpositively curved metric spaces, we will construct the harmonic map from elementary methods here.
To carry this out, we choose balls $`B_1,\mathrm{},B_n`$ on $`S`$ so that:
* the balls are topologically trivial
* the restriction $`f_0_{\widehat{B}_j}`$ of $`f_0`$ to a lift $`\widehat{B}_j`$ of $`B_j`$ is not a constant map for $`j=1,\mathrm{},n`$, and
* the set $`\{B_1,\mathrm{}B_n\}`$ of balls is an open cover of $`S`$
Thus we have that each lift of $`B_j`$ is disjoint from every other lift of $`B_j`$, and the union of all the lifts of all the balls $`\{B_1,\mathrm{}B_n\}`$ covers $`\stackrel{~}{S}`$.
Now for each lift $`\widehat{B}_1`$ of $`B_1`$ the image $`f_0(\widehat{B}_1)`$ is a finite subtree of $`T`$. This follows from the fact that, for a basepoint $`b_1\widehat{B}_1`$, the image $`f_0(\widehat{B}_1)`$ lies in a compact subset $`K`$ of the space of directions at $`f_0(b_1)`$, and as this space of directions $`K`$ is discrete (from the definition of R-tree), it is also finite.
It is straightforward that there exists a unique harmonic map $`\widehat{f_0}:\widehat{B}_1T`$ so that $`\widehat{f_0}_{\widehat{B}_1}=f_0_{\widehat{B}_1}`$ (see the Appendix of \[W1\] for existence. To see uniqueness, apply the method of Cor. 3.2 of \[W3\] (see also §4 of \[W1\]): the distance between any pair of solutions would be subharmonic on $`\widehat{B}_1`$ and vanishing on $`\widehat{B}_1`$ – thus any pair of solutions coincide.). Moreover, if $`h(\widehat{B}_1)`$ is any other lift of $`B_1`$, the uniqueness of the harmonic map then forces $`\widehat{f_0}_{h(\widehat{B}_1)}=h\widehat{f_0}_{\widehat{B}_1}`$. Let $`\varphi _1`$ denote the map from the complete lift of $`B_1`$ to $`T`$. Then $`\varphi _1`$, being nonconstant, also has the following properties:
* $`\varphi _1`$ is projection along the vertical measured foliation of its Hopf differential, and
* $`\varphi _1`$ is $`C^{\mathrm{}}`$ on the interior of its domain (off of the zeroes of the Hopf differential of $`\varphi _1`$)
Set
$$f_1(z)=\{\begin{array}{cc}\varphi _1(z)\hfill & z\text{ lift of }B_1\hfill \\ f_0\hfill & \text{otherwise}\hfill \end{array}$$
Then $`f_1`$ is equivariantly homotopic to $`f_0`$, and is a $`C^{\mathrm{}}`$ equivariant projection (as above) along a measured foliation on the domain of $`\varphi _1`$.
We repeat this procedure for lifts of the ball $`B_2`$, using $`f_1`$ as the original map instead of $`f_0`$. We then obtain a map $`f_2`$. The situation is most interesting when $`B_1B_2\mathrm{}`$, as then the boundary values for $`\varphi _2`$ are defined by values of $`\varphi _1`$, which may not agree with those of the original $`f_0`$.
The main observation we need to make is the following: along most of a small neighborhood of $`\widehat{B}_2\widehat{B}_1`$ we have that $`\varphi _1_{\widehat{B}_1\widehat{B}_2}`$ and $`\varphi _2_{\widehat{B}_2}`$ extend to be a well-defined Lipschitz projection along a well-defined Lipschitz measured foliation. To see this note that $`\varphi _1_{\widehat{B}_1}`$ is $`C^{1,\alpha }`$ and the measure of the vertical foliation of the Hopf differential of $`\varphi _1`$ is defined by distance between image points in $`T`$ (see §2.5). As this also holds for $`\varphi _2_{B_2}`$, and $`\overline{\widehat{B}_1}\overline{\widehat{B}_2}`$ is compact, the claim follows, except at (a discrete set of) places where the boundary values $`f_1_{\overline{B_2}}`$ double back and result in small arcs in both $`\overline{B_1}`$ and $`\overline{B_2}`$ which close up in $`\widehat{B}_1\widehat{B}_2`$.
We follow the same procedures iteratively for lifts of $`B_3,\mathrm{},B_n`$ obtaining an equivariant map $`f_n:\stackrel{~}{S}T`$ which is a projection along a Lipschitz measured foliation except for a discrete set of places where the leaves are closed and homotopically trivial.
In these places, we do an equivariant surgery to the map. For any region consisting of a union of concentric closed leaves, consider the closure of the largest such region. We then collapse the region to a segment which maps to the point defined by the boundary leaves. Call the new (collapsed) map $`F:\stackrel{~}{S}T`$. It is evidently an equivariant map along a measured foliation with singularities that are $`k`$-pronged.
In \[W2\](Prop. 3.1), an elementary proof shows that the piecewise harmonic map $`F:\text{H}^2T`$ as above is equivariantly homotopic to a harmonic map $`f:\text{H}^2T`$. (This proof only requires that there are two elements of $`\mathrm{\Gamma }`$ whose axes in $`T`$ have unbounded intersection. This property is much weaker than requiring that the action be small, but, for our purposes, follows from Lemma 2.1 above.) Moreover, attached to $`f`$ is a holomorphic quadratic differential $`\stackrel{~}{\mathrm{\Phi }}_0`$, the Hopf differential of $`f`$, with the following properties (see \[W2; §2.2\]):
* The vertical measured foliation of $`\stackrel{~}{\mathrm{\Phi }}_0`$ is equivalent to $`(\stackrel{~}{},\stackrel{~}{\mu })`$.
* The leaf space of the vertical foliation of $`\stackrel{~}{\mathrm{\Phi }}_0`$ is $`R`$, and the vertical measure pushes down (say via $`\pi :\text{H}^2R`$) to the metric on $`R`$. This map is harmonic.
* On neighborhoods $`B\stackrel{~}{S}`$ which are disjoint from $`\stackrel{~}{\mathrm{\Phi }}^1(0)`$, the map $`f|_B`$ agrees with $`\pi |_B`$ up to an isometry, while $`\pi |_B`$ is the projection $`zRez`$ in a natural coordinate system.
This last property is quite important for the sequel, so we recall some the details from, for instance, \[W1; p. 273\] and \[W2; p. 117\]. By §2.2, there is a canonical coordinate $`\zeta =\xi +i\eta `$ so that $`\mathrm{\Phi }_0=d\zeta ^2`$ on $`B`$. In its guise as a Hopf differential, of course, the definitions from §3.1 provides that $`\mathrm{\Phi }_0=f_{}\xi ^2f_{}\eta ^2+2i<f_{}\xi ,f_{}\eta >`$. Combining these two descriptions of $`\mathrm{\Phi }_0`$ and using that B is one-dimensional, we find that $`f|_B`$ is isometric to the map $`\zeta Re\zeta =\xi `$.
### 4.2 Definition of the morphism $`\varphi :RT`$
Define an associated harmonic projection $`\pi :\stackrel{~}{S}R`$ via the construction in Example 3.2.4. Define also a map $`\varphi :RT`$ by $`\varphi =f\pi ^1`$. We claim that $`\varphi `$ is a morphism. To see this, let $`I`$ denote a nondegenerate segment on the tree $`R`$; we must find a non-degenerate subsegment $`JI`$ for which $`\varphi |_J`$ is an isometry. Well, as $`R`$ is defined via projection $`\pi :\stackrel{~}{S}R`$, we can find an arc $`\gamma \stackrel{~}{S}`$ with $`\pi (\gamma )=I`$. Here $`\gamma `$ is quasi-transverse to $`\stackrel{~}{}`$ (in the sense of \[HM\],p. 231) and $`\mu (\gamma )=\mathrm{}_R(I)`$. On any subarc $`\gamma ^{}`$ of $`\gamma `$ which avoids the zeros of $`\stackrel{~}{\mathrm{\Phi }}_0`$, we may write (as we did at the end of the previous subsection) $`\stackrel{~}{\mathrm{\Phi }}_0=dz^2`$ for a choice of conformal coordinate in a neighborhood of $`\gamma ^{}`$, and (again as in the previous subsection) $`f`$ is an isometric submersion. Then for $`J=\pi (\gamma ^{})\pi (\gamma )=I`$, we have that $`\varphi |_J=f|_\gamma ^{}`$ which is an isometry by construction.
Finally, $`\varphi `$ is surjective by the minimality hypothesis, and $`\varphi `$ is equivariant since $`f`$ is equivariant.
## 5 Proving that $`\varphi `$ doesn’t fold
### 5.1 No edge folds
It is a direct consequence of harmonicity of $`\varphi `$ that $`\varphi `$ does not fold at edge points. This is actually implicit in the proof above that $`\varphi `$ is a morphism, but we give a slightly different proof in the next proposition, to which we will refer back several times in the sequel.
###### Proposition 5.1 (no edge-point folds)
The morphism $`\varphi :RT`$ does not fold at an edge point $`xR`$.
Proof: The pre-image of an edge point is a nonsingular leaf of the foliation $`\stackrel{~}{}`$. Any point $`z_0`$ on this leaf has a neighborhood $`𝒩`$ foliated by non-singular arcs of leaves, and admits a conformal coordinate $`z=x+iy`$ with the foliation parallel to $`\mathrm{ker}(\text{Re}dz)`$. If $`\varphi :RT`$ were to fold at an edge point $`\pi (z_0)`$, then the harmonic map on the neighborhood $`𝒩`$ would necessarily have the form $`z|\text{Re}z|`$, which is, of course, not harmonic.
Alternatively, using the same notation for the morphism $`\varphi `$ folding at an edge point $`\pi (z_0)`$, letting $`p_0`$ denote the point $`p_0=\varphi \pi (z_0)`$, we may apply the maximum principle to the function $`h=f^{}(d_T(p_0,))`$ on a neighborhood of $`z_0`$. Here $`d_T(p_0,)`$ is convex on $`f(𝒩)`$, while $`f^{}(d_T(p_0,))`$ is not subharmonic on $`𝒩`$, contradicting Proposition 3.2. $``$
Note that, at this point, we have shown that for any small $`\mathrm{\Gamma }`$-action on an R-tree $`T`$, there is an action on a tree $`R`$, dual to a measured foliation, and a $`\mathrm{\Gamma }`$-equivariant morphism $`\varphi :RT`$ which folds only at vertex points.
### 5.2 No vertex folds
In this section we will show that, when the action of $`\mathrm{\Gamma }`$ on $`T`$ is small, the morphism $`\varphi `$ is an isometry. A crucial feature of our argument will be a lemma that says that for actions $`\mathrm{\Gamma }\times TT`$ which are not small, the choice of tree $`R`$ is not uniquely determined.
###### Proposition 5.2 (no vertex folds)
If the action $`\mathrm{\Gamma }\times TT`$ is small, then the morphism $`\varphi `$ does not fold at a vertex point $`vR`$.
The rest of this section is devoted to proving Proposition 5.2.
#### 5.2.1 Vertex fold gives bad family
We begin with the following generalization of Lemma 3.1.
###### Lemma 5.3
With notation as above, the following conditions are equivalent:
1. The action of $`\mathrm{\Gamma }`$ on $`T`$ is dual to the measured foliation $``$.
2. The morphism $`\varphi :RT`$ is an isometry.
3. The Hubbard-Masur section $`\psi _{}:\text{Teich}(S)Q(S)`$ for $``$ is the same as the Hopf differential section $`\mathrm{\Phi }:\text{Teich}(S)Q(S)`$ for $`T`$.
Proof: As $`R`$ is the dual tree of $`\stackrel{~}{}`$, it is clear that (2) implies (1). Lemma 3.1 states that (1) implies (3).
Now we prove that (3) implies (2). If $`\varphi `$ is not an isometry, then $`\varphi `$ must fold at some vertex point $`vR`$, by Proposition 5.1. Say $`\varphi (e_1)=\varphi (e_2)`$ for edges $`e_1,e_2`$.
We may assume that $`R`$ has vertices of valence at least 4: otherwise a vertex fold at a vertex $`vR`$ would be the fold of a $`3`$-pronged star to an interval or half-interval. Thus the map $`f`$ would restrict, in a neighborhood of the pre-image of $`vR`$, to a harmonic function on a disk where the corresponding Hopf differential has a 3-pronged zero. This is impossible, as harmonic functions are locally $`\text{Re}(cz^k)+O(z^{k+1})`$, for $`k`$ an integer.
\[Alternatively, the preimage of $`\pi ^1(v)`$ is a tree with discrete trivalent singularities. Near the singularities, this tree locally disconnects $`\stackrel{~}{S}`$ into three sectors, with the harmonic map $`f:\stackrel{~}{S}T`$ folding the image of one sector onto the image of an adjacent sector. Yet as the sectors meet along an edge, the proof of Proposition 5.1 applies to yield a contradiction.\]
Now consider the Hubbard-Masur section $`\psi _{}:\text{Teich}(S)Q(S)`$ for the foliation $`(,\mu )`$. We are assuming that $`\psi _{}(S)`$ has zeroes of order at least two. Let $`s_1,s_2`$ be the sectors of $``$ corresponding to the edges $`e_1,e_2`$. By Proposition 2.3 there is another quadratic differental $`q^{}=\psi _{}(S^{})`$ so that the sectors of the vertical foliation of $`q^{}`$ corresponding to $`s_1,s_2`$ have closures which meet along an edge. Since by assumption $`\psi _{}`$ is the same as the Hopf differential section $`\mathrm{\Phi }`$, this is a contradiction: it violates the maximum principle for the map $`f^{}`$ (defined as projection along the foliation of $`\psi _{}(S^{})=\mathrm{\Phi }(S^{}`$)), again as the map would locally have the form $`z|\text{Re}z|`$. $``$
Proof of Proposition 5.2: We now suppose, in expectation of reaching a contradiction, that the given action is not dual to a measured foliation, i.e. that $`\varphi `$ is not an isometry. The equivalence of (1) and (3) in Lemma 5.3 then implies that there is a family $`\{S_t\},t\text{R}`$ of distinct Riemann surfaces for which $`\psi _{}(S_t)\mathrm{\Phi }(S_t)`$ for $`t>0`$ and $`\psi _{}(S_0)=\mathrm{\Phi }(S_0)`$ (here $``$ is defined by setting $`\psi _{}(S_0)=\mathrm{\Phi }(S_0)`$ for some base point $`S_0`$), as we may as well assume for notational convenience that the two sections differ in a neighborhood of $`S_0`$: here we get a family of surfaces where the sections $`\psi _{}`$ and $`\mathrm{\Phi }`$ disagree rather than just a pair of points because the sections $`\psi _{}`$ and $`\mathrm{\Phi }`$ are continuous.
To set notation, we rephrase this as follows: there is a family $`\{S_t\}`$ of distinct Riemann surfaces and corresponding $`\mathrm{\Gamma }`$-equivariant harmonic maps $`f_t:\stackrel{~}{S}_tT`$, Hopf differentials $`\stackrel{~}{\mathrm{\Phi }}_t`$, vertical foliations $`_t`$, and projections $`\pi _t:\stackrel{~}{S}_tR_t`$ to R-trees with small $`\mathrm{\Gamma }`$-actions and universal covering maps $`p_t:\text{H}S_t`$ (choosing the notation so that $`t=0`$ corresponds to the original action). Note that the trees $`R_t`$ and morphisms $`\varphi _t`$ are distinct, and that the foliations $`_t`$ represent different points in $`𝒫(S)`$. If this were not true then $`\psi _{}(S_{t_1})=\mathrm{\Phi }(S_{t_1})`$ for some $`t_1>0`$, which would force the sections $`\psi _{}`$ and $`\mathrm{\Phi }`$ to agree over $`S_{t_1}`$, contrary to the definition of the family $`S_t`$.
The heart of our argument is the case when the foliations $`_t`$ are orientable and minimal. We begin with a reduction towards that case.
#### 5.2.2 Some leaf is not closed
Let $`eET`$ denote a point of $`T`$ which is not the image of a vertex in $`R_0`$ and which lies on the folded edge $`E`$ of $`T`$. We consider the leaves of $`_0`$ containing $`p_0f_0^1(e)S_0`$. Since $`e`$ lies on the folded edge $`E`$ there are at least two of these. Each such leaf which is a closed curve represents a (conjugacy class of) element of $`\mathrm{\Gamma }`$ which fixes the edge $`ET`$.
If each of these two leaves were closed, then they must represent the same element of $`\pi _1(S)`$: being simple closed curves, they do not represent powers of a common element of $`\pi _1(S)`$, hence some powers of these two elements in $`\pi _1(S)`$ must generate a free group since $`S`$ is closed and hyperbolic; this free subgroup of $`\pi _1(S)`$ stabilizes $`E`$, contradicting smallness. But these two closed leaves are not even freely homotopic. If they were then they would bound an annulus $`A`$ on $`S_0`$. Since $`A`$ has Euler-characteristic zero and the boundary components are leaves of $`_0`$, no singularity of $`_0`$ lies in $`A`$. Hence the foliation $`_0`$ on this annulus would be by closed curves parallel to the boundary and the harmonic map $`\pi |_A`$ restricted to this annulus would map to an interval, with constant boundary values. This forces the map to be everywhere constant, so that the Hopf differential vanishes on $`A`$, hence everywhere, an absurdity.
#### 5.2.3 The model case
So we may assume that one of the components $`\mathrm{}`$ of $`p_0f_0^1(e)`$ is not closed. Then consider a small arc $`\alpha S`$ transverse to $`\mathrm{}`$ and to $`_0`$. As the leaf $`\mathrm{}`$ is not closed, it is dense in a subsurface which we might as well take to contain $`\alpha `$ (after maybe reducing the size of $`\alpha `$ -see \[St\], §11). Indeed, we can find a finite number of edge points $`e_1,\mathrm{},e_n`$ so that the trajectories $`p_0f_0^1(e_i)`$ have closure equal to all of $`S_0`$.
Again, let $`\alpha `$ denote a small half-open arc transverse to $`_t`$ on $`S_t`$ with its endpoint on the singularity $`q_0S`$; we also assume that $`f_t(\alpha )E`$, the folded edge, and that $`\alpha `$ is chosen small enough to ensure that $`\varphi _t|_{\pi _t(\alpha )}`$ is an isometry. By 5.2.2, we may assume that the nonsingular leaves through $`\alpha `$ are not closed on $`S_t`$. (If a non-singular leaf were closed, it would be contained in a neighborhood of non-singular closed leaves (\[St\], §9.3) and so there would be no leaf through $`\alpha `$ which would also be dense in a subsurface containing $`\alpha `$. On the other hand, if every neighborhood of $`q_0`$ in $`\alpha `$ had regular closed leaves, since there are but a finite number of (maximal) ring domains (i.e. maximal neighborhoods of regular closed leaves) in $`_t`$, we see that a neighborhood of $`q_0`$ in $`\alpha `$ is contained in one of these ring domains. If this were true for all arcs $`\alpha `$ as above with $`f_t(\alpha )E`$, we would be in the situation of 5.2.2, a contradiction.)
We begin with the model case of $`_t`$ being orientable and minimal, i.e., every non-singular leaf is dense. The general case will follow from technical modifications to the proof in this case, but the essential ideas will be the same as in this model case.
Now, under the assumption that $`_t`$ is minimal and orientable, we see that the first return map $`P_t:\alpha \alpha `$ determines an interval exchange map $`\sigma _t:\alpha \alpha `$ on $`\alpha `$ (see \[St\], p. 58). Moreover, one can reconstruct the measured foliation $`(_t,\mu _t)`$ directly from the interval exchange map $`\sigma _t:\alpha \alpha `$. We recall that this interval exchange map $`\sigma _t`$ is determined by looking at the largest open subintervals $`R_i(t)`$ of $`\alpha `$ on which $`P_t`$ is continuous. The endpoints $`\{x_0(t)=q_0,x_1(t),\mathrm{},x_N(t)\}`$ of these subintervals are contained in singular leaves of $`_t`$, and hence (have lifts to $`\stackrel{~}{S}`$ which) project to vertex points of the tree $`R_t`$.
We know that the set of vertex points in $`R_t`$ is totally disconnected, as they are the image of the countable discrete set in H of zeroes of $`\mathrm{\Phi }_t`$. It is also easy to see from this that the set of vertex points of the tree $`\varphi _t(R_t)`$ in $`T`$ is also totally disconnected. We now assume, postponing the proof until the end of this subsection, that for each $`t`$ there is some vertex point $`vR_t`$ such that $`\varphi _t(v)`$ a vertex point.
Continuity argument: Our main observation is that, since the $`\mathrm{\Gamma }`$-equivariant maps $`f_t:\text{H}^2T`$ are continuous in $`t`$, we see that if $`f_t(\stackrel{~}{x_i(t)})`$ is a vertex in $`T`$, then as the vertices in $`T`$ are a totally disconnected set, the family $`f_t(\stackrel{~}{x_i(t)})`$ is constant in $`t`$. By the previous paragraph, there must be at least one endpoint $`x_i(t)`$ whose lift $`\stackrel{~}{x_i(t)}`$ projects to a vertex in $`R_t`$. Since $`_t`$ is minimal and $`f_t`$ is equivariant, we have that $`\mathrm{\Gamma }f_t(\stackrel{~}{x_i(t)})=\mathrm{\Gamma }f(\stackrel{~}{x_i})`$ is dense in $`f(\stackrel{~}{\alpha })`$, for lifts $`\stackrel{~}{x_i(t)}`$ and $`\stackrel{~}{\alpha }`$ with $`\stackrel{~}{x_i(t)}\stackrel{~}{\alpha }`$. Letting $`\mathrm{\Gamma }_{x_i(t)}=\stackrel{~}{\alpha }\pi _t^1(\mathrm{\Gamma }\pi _t\stackrel{~}{x_i(t)})`$, we see that $`f_t|_{\mathrm{\Gamma }_{x_i(t)}}`$ is constant in $`t`$, which forces $`f_t(\stackrel{~}{x_j(t)})`$ to be constant in $`t`$ for each $`j`$.
Since the measure of $`\stackrel{~}{\alpha }`$ between consecutive vertices $`\stackrel{~}{x_i(t)}`$ and $`\stackrel{~}{x_{i+1}(t)}`$ (for $`i=0,\mathrm{},N1`$) is determined by the distance $`d_T`$ ($`f_t(\stackrel{~}{x_i(t)})`$, $`f_t(\stackrel{~}{x_{i+1}(t)}`$) in the tree $`T`$, we see that these measures are also constant. Of course, after projecting from the cover $`\stackrel{~}{S}`$ to the surface $`S`$, we see that the endpoints $`x_i(t)\alpha `$ are also constant in $`t`$.
Finally, observe that the first return maps $`P_t:\alpha \alpha `$ vary continuously in $`t`$ on the interiors of the intervals $`R_i(t)`$ (and are affine there); hence, since the endpoints $`x_i(t)`$ are constant in $`t`$, we see that the maps $`P_t`$ are constant in $`t`$ as well. We conclude that the interval exchange maps $`\sigma _t`$ are constant in $`t`$, so that $`(_t,\mu _t)=(_0,\mu _0)`$ after we reconstruct $`(_t,\mu _t)`$ from $`\sigma _t:\alpha \alpha `$. Hence we are done by Lemma 5.3.
Proof that some vertex point maps to a vertex point: Since this property is preserved under perturbations of the map, it is enough to prove the statement for some $`t`$.
Suppose this were not the case. Then every vertex point of every $`R_t`$ maps to an edge point of $`T`$. Hence by Lemma 2.1 some edge point of each $`R_t`$ maps to a vertex point of $`T`$. Since there are finitely-many $`\mathrm{\Gamma }`$-orbits of vertex points, there exists $`\delta _t>0`$ so that, on a $`\mathrm{\Gamma }`$fundamental domain of $`R_t`$, any such edge point of $`R_t`$ has distance at least $`\delta _t`$ from any vertex point of $`R_t`$. For $`t`$ small, we may take all $`\delta _t>\delta `$, for some fixed $`\delta >0`$.
We first claim that by making a small perturbation in Teichmüller space from $`S`$ to $`S_t`$ we may assume that $`_t`$ has a closed leaf $`\lambda `$ representing an edge point $`xR_t`$ within a $`\delta /6`$-neighborhood of some vertex point $`v_t`$; necessarily, then there is a whole nondegenerate edge $`E`$ containing $`x`$ which is both within a distance $`\delta /3`$ of $`v_t`$ and fixed by an element $`g\mathrm{\Gamma }`$. This first claim follows from essentially the same argument we used in the continuity argument above: take a small arc which abuts the vertex $`v_tR_t`$, and consider the image $`\alpha `$ on $`S_t`$ of a lift of that arc. The foliation $`_t`$ is determined by the interval exchange defined by the first return map on that arc $`\alpha `$. In particular, perturbations of $`_0`$ are given by perturbations of that first return map, and we can find such a perturbation $`_0`$ so that $`_t`$ has a closed leaf through $`\alpha `$.
Now we make a few observations about our situation: since (1) all vertices are being folded away to edge points creating edges of radius at least $`\delta /2`$ from the image of $`v_t`$, but (2) on the surfaces $`S_t`$, no pair of adjacent sectors are having their $`R_t`$ images folded together (by the argument late in the proof of Proposition 3.2), we see that for any point $`e^{}`$ in any edge $`E^{}`$ within $`\delta /2`$ of the image of $`v_t`$ in $`T`$, we must have at least two distinct leaves on $`S_t`$ whose lifts project to $`e^{}`$. But this contradicts smallness, as we showed in §5.2.2. Hence some vertex point maps to a vertex point.
Next we begin to loosen the hypotheses of the model case so as to eventually find ourselves in the general case, where $`_t`$ may be non-orientable and have several minimal components.
#### 5.2.4 Nonorientable case
Let us first drop the assumption that $`_t`$ should be orientable. This is merely a matter of generalizing the correspondence between measured foliations $`(_t,\mu _t)`$ and interval exchange maps $`S_t`$. The idea here goes back to Strebel (see \[St\]). We regard one side of $`\alpha `$ as $`\alpha ^+`$ and the other side as $`\alpha ^{}`$: if $`_t`$ is orientable, then the rectangles $`R_i(t)`$ have one edge on $`\alpha ^+`$ and another on $`\alpha ^{}`$, but if $`_t`$ is not orientable, a rectangle may have both edges on, say, $`\alpha ^+`$. Yet, if we now regard the first return map $`P_t`$ as a map $`P_t:\alpha ^+\alpha ^{}\alpha ^+\alpha ^{}`$, we can consider an associated interval exchange map $`S_t:\alpha ^+\alpha ^{}\alpha ^+\alpha ^{}`$ from which we can reconstruct $`(_t,\mu _t)`$. The endpoints $`\{x_i(t)\}`$ of the intervals $`R_i(t)`$ on $`\alpha ^+\alpha ^{}`$ still (have lifts which) map continuously into the disconnected set of vertices (constant in $`t`$) of $`T`$, so then, as before the endpoints $`\{x_i(t)\}`$, and the map $`P_t`$, $`S_t`$ are constant in $`t`$. We conclude that the measured foliations are also constant in $`t`$.
#### 5.2.5 Breaking the model case into pieces
We come finally to the most general part, where we no longer require that $`_t`$ is minimal. Then for $`_0`$ choose a collection of closed arcs $`\alpha _1,\mathrm{}\alpha _n`$ which are transverse to $`_0`$, and whose $`_0`$-orbits both cover $`\text{H}^2/\mathrm{\Gamma }_0`$ and intersect at most along some compact singular leaves. At this point, we also require the intervals $`\alpha _i`$ to have corresponding interval exchange maps for $`_0`$ which are either irreducible, i.e. we cannot (non-trivially) decompose $`\alpha _i=\alpha _i^{}\alpha _i^{\prime \prime }`$ with the interval exchange map $`\sigma _i`$ for $`\alpha _i`$ having a restriction $`\sigma _i|_{\alpha _i^{}}:\alpha _i^{}\alpha _i^{}`$ which preserves the proper subinterval $`\alpha _i^{}`$, or correspond to a single cylinder in $`_0`$ , so that the interval exchange is the identity on a single cylinder.
We claim that the measured foliation $`_t`$ on the whole surface $`\mathrm{\Sigma }(t)`$ is constant in $`t`$. This will give a contradiction by Lemma 5.3, proving the theorem.
Let $`\mathrm{\Sigma }_i(t)`$ be the subsurface of $`S_t`$ obtained by taking the closure of the orbit of $`\alpha _i`$ along the leaves of $`_t`$. Our restrictions on $`\{\alpha _i\}`$ have the effect of forcing either $`\mathrm{\Sigma }_i(0)`$ to be a cylinder or a surface on which $`_0`$ is minimal. We observe that the argument given earlier for the cases where $`_0`$ was minimal on the closed surface $`\text{H}^2/\mathrm{\Gamma }_0`$ continue to hold for the case where $`_0`$ is minimal on $`\mathrm{\Sigma }_i(0)`$. In particular, for $`\mathrm{\Sigma }_i(0)`$ a subsurface with almost every leaf dense, we see that the interval exchange maps $`\sigma _i(t)`$ must be constant in $`t`$. Yet, it is part of the basic construction of measured foliations from interval exchange maps that the topology of $`\mathrm{\Sigma }_i(t)`$ (as well as $`_t`$) is determined from the map $`\sigma _i(t)`$ (see, e.g., \[Ma1\]). Thus, as $`\sigma _i(t)=\sigma _i(0)`$, we see that $`\mathrm{\Sigma }_i(t)`$ is homeomorphic to $`\mathrm{\Sigma }_i(0)`$.
Now each boundary circle of each $`\mathrm{\Sigma }_i(t)`$ is a leaf of the foliation on that subsurface. This leaf may be taken to be singular as it would otherwise be an interior leaf of a cylinder of non-singular homotopic leaves, counter to the construction of $`\{\alpha _i\}`$. Thus the continuity argument also shows that the foliations on the cylindrical subsurfaces $`\mathrm{\Sigma }_j(t)`$ are constant in $`t`$. Hence the measured foliation on each subsurface $`\mathrm{\Sigma }_i(t)`$ is constant in $`t`$. Finally, whenever two subsurfaces $`\mathrm{\Sigma }_i(t)`$ and $`\mathrm{\Sigma }_j(t)`$ have a common boundary component $`C(t)`$, the continuity argument shows that $`C(t)`$ cannot become a cylinder at any time $`t`$ as this would require the single vertex $`f_t(C(t))`$ to continuously deform into a non-trivial family of pairs of vertices, an absurdity. So we see that the identification of all the boundary components of all the $`\mathrm{\Sigma }_i(t)`$ are constant over $`t`$, so that $`_t`$ is constant.
## 6 Appendix
This appendix is dedicated to a proof of Proposition 2.3, which is partly implicit in \[HM\] and partly a “folklore theorem”. We provide here an elementary, geometric, and self-contained proof due almost entirely to Howard Masur (personal communication), who graciously permitted us to reproduce it here.
The proof can be reduced to the following claim: If either
1. $`q`$ has a pair $`\{z_1,z_2\}`$ of distinct zeros connected by an arc $`A`$ of a leaf of its vertical foliation, or
2. $`q`$ has a $`k`$-pronged singularity at $`z_3`$, and $`2`$ arbitrary sectors $`s_1,s_2`$ of this $`k`$-prong are specified,
then there is a Riemann surface $`S^{}`$ and a holomorphic quadratic differential $`q^{}`$ on $`S^{}`$ so that the vertical foliation of $`q^{}`$ is measure equivalent to $`(,\mu )`$ and
1. (in case (1) above) the zeros of $`q^{}`$ corresponding to $`\{z_1,z_2\}`$ coincide, or
2. (in case (2) above) the images of the sectors $`s_1,s_2`$ under the equivalence of foliations meet along an arc.
The proposition follows from the claim as follows. First apply (1) above to get $`s_1`$ and $`s_2`$ as sectors abutting a common singularity $`z`$. Then apply (2) above and we are done.
We are left to prove the claim.
Single cylinder case. We first prove the claim for Jenkins differentials, i.e. those differentials whose vertical foliations are but one foliated open (right Euclidean) cylinder $`C`$ with all singularities lying in $`\overline{C}`$. Here $`S`$ can be thought of as an identification space $`\pi :\overline{C}S`$, with identifications being made on $`\overline{C}`$. Let $`C_1,C_2`$ denote the 2 components of $`\overline{C}`$. Note that the graph $`L`$ lies in $`\overline{C}`$, and all the singularities on a single component of $`\overline{C}`$ are connected by $`L`$. Moreover, there are natural correspondences between topological or geometric operations on the surface $`S`$ and topological or geometric operations on $`\overline{C}`$. This means that if we continuously deform $`C`$ to another right Euclidean cylinder $`C^{}`$ (so that there is a canonical correspondence of identifications on $`\overline{C^{}}`$), then the canonical quadratic differential $`q^{}`$ on $`C^{}`$ (defined so that the metric $`|q^{}|`$ agrees with the Euclidean metric on $`C^{}`$ and all of whose vertical leaves are parallel to $`\overline{C^{}}`$) descends to a quadratic differential $`q^{}`$ on the identified surface $`S^{}`$ with the vertical foliation of $`q^{}`$ on $`S^{}`$ being Whitehead equivalent to the vertical foliation of $`q`$ on $`S`$.
To prove (1) and (2) above, we will first perform the desired operation on $`C`$ to obtain a new Euclidean cylinder $`C^{}`$ with canonically determined quadratic differential $`q^{}`$ as above. The important thing to check in each case is that we can do this so that the resulting Euclidean lengths $`\mathrm{}(C_1^{})`$ and $`\mathrm{}(C_2^{})`$ of the 2 components of $`\overline{C^{}}`$ are equal. This imediately implies that the identification $`\pi `$ determines an identification $`\pi ^{}:C^{}S^{}`$ to a Riemann surface $`S^{}`$, and that the canonical quadratic differential on $`C^{}`$ descends to a quadratic differential $`q^{}`$ on $`S^{}`$. By construction $`q^{}`$ has vertical foliation measure equivalent to that of the vertical foliation of $`q`$.
Consider case (1). Let $`A_1,A_2\overline{C}`$ denote the 2 components of $`\pi ^1(A)`$, where we recall that $`A`$ is the arc of the vertical foliation we wish to collapse. Note that $`\mathrm{}(A_1)=\mathrm{}(A_2)`$.
Case 1a: $`A_1`$ and $`A_2`$ lie in different components of $`\overline{C}`$. In this case contract both $`A_1`$ and $`A_2`$ to a point to give a Euclidean cylinder $`C^{}`$. Since $`C_1^{}`$ and $`C_2^{}`$ have the same Euclidean length, so we are done by the above.
Case 1b: $`A_1,A_2`$ lie on the same component of $`\overline{C}`$. First note that the arcs of $`L`$ have preimages in $`\overline{C}`$ which come in pairs, as neighborhoods of the arcs on the identification space have full neighborhoods, while neighborhoods of arcs on $`\overline{C}`$ have only half-neighborhoods. Hence since $`A_1,A_2`$ lie on the same component of $`\overline{C}`$, we must be able to find some collection of pairs of arcs on the other component the sum of whose lengths is at least that of the sum of the lengths of $`A_1`$ and $`A_2`$. (This is just a pigeon-hole principle: the arcs come in pairs whose lengths are equal and for which total lengths of all the arcs are the sum of the lengths of the boundary components of $`\overline{C}`$, yet each of these boundary components have the same lengths, so the fact that $`A_1`$ and $`A_2`$ contribute solely to one component of $`\overline{C}`$ forces some other family of pairs to contribute at least as much solely to the other component of $`\overline{C}`$.) Thus we act as before, contracting $`A_1`$ and $`A_2`$ on one component of $`\overline{C}`$ and simultaneously some other pairs of arcs the same amount on the other component of $`\overline{C}`$. It is quite important here that the contraction of the other components has no effect on our claim or our purpose; the proof of the first part of the claim concludes as before.
Now to prove part (2) of the claim. Under the identification map $`\pi :CS`$, each sector $`s_i,i=1,2`$ has a unique pre-image on $`C`$ as a neighborhood $`U_i`$ of a vertex $`v_i`$.
Case 2a: $`v_1`$ and $`v_2`$ lie on different components of $`\overline{C}`$.
We split the vertex $`v_1`$ into a pair of vertices $`v_{1,1}`$ and $`v_{1,2}`$ connected by an arc $`A_1`$, and we split the vertex $`v_2`$ into a pair of vertices $`v_{2,1}`$ and $`v_{2,2}`$ connected by an arc $`A_2`$ of the same length as $`A_1`$. We then re-identify the cylinder as before, with the only changes being that instead of identifying $`v_1`$ to $`v_2`$, we send $`A_1`$ isometrically onto $`A_2`$ (there is a unique way to do this which preserves the ordering of the sectors). The resulting surface gives $`S^{}`$ and $`q^{}`$ as required.
Case 2b: $`v_1`$ and $`v_2`$ lie on the same component, say $`C_1`$, of $`\overline{C}`$. Split $`v_1`$ and $`v_2`$ as in Case 2b. We now do a further deformation to make $`\mathrm{}(C_1^{})=\mathrm{}(C_2^{})`$.
If for some compact singular arc $`BL`$, we have both components of $`\pi ^1(B)`$ lying on $`C_2`$, then by lengthening $`B`$ we could achieve $`\mathrm{}(C_1^{})=\mathrm{}(C_2^{})`$. If this isn’t true, then by the pigeon-hole priciple, for each such $`B`$ we know $`\pi ^1(B)`$ has one component on $`C_1`$ and one on $`C_2`$.
Now observe that any singularity on the surface $`S`$ with, say, $`k`$ sectors, admits a cyclic ordering of these sectors $`s_1,\mathrm{}s_k`$ (where the closure of $`s_2`$ meets the closure of $`s_1`$ on one side and the closure of $`s_3`$ on the other side, and so on). Since we are in a case where each edge incident to a singularity on $`S`$ has preimages on both boundary components $`C_1`$ and $`C_2`$ of $`\overline{C}`$, and since sectors have preimages near components of $`\overline{C}`$ where their bounding arcs have preimages, we see that the sectors $`s_1,\mathrm{},s_k`$ also alternate between having preimages in $`C_1`$ and in $`C_2`$. This implies that all of the singularities on the surface $`S`$ have an even number of sectors.
We now claim that there are vertices $`w_1,w_2`$ in $`C_2`$ which are still identified by the identification rules, even after the splitting of $`v_1`$ and $`v_2`$. (Here the subtlety is that by first splitting $`v_1`$ and $`v_2`$, we have changed the identification rules, and hence the orbits of identified vertices on $`\overline{C}`$. Our vertices $`w_1`$ and $`w_2`$ must not only then correspond to each other by the original identification rules, but they must also lie in the same new orbit of vertices on $`\overline{C}`$, after the splitting of $`v_1`$ and $`v_2`$.) This finishes the proof of case 2b since we then split $`w_1`$ and $`w_2`$ to make $`\mathrm{}(C_1^{})=\mathrm{}(C_2^{})`$.
To see that there are such vertices $`w_1`$ and $`w_2`$, we recall that the total multiplicity of zeroes of a holomorphic quadratic differential on a Riemann surface $`S`$ of genus $`g`$ is equal to $`4g4`$ (Riemann-Roch). Thus, since in the case under consideration all of the singularities have an even number of sectors (and hence an even order of zero), we see that there is either one singularity $`z_0`$ with at least six sectors, or several singularities which all have at least four sectors. In the first case we see that any initial splitting of $`z_0`$ (by splitting a pair of vertices $`v_1`$ and $`v_2`$ on the same component $`C_1`$ of $`\overline{C}`$) would leave a topological foliation with two singularities of which at least one would have four sectors, with two of those sectors having preimages on $`C_2`$: we would then split the vertices of those sectors, say $`w_1`$ and $`w_2`$ to finish the case. In the second case, there is at the outset a singularity on $`S`$ whose preimages do not include $`v_1`$ and $`v_2`$, and which has at least a pair of sectors with preimages on $`C_2`$, as desired.
General case. We prove the general case by the now standard technique of approximating. By \[Ma2\] we may approximate $`q`$ by a sequence $`\{q_n\}`$ of Jenkins differentials on $`S`$. In case (1), let $`A`$ denote the arc of the vertical foliation of $`q`$ which we wish to contract. As $`q_n`$ approximates $`q`$, there is an arc $`A_n\overline{C}_n`$ which approximates $`A`$. Furthermore, there is a contractible neighborhood $`U`$ on the underlying differentiable surface which is a neighborhood of the arc $`A`$ and all arcs $`A_n`$ for $`n`$ sufficiently large.
Now, the Riemann surface $`S`$ is an identification space of each cylinder $`\overline{C}_n`$, with identifications being made on $`\overline{C}_n`$. As in the single cylinder case above, we can form new Riemann surfaces $`S_n^{}`$ equipped with quadratic differentials $`q_n^{}`$ by contracting the arcs $`A_n\overline{C}_n`$ and identifying as before; here the arc $`A_n`$ on $`S_n`$ bounded by a pair of low order zeros is replaced on $`S_n^{}`$ by a single higher order zero, say $`z_n^{}`$.
The important thing to notice about this operation is that the complement $`V=U^c`$ of the neighborhood $`US`$ is approximated by the closure of an open set $`V_n`$ on $`S_n^{}`$ which only avoids a small neighborhood of the high order zero $`z_n^{}`$; moreover, the conformal structures on $`V_n`$ compare uniformly to the conformal structure on $`V`$, hence to each other. Hence, by passing to a subsequence if necessary, we have that that $`S_n^{}`$ converges in (the interior of) Teichmuller space to a Riemann surface $`S^{}`$. It also follows that $`q_n^{}`$ converges to a holomorphic quadratic differential, say $`q^{}`$, on $`S^{}`$; as $`q_n`$ approximates $`q`$ and $`q_n^{}`$ is measure equivalent to $`q_n`$, we see that $`q^{}`$ is measure equivalent to $`q`$. Moreover, as the foliation of $`q_n^{}`$ results from a Whitehead move applied to the foliation of $`q_n`$ (which contracts $`A_n`$ to a point), the foliation of $`q^{}`$ is obtained from the foliation of $`q`$ via a Whitehead move which contracts $`A`$.
Case (2) is virtually identical: we still have uniform convergence of the conformal structures outside the pair(s) of neighborhoods of the vertices (or arcs) we are splitting to pairs of vertices connected by an arc.
$``$
Benson Farb:
Dept. of Mathematics, University of Chicago
5734 University Ave.
Chicago, Il 60637
E-mail: farb@math.uchicago.edu
Michael Wolf:
Dept. of Mathematics, Rice University
Houston, TX 77251
E-mail: mwolf@math.rice.edu
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# Contents
## 1 Introduction
BPS D-branes enjoy a double life. On the one hand as a conformal field theory described by open strings with Dirichlet boundary conditions , and on the other hand as classical solitons of supergravity<sup>1</sup><sup>1</sup>1See and references therein.. For a single D-brane the regimes of validity of these two descriptions are complementary. The conformal field theory description is valid at weak coupling; whereas the supergravity solution corresponds to a strongly interacting gravitational system, which corresponds to strong coupling. On the other hand, for a large number of branes $`N`$ the system can be well described by a classical solution at weak coupling . The two main properties of the BPS D-branes that assure this consistent dual behaviour are the fact that BPS-branes preserve a fraction of supersymmetry, hence they are stable under variations of the string coupling; and the fact that one can consider a superposition of a large number of D-branes in such a way that they still preserve a fraction of supersymmetry. These properties make possible, for instance, the entropy calculation of black-holes using D-branes . On the other hand, D-branes can be studied in more general situations, for which spacetime supersymmetry is not preserved. This is the case of the non-BPS D-branes . These branes have an exact conformal field theory description, whose consistency conditions do not rely on spacetime supersymmetry . Having such a precise description for the non-BPS D-branes in the conformal field theory side, it is natural to wonder whether the non-BPS D-branes also enjoy a double life and have a description in terms of a classical solution of a certain effective (super)gravity theory.
It has been recently suggested in that the unstable non-BPS branes of type II theories might have a gravitational counterpart described by a gravitational sphaleron<sup>2</sup><sup>2</sup>2See and references therein.. These are unstable solutions with finite energy which interpolate between two (possibly distinguishable) vacuum configurations. However, these solutions are unstable and probably do not remain valid classically. On the other hand, we expect that consistent classical solutions will be related to non-BPS branes which have some stability properties similar to the BPS D-branes.
In order to find the appropriate conditions for constructing a classical solution for a non-BPS D-brane we must find out which properties of the D-branes are such that (1) they assure the validity of a classical supergravity description and (2) they are not necessarily related to fractional supersymmetry. These properties certainly comprise stability, which ensures that the state survives at strong coupling. Moreover, stability can be based on grounds different from supersymmetry, for instance, being the lightest object carrying a certain quantum number . Another key property is the fact that we can superpose an arbitrary number of parallel D-branes, which is the same as having a no-force condition at all distances . This is a consequence of the BPS property, which can exist independently of supersymmetry. Therefore, it seems natural that in order to find a classical solution for a non-BPS D-brane, this brane should be stable and enjoy a no-force property. Stable non-BPS D-branes have been found in different theories . However, only non-BPS D-branes in a certain orbifold of type II theories are known to enjoy both stability and the no-force property . There it was shown that at a particular critical radius of the compact orbifold the non-BPS branes develop a Bose-Fermi degeneracy at one loop. The critical radius is in fact the value beyond which the non-BPS D-brane becomes unstable and can decay into a pair of BPS D-branes.
In this article we study the particular case of a stable non-BPS D-particle in type IIB string theory<sup>3</sup><sup>3</sup>3The gravity duals of BPS branes in orbifolds have been considered before in . orbifolded by $`(1)^{F_L}_4`$. This D-particle is a truncated D-brane , is charged electrically under a twisted R-R 1-form field and it is the strong-weak coupling dual of a non-BPS state in the orientifold $`\mathrm{\Omega }_4`$of type IIB, charged under the $`U(1)`$ field of the D5-O5 system . The coupling of the non-BPS D-particle to a twisted R-R vector field is the origin of the stability of this non-BPS D-particle . Moreover, at a critical radius of the compact orbifold this D-particle meets all the requirements suggesting the existence of a classical solution. We use the technique of the boundary state in the covariant formalism<sup>4</sup><sup>4</sup>4For a recent review see . to describe the non-BPS D-particle. We analyse the interaction potential between two non-BPS D-particles in relative motion. Remarkably, we find no force at order $`v^2`$ as for BPS D-branes. At the critical radius, the static force is moreover vanishing, hence the non-BPS D-particle presents a BPS-like behaviour, up to $`v^4`$ corrections. Unlike for BPS branes, we find no matching between the $`v^4`$ terms in the open and closed string description.
Using the boundary state for BPS D-branes, the long distance behaviour of the classical massless fields generated by the D-brane was computed in<sup>5</sup><sup>5</sup>5The same results were obtained earlier in using different techniques. , and it was shown that the asymptotic behaviour of the corresponding classical solution is precisely recovered. A BPS D$`p`$-brane is described by a boundary state with two parts, the NS-NS part and the R-R part. They generate the asymptotic behaviour of NS-NS massless fields (metric and dilaton), and of the R-R massless fields (R-R $`(p+1)`$-form potential), respectively. In this paper we implement the same technique to obtain the long distance behaviour of the non-BPS D-particle geometry. This is given by the asymptotic form of a metric and a dilaton propagating in the bulk, and a twisted R-R 1-form potential propagating in the orbifold fixed plane.
One difficulty about recovering the full solution for the non-BPS D-particle from its asymptotic form is that, in principle, there might be many different possible geometries with the same asymptotic structure. In order to restrict these geometries we will assume that the no force condition also takes place when one consider the complete geometry at the critical radius, as it happens for BPS branes . That is, we take into account the extra pieces that one would need to add to the asymptotic behaviour in order to recover the full form of metric. On the other hand, although the Bose-Fermi degeneracy occurs at any distance between the non-BPS D-branes , at short distances, open strings loops might spoil this property. In fact, it has been recently proposed in that even at 1-loop in the open string theory, the no-force is removed in favour of another vacuum configuration in which the branes attract each other. For these reasons, our assumptions will be only acceptable for distances much larger than the string scale, which is also the range of validity of the classical solution for a BPS D-brane .
Using the non-BPS D-particle as a probe in the background of another non-BPS D-particle we recover the no-force behaviour at the critical radius. Moreover, under the assumptions presented above we are able to obtain part of the complete metric, dilaton and twisted R-R 1-form generated by the D-particle. The velocity dependent part of the brane action multiplies a flat metric, agreeing with the vanishing of the $`v^2`$ in the interaction potential. Extending this property for the complete geometry of the non-BPS D-particle source, we are able to find more properties of the classical geometry. We expect that the solution is consistent classically at the critical radius only, since it is only there where one can consider a superposition of a large number of D-particles $`N`$.
We find a diagonal metric with $`SO(5)\times SO(4)`$ symmetry (in Einstein frame):
$$ds_{10,E}^2=g_{00}(y)dt^2+g_{mn}(y)dy^mdy^n+g_{ij}(y)dx^idx^j,$$
where we do the split $`\mu =(0,m,i)`$ according to the orbifold symmetry, and $`_4`$ acts on the $`i`$ directions. By using the assumptions mentioned above we are able to find the form of two of the components of the metric at the critical radius:
$`g_{00}(y)`$ $`=`$ $`\left(1+{\displaystyle \frac{\kappa _6T_0}{2a\mathrm{\Omega }_4}}(2\pi ^2\alpha ^{})^1{\displaystyle \frac{1}{|y|^3}}+\mathrm{}\right)^{\frac{7}{6}a},`$
$`g_{mn}(y)`$ $`=`$ $`\left(1+{\displaystyle \frac{\kappa _6T_0}{2a\mathrm{\Omega }_4}}(2\pi ^2\alpha ^{})^1{\displaystyle \frac{1}{|y|^3}}+\mathrm{}\right)^{\frac{1}{6}a}\delta _{mn}.`$
The expressions are given in terms of a parameter $`a`$, and some possible extra dependence in $`|y|^n`$, $`n<3`$ denoted by $`\mathrm{}`$ which remain to be determined. We find no expression for $`g_{ij}`$ since, as will be explained in Section 4, the D-particle cannot probe the precise geometry in these directions. Moreover, the form of the dilaton and twisted R-R potential at the critical radius are found to be:
$`\mathrm{e}^\varphi `$ $`=`$ $`\left(1+{\displaystyle \frac{\kappa _6T_0}{2a\mathrm{\Omega }_4}}(2\pi ^2\alpha ^{})^1{\displaystyle \frac{1}{|y|^3}}+\mathrm{}\right)^a,`$
$`𝒞_0^{(1)}`$ $`=`$ $`\left(1+{\displaystyle \frac{\kappa _6Q_0}{4a\mathrm{\Omega }_4}}{\displaystyle \frac{1}{|y|^3}}+\mathrm{}\right)^{\frac{4}{3}a}1.`$
Here $`T_0`$ is the tension of the D-particle and is related to its charge $`Q_0`$ by $`T_0=Q_0\pi ^2\alpha ^{}`$, hence at the critical radius the fields above are given in terms of a single function.
This article is organised as follows. In Section 2 we carry out the construction of the covariant boundary state for the non-BPS D-particle in the orbifold of type IIB generated by $`(1)^{F_L}_4`$, for the non-compact and compact cases. We also evaluate the amplitude for D-particles in relative motion and analyse the long and short distance behaviour of the interaction. In Section 3 we evaluate the asymptotic behaviour of the massless fields excited by the non-BPS D-particle, for the non-compact and compact cases. In Section 4 we recover the no-force property at the critical radius and derive part of the complete classical solution, by using the assumptions mentioned above. For completeness, we also include the derivation of the zero-mode of the twisted R-R boundary state in Appendix A, and the explicit form of the complete covariant boundary state for the D-particle in Appendix B.
## 2 Stable non-BPS D-particle in IIB/$`(1)^{F_L}_4`$
The stable non-BPS D-particle of type IIB string theory orbifolded by $`(1)^{F_L}_4`$ was described in . It corresponds to the S-dual of a massive particle-like state in the system of a D5-brane on top of an orientifold 5-plane, which is stable but non-BPS . Its stability is due to the fact that this particle is the lightest state charged under the $`U(1)`$ gauge field of the D5-O5 worldvolume theory. Similarly, the stability of the D-particle is due to its charge under a twisted R-R 1-form, which can be identified with the $`U(1)`$ gauge field in the worldvolume of a NS-5 brane on top of the orbifold fixed plane.
Let us consider first the non-compact theory, i.e. type IIB in Minkowski space orbifolded by $`(1)^{F_L}_4`$. The operator $`_4`$ corresponds to a reflection in the directions $`x^i`$, $`i=6,7,8,9`$. In the non-compact case, the orbifold contains a fixed plane at $`x^6=x^7=x^8=x^9=0`$, hence the orbifold breaks the $`SO(1,9)`$ symmetry down to $`SO(1,5)\times SO(4)`$. The operator $`F_L`$ is the spacetime fermion number of the left-sector, hence $`(1)^{F_L}`$ changes the sign of the of the R-R groundstate, without having any action on the oscillators. The closed string spectrum of the orbifold theory consists of an untwisted sector, given by type IIB states which are invariant under $`(1)^{F_L}_4`$, and a twisted sector localised on the orbifold fixed plane. If we split the coordinates as $`X^\mu =(X^\alpha ,X^i)`$, where $`X^\alpha `$, $`\alpha =0,\mathrm{},5`$, is longitudinal to the fixed plane, and $`X^i`$, $`i=6,\mathrm{},9`$, transverse, the oscillators of the twisted sector have the following modding:
$`\text{twisted NS}:`$ $`\{\begin{array}{cc}\alpha _n^\alpha ,\hfill & n\mathrm{Z}\mathrm{Z}\hfill \\ \alpha _r^i,\hfill & r\mathrm{Z}\mathrm{Z}+1/2\hfill \end{array}\{\begin{array}{cc}\psi _r^\alpha ,\hfill & r\mathrm{Z}\mathrm{Z}+1/2\hfill \\ \psi _n^i,\hfill & n\mathrm{Z}\mathrm{Z}\hfill \end{array}`$ (2.5)
$`\text{twisted R}:`$ $`\{\begin{array}{cc}\alpha _n^\alpha ,\hfill & n\mathrm{Z}\mathrm{Z}\hfill \\ \alpha _r^i,\hfill & r\mathrm{Z}\mathrm{Z}+1/2\hfill \end{array}\{\begin{array}{cc}\psi _n^\alpha ,\hfill & n\mathrm{Z}\mathrm{Z}\hfill \\ \psi _r^i,\hfill & r\mathrm{Z}\mathrm{Z}+1/2\hfill \end{array}`$ (2.10)
There are 4 fermionic zero-modes in the NS sector in the directions $`i=6,7,8,9`$, and 6 fermionic zero-modes in the R-sector in the directions $`\alpha =0,1,\mathrm{},5`$. Since the intercept vanish for both sectors, the twisted groundstate will be given by these zero-modes. The twisted NS-NS sector gives a vector of $`SO(4)`$, or equivalently 4 scalars under $`SO(1,5)`$. On the other hand, the twisted R-R sector gives rise to a vector under $`SO(1,5)`$. From the point of view of the fixed plane, the supersymmetries are generated by two Weyl supercharges of different chirality. This corresponds to $`(1,1)`$ supersymmetry in 6 dimensions. In fact, the massless twisted sector can be identified with the worldvolume degrees of freedom of a NS-5 brane of type IIB .
### 2.1 The D-particle Boundary State
In this section we construct the boundary state for the stable non-BPS D-particle in type IIB orbifolded by $`(1)^{F_L}_4`$. This has been carried out in the light-cone gauge in . Here we use the covariant formalism for the boundary state, which has not been implemented before for non-BPS D-branes in orbifolds. The D-particle boundary state is made up of an untwisted NS-NS part and a twisted R-R part<sup>6</sup><sup>6</sup>6We use the subindex NS and R in the boundary states for short. :
$$|D0𝒾=|D0𝒾_{\mathrm{NS},\mathrm{U}}+|D0𝒾_{\mathrm{R},\mathrm{T}}.$$
(2.11)
The NS-NS boundary state is defined as the GSO-invariant combination of boundary states $`|D0,\eta _{\mathrm{NS},\mathrm{U}}`$, with $`\eta =\pm 1`$, which turns out to be :
$$|D0𝒾_{\mathrm{NS},\mathrm{U}}=𝒫_{\mathrm{GSO},\mathrm{U}}|D0,+_{\mathrm{NS},\mathrm{U}}=\frac{1}{2}(|D0,+𝒾_{\mathrm{NS},\mathrm{U}}|D0,𝒾_{\mathrm{NS},\mathrm{U}}),$$
(2.12)
where the GSO projector in the NS-NS sector is given by
$$𝒫_{\mathrm{GSO},\mathrm{U}}=\frac{1}{4}\left(1(1)^{F+G}\right)\left(1(1)^{\stackrel{~}{F}+\stackrel{~}{G}}\right),$$
(2.13)
with $`F`$ and $`G`$ the (worldsheet) fermion and superghost number operators, respectively:
$$F=\underset{m=1/2}{\overset{\mathrm{}}{}}\psi _m\psi _m,G=\underset{m=1/2}{\overset{\mathrm{}}{}}\left(\gamma _m\beta _m+\beta _m\gamma _m\right),$$
(2.14)
and $`\beta `$, $`\gamma `$ are the superghosts. For the right movers these operators are analogously defined. The state $`|D0,\eta _{\mathrm{NS},\mathrm{U}}`$ takes the form :
$$|D0,\eta _{\mathrm{NS},\mathrm{U}}=\frac{T_0}{2}|D0_X|D0_{\mathrm{gh}}|D0_\psi ,\eta _{\mathrm{NS}}|D0_{\mathrm{sgh}},\eta _{\mathrm{NS}},$$
(2.15)
where $`T_0`$, a constant related to the tension of the D-particle, is to be determined later. There is a bosonic and a fermionic part, ($`|D0_X`$ and $`|D0_\psi ,\eta _{\mathrm{NS}}`$), and also a ghost and a superghost part ($`|D0_{\mathrm{gh}}`$ and $`|D0_{\mathrm{sgh}},\eta _{\mathrm{NS}}`$). The boundary state $`|D0,\eta _{\mathrm{NS},\mathrm{U}}`$ is very similar to the NS-NS boundary state for type II branes. In fact, the only piece modified by the orbifold with respect to the type II case is the zero-mode part of $`|D0_X`$. For the bosonic untwisted boundary state the most general conditions invariant under the orbifold symmetry are
$`_\tau X^0|_{\tau =0}|D0_X`$ $`=`$ $`0,`$
$`X^p|_{\tau =0}|D0_X`$ $`=`$ $`y^p,p=1,\mathrm{},5,`$ (2.16)
$`X^i|_{\tau =0}|D0_X`$ $`=`$ $`0,i=6,\mathrm{},9,`$
from which we can deduce that
$$|D0_X=\delta ^{(5)}(\widehat{q}^py^p)\delta ^{(4)}(\widehat{q}^i)\mathrm{exp}\left[\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}\alpha _nS\stackrel{~}{\alpha }_n\right]|k=0.$$
(2.17)
Note that the D-particle position in the directions $`i=6,\mathrm{},9`$ is restricted by the orbifold symmetry to be on the fixed plane. The other pieces of the untwisted boundary state are not modified by the orbifold and are the same as for the type II case. These are explicitly given in Appendix B. Moreover, since there are nine Dirichlet directions for the case of the D-particle we have<sup>7</sup><sup>7</sup>7We take the Minkowski metric to be mostly plus. $`S_{\mu \nu }=\delta _{\mu \nu }`$.
The twisted R-R boundary state that we denote by $`|D0_{\mathrm{R},\mathrm{T}}`$ is defined as the GSO invariant combination of boundary states $`|D0,\eta _{\mathrm{R},\mathrm{T}}`$ and is given by :
$$|D0_{\mathrm{R},\mathrm{T}}=𝒫_{\mathrm{GSO},\mathrm{T}}|D0,+_{\mathrm{R},\mathrm{T}}=\frac{1}{2}\left(|D0,+_{\mathrm{R},\mathrm{T}}+|D0,_{\mathrm{R},\mathrm{T}}\right).$$
(2.18)
In this sector the GSO-operator is given by:
$$𝒫_{\mathrm{GSO},\mathrm{T}}=\frac{1}{4}\left(1+(1)^{F+G}\right)\left(1(1)^{\stackrel{~}{F}+\stackrel{~}{G}}\right),$$
(2.19)
with $`F`$ and $`G`$ the (worldsheet) fermion and superghost number operators in the twisted R-sector, respectively:
$`(1)^F`$ $`=`$ $`\mathrm{\Psi }(1)^{\underset{m=1}{\overset{\mathrm{}}{}}\psi _m^\alpha \eta _{\alpha \beta }\psi _m^\beta }(1)^{\underset{r=1/2}{\overset{\mathrm{}}{}}\psi _r^i\delta _{ij}\psi _r^j},`$
$`G`$ $`=`$ $`\gamma _0\beta _0{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}\left(\gamma _m\beta _m+\beta _m\gamma _m\right),`$ (2.20)
and similarly for the right movers. Here $`\mathrm{\Psi }`$ (and $`\stackrel{~}{\mathrm{\Psi }}`$) represent the zero-mode parts of the GSO-projectors, which are explicitly given in Appendix A. The twisted R-R boundary state $`|D0,\eta _{\mathrm{R},\mathrm{T}}`$ is given by:
$$|D0,\eta _{\mathrm{R},\mathrm{T}}=\frac{Q_0}{2}|D0_X_\mathrm{T}|D0_{\mathrm{gh}}|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}|D0_{\mathrm{sgh}},\eta _\mathrm{R},$$
(2.21)
where $`Q_0`$ is a normalisation factor related to the charge density of the brane and to be determined below. Notice that since this D-particle is non-BPS, $`Q_0T_0`$, unlike the BPS D-branes. The (super)ghosts are not affected by the orbifold , hence the corresponding pieces have the same form as for a type II R-R boundary state. Since in the twisted R-R sector the bosons have integer modding along the orbifold fixed plane and half-integer modding along the orbifolded directions, the state (2.21) have zero-modes along the fixed plane only. Accordingly, the twisted bosonic part takes the form:
$$|D0_X_\mathrm{T}=\delta ^{(5)}(\widehat{q}^py^p)\mathrm{exp}\left[\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}\alpha _n^\alpha \delta _{\alpha \beta }\stackrel{~}{\alpha }_n^\beta +\underset{n=1/2}{\overset{\mathrm{}}{}}\frac{1}{n}\alpha _n^i\delta _{ij}\stackrel{~}{\alpha }_n^j\right]|k=0.$$
(2.22)
The fermions in the twisted R-R sector have integer modding along the fixed plane and half-integer modding along the orbifolded directions. As a consequence (2.21) have fermionic zero-modes on the fixed plane only. The fermionic overlap conditions now read:
$`\left(\psi _m^\alpha i\eta S^\alpha {}_{\beta }{}^{}\stackrel{~}{\psi }_{m}^{\beta }\right)|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}`$ $`=`$ $`0,m\mathrm{Z}\mathrm{Z},`$
$`\left(\psi _r^ii\eta S^i{}_{j}{}^{}\stackrel{~}{\psi }_{r}^{j}\right)|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}`$ $`=`$ $`0,r\mathrm{Z}\mathrm{Z}+{\displaystyle \frac{1}{2}}.`$ (2.23)
From them we can deduce:
$$|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}=\mathrm{exp}\left[\mathrm{i}\eta \underset{m=1}{\overset{\mathrm{}}{}}\psi _m^\alpha \delta _{\alpha \beta }\stackrel{~}{\psi }_m^\beta \mathrm{i}\eta \underset{m=1/2}{\overset{\mathrm{}}{}}\psi _m^i\delta _{ij}\stackrel{~}{\psi }_m^j\right]|D0,\eta _{\mathrm{R},\mathrm{T}}^{(0)}.$$
(2.24)
The zero-mode part $`|D0,\eta _{\mathrm{R},\mathrm{T}}^{(0)}`$ is determined from the zero mode overlap conditions in (2.23) and is constructed upon the twisted R-R groundstate:
$$|D0,\eta _{\mathrm{R},\mathrm{T}}^{(0)}=M_{ab}|a_\mathrm{T}\stackrel{~}{|b}_\mathrm{T},$$
(2.25)
where $`a`$ and $`b`$ are spinor indices of $`SO(1,5)`$. The explicit form of $`M_{ab}`$ and its derivation are given in Appendix A.
Having constructed the covariant boundary state for the non-BPS D-particle, we determine next its tension and charge using the open-closed string consistency condition for boundary states. The interaction between two D-particles separated by a distance $`y`$ in the fixed plane is given by the amplitude between two boundary states located at a relative distance $`y`$ with respect to each other, with the insertion of a closed string propagator
$$D0|𝒟|D0,$$
(2.26)
with
$$𝒟_a=\frac{\alpha ^{}}{4\pi }_{|z|1}\frac{d^2z}{|z|^2}z^{L_0a}\overline{z}^{\stackrel{~}{L}_0a},$$
(2.27)
where $`a=1/2`$ in the untwisted NS-NS sector and $`a=0`$ in the twisted R-R sector. Moreover, one needs to take the appropriate moddings in the expression for $`L_0`$ and $`\stackrel{~}{L}_0`$ in the twisted sector. Also care must be taken in computing the matrix elements involving the zero modes of the superghosts and fermions in the twisted R-R sector. This is because, in general, the superghost zero modes produce infinite number of terms with any superghost number contribution, hence a regularisation is needed. We use the same regularisation as in . Defining $`|D0,\eta _{\mathrm{R},\mathrm{T}}^{(0)}=|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}^{(0)}|D0_{\mathrm{sgh}},\eta _\mathrm{R}^{(0)}`$ we give below the regularised result:
$${}_{\mathrm{R},\mathrm{T}}{}^{(0)}D0,\eta _1|D0,\eta _2_{\mathrm{R},\mathrm{T}}^{(0)}=\underset{\rho 1}{lim}{}_{\mathrm{R},\mathrm{T}}{}^{(0)}D0,\eta _1|\rho ^{2F_0+2G_0}|D0,\eta _2_{\mathrm{R},\mathrm{T}}^{(0)}=4\delta _{\eta _1\eta _2,1},$$
(2.28)
where $`F_0`$ and $`G_0`$ are the zero-mode parts of the operators $`F`$ and $`G`$, implicitly given in (2.20). Making a change of variables according to $`|z|=\mathrm{e}^\pi \mathrm{}`$ and $`d^2z=\pi \mathrm{e}^{2\pi \mathrm{}}d\mathrm{}d\varphi `$ we find:
$`D0|𝒟|D0`$ $`=`$ $`{\displaystyle \frac{V_1\alpha ^{}\pi }{16}}(2\pi ^2\alpha ^{})^{9/2}{\displaystyle _0^{\mathrm{}}}d\mathrm{}\mathrm{}^{9/2}\mathrm{e}^{\frac{y^2}{2\pi \alpha ^{}\mathrm{}}}\times `$
$`\times \left\{(T_0)^2{\displaystyle \frac{f_3^8(\mathrm{e}^\pi \mathrm{})f_4^8(\mathrm{e}^\pi \mathrm{})}{f_1^8(\mathrm{e}^\pi \mathrm{})}}(Q_0)^2(2\pi ^2\alpha ^{}\mathrm{})^2{\displaystyle \frac{f_2^4(\mathrm{e}^\pi \mathrm{})f_3^4(\mathrm{e}^\pi \mathrm{})}{f_1^4(\mathrm{e}^\pi \mathrm{})f_4^4(\mathrm{e}^\pi \mathrm{})}}\right\},`$
where $`V_1`$ is the (infinite) length of the D-particle worldline and $`f_i`$ are functions of $`\mathrm{e}^\pi \mathrm{}`$ defined in the usual manner. We can make a worldsheet transformation, $`\mathrm{}=1/\tau `$, to express the above closed string channel result in the open string channel. Note that we have considered both closed string and open string to have same periodicity in the spatial direction of the worldsheet. Open-closed string consistency then requires that the result must be equal to the 1-loop amplitude of the open strings stretched between the two D-particles. However, care must be taken to allow only states invariant under the orbifold projection to propagate in the loop . The open string states on a non-BPS D-particle are labelled by Chan-Paton factors<sup>8</sup><sup>8</sup>8This can be derived from a D0 anti-D0 system in type IIA following the prescription given in . 1l and $`\sigma _1`$. States with Chan-Paton factor 1l have usual GSO projection and are even under $`_4`$. On the other hand, states with Chan-Paton factor $`\sigma _1`$ have opposite GSO projection and are odd under $`_4`$. Moreover, the symmetry $`(1)^{F_L}`$ does not have any action on the open string oscillators. Accordingly, the open string 1-loop amplitude is given by:
$`𝒜`$ $`=`$ $`2V_1{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\tau }{2\tau }}\mathrm{Tr}_{\mathrm{NS}\mathrm{R}}\{{\displaystyle \frac{1}{4}}(1(1)^{F+G})(1+_4)e^{2\pi \tau L_0}`$ (2.30)
$`+{\displaystyle \frac{1}{4}}(1+(1)^{F+G})(1_4)e^{2\pi \tau L_0}\}`$
$`=`$ $`2V_1{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\tau }{2\tau }}\mathrm{Tr}_{\mathrm{NS}\mathrm{R}}\left\{{\displaystyle \frac{1}{2}}\left(1(1)^{F+G}_4\right)e^{2\pi \tau L_0}\right\},`$
where $`F`$ and $`G`$ are the (worldsheet) fermion and superghost number operators for the open string. The tachyon is projected out in the trace and this renders the non-BPS D-particle stable . We obtain the following 1-loop amplitude:
$`𝒜`$ $`=`$ $`{\displaystyle \frac{1}{2}}V_1(8\pi ^2\alpha ^{})^{1/2}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\tau }{\tau ^{3/2}}}\mathrm{e}^{\frac{y^2\tau }{2\pi \alpha ^{}}}\times `$
$`\times \left\{{\displaystyle \frac{f_3^8(\mathrm{e}^{\pi \tau })f_2^8(\mathrm{e}^{\pi \tau })}{f_1^8(\mathrm{e}^{\pi \tau })}}4{\displaystyle \frac{f_4^4(\mathrm{e}^{\pi \tau })f_3^4(\mathrm{e}^{\pi \tau })}{f_1^4(\mathrm{e}^{\pi \tau })f_2^4(\mathrm{e}^{\pi \tau })}}\right\}.`$
As mentioned above, open-closed string consistency allows us to fix the normalisation of the boundary state:
$$T_0=8(\alpha ^{})^{3/2}\pi ^{7/2},Q_0=8\pi \sqrt{\pi \alpha ^{}}.$$
(2.32)
The tension of the non-BPS D-particle is given by $`T_0`$ and $`Q_0`$ is related to its electric charge. Using that in ten dimensions $`\kappa _{10}=8\pi ^{7/2}g_s(\alpha ^{})^2`$, and that for the orbifold $`\kappa _{orb}=\sqrt{2}\kappa _{10}`$, we find the mass per unit volume of the D-particle:
$$M_0=\frac{T_0}{\kappa _{orb}}=\frac{1}{g\sqrt{2\alpha ^{}}},$$
(2.33)
which agrees with the mass of the D-particle found in . In analogy with the BPS D-branes, we can define an electric charge with respect to the (twisted) R-R field as $`\mu _0=\sqrt{2}Q_0=8\pi \sqrt{2\pi \alpha ^{}}`$. It is interesting to observe that the tension $`T_0`$ of this non-BPS D-particle is the same as that of a BPS D-particle of type IIA theory in ten dimensions<sup>9</sup><sup>9</sup>9The unstable non-BPS D-particle of type IIB has a tension $`\sqrt{2}`$ times bigger than the type IIA BPS D-particle. However, in the orbifold case, there is an extra factor $`1/2`$ in the open string amplitude coming from the projection operator, such that the $`\sqrt{2}`$ factor is compensated.. However, their charges are different. In fact, $`\mu _0`$ is exactly twice the charge of the BPS D-particle of type IIA in six dimensions.
### 2.2 The D-particle in the Compact Orbifold
Let us consider now the boundary state of a non-BPS D-particle in the case of the compact orbifold, i.e. type IIB on $`T^4/`$$`(1)^{F_L}_4`$. The coordinates $`x^i`$, $`i=6,7,8,9`$, are periodic and there are 16 fixed planes instead of one, located at $`x^i=0,\pi R_i`$. We consider a non-BPS D-particle located on one of the fixed planes, that we choose $`x^i=0`$, $`i=6,7,8,9`$, for simplicity.
The boundary state is now constructed upon a groundstate with zero winding and zero momentum. The only piece modified in the boundary state with respect to the non-compact case is the bosonic part and the overall normalisation. These two modifications will only affect the NS-NS part of the boundary state. The bosonic part is modified to include Kaluza-Klein modes in the directions $`i=6,\mathrm{},9`$:
$$|D0_X=\delta ^{(5)}(\widehat{q}^py^p)\underset{i=6}{\overset{9}{}}\underset{n_i\mathrm{Z}\mathrm{Z}}{}\mathrm{e}^{i\widehat{q}_{n}^{}{}_{}{}^{i}\frac{n_i}{R_i}}\mathrm{exp}\left[\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}\alpha _nS\stackrel{~}{\alpha }_n\right]|k=0,n=0.$$
(2.34)
The states $`|n_i=\mathrm{exp}(iq_n^i\frac{n_i}{R_i})|n=0`$ are normalised as
$$n_i|n_i^{}=\mathrm{\Phi }_i\delta _{n_in_i^{}},$$
(2.35)
where $`\mathrm{\Phi }_i`$ is the self-dual volume which satisfies
$$\underset{R_i\mathrm{}}{lim}\mathrm{\Phi }_i=2\pi R_i,\underset{R_i0}{lim}\mathrm{\Phi }_i=\frac{2\pi \alpha ^{}}{R_i}.$$
(2.36)
We consider the new overall normalisation factor as a constant times the normalisation for the uncompactified case:
$$|D0,\eta _{\mathrm{NS},\mathrm{U}}=\frac{T_0}{2}𝒩|D0_X|D0_{\mathrm{gh}}|D0_\psi ,\eta _{\mathrm{NS}}|D0_{\mathrm{sgh}},\eta _{\mathrm{NS}},$$
(2.37)
The factor $`𝒩`$ can be obtained by open-closed string consistency and is such that in the decompactification limit $`R_i\mathrm{}`$, we recover the non-compact boundary state . The amplitude between two non-BPS D-particles in the closed string channel is given by:
$`D0|𝒟|D0`$ $`=`$ $`{\displaystyle \frac{V_1\alpha ^{}\pi }{16}}(2\pi ^2\alpha ^{})^{5/2}{\displaystyle _0^{\mathrm{}}}d\mathrm{}\mathrm{}^{5/2}\mathrm{e}^{\frac{y^2}{2\pi \alpha ^{}\mathrm{}}}\times `$
$`\times \{(T_0)^2𝒩^2\left({\displaystyle \underset{i=6}{\overset{9}{}}}\mathrm{\Phi }_i\right){\displaystyle \underset{i=6}{\overset{9}{}}}{\displaystyle \underset{n_i\mathrm{Z}\mathrm{Z}}{}}\mathrm{e}^{\frac{\pi \alpha ^{}\mathrm{}}{2}(\frac{n_i}{R_i})^2}{\displaystyle \frac{f_3^8(\mathrm{e}^\pi \mathrm{})f_4^8(\mathrm{e}^\pi \mathrm{})}{f_1^8(\mathrm{e}^\pi \mathrm{})}}`$
$`(Q_0)^2{\displaystyle \frac{f_2^4(\mathrm{e}^\pi \mathrm{})f_3^4(\mathrm{e}^\pi \mathrm{})}{f_1^4(\mathrm{e}^\pi \mathrm{})f_4^4(\mathrm{e}^\pi \mathrm{})}}\}.`$
Using the worldsheet duality $`\mathrm{}=\frac{1}{\tau }`$ and the Poisson resummation formula
$$\underset{n\mathrm{Z}\mathrm{Z}}{}\mathrm{e}^{\frac{\pi \alpha ^{}\mathrm{}}{2}(\frac{n}{R})^2}=\sqrt{\frac{2}{\mathrm{}\alpha ^{}}}R\underset{m\mathrm{Z}\mathrm{Z}}{}\mathrm{e}^{\frac{2\pi }{\alpha ^{}\mathrm{}}(mR)^2},$$
(2.39)
one obtains
$`D0|𝒟|D0`$ $`=`$ $`{\displaystyle \frac{V_1\alpha ^{}\pi }{16}}(2\pi ^2\alpha ^{})^{5/2}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\tau }{\tau ^{3/2}}}\mathrm{e}^{\frac{y^2\tau }{2\pi \alpha ^{}}}\times `$
$`\times \{(T_0)^2𝒩^2\left({\displaystyle \frac{2}{\alpha ^{}}}\right)^2\left({\displaystyle \underset{i=6}{\overset{9}{}}}R_i\mathrm{\Phi }_i\right){\displaystyle \underset{i=6}{\overset{9}{}}}{\displaystyle \underset{m_i\mathrm{Z}\mathrm{Z}}{}}\mathrm{e}^{\frac{2\pi \tau }{\alpha ^{}}(m_iR_i)^2}{\displaystyle \frac{f_3^8(\mathrm{e}^{\pi \tau })f_2^8(\mathrm{e}^{\pi \tau })}{f_1^8(\mathrm{e}^{\pi \tau })}}`$
$`(Q_0)^2{\displaystyle \frac{f_4^4(\mathrm{e}^{\pi \tau })f_3^4(\mathrm{e}^{\pi \tau })}{f_1^4(\mathrm{e}^{\pi \tau })f_2^4(\mathrm{e}^{\pi \tau })}}\}.`$
Open-closed string consistency imposes that this amplitude must be equal to the open string 1-loop amplitude for the compactified case
$`𝒜`$ $`=`$ $`2V_1{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\tau }{2\tau }}\mathrm{Tr}_{\mathrm{NS}\mathrm{R}}\left\{{\displaystyle \frac{1}{2}}\left(1+(1)^F_4\right)\mathrm{e}^{2\pi \tau L_0}\right\}`$
$`=`$ $`{\displaystyle \frac{1}{2}}V_1(8\pi ^2\alpha ^{})^{1/2}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\tau }{\tau ^{3/2}}}\mathrm{e}^{\frac{y^2\tau }{2\pi \alpha ^{}}}\times `$
$`\times `$ $`\left\{{\displaystyle \underset{i=6}{\overset{9}{}}}{\displaystyle \underset{m_i\mathrm{Z}\mathrm{Z}}{}}\mathrm{e}^{\frac{2\pi \tau }{\alpha ^{}}(m_iR_i)^2}{\displaystyle \frac{f_3^8(\mathrm{e}^{\pi \tau })f_2^8(\mathrm{e}^{\pi \tau })}{f_1^8(\mathrm{e}^{\pi \tau })}}4{\displaystyle \frac{f_4^4(\mathrm{e}^{\pi \tau })f_3^4(\mathrm{e}^{\pi \tau })}{f_1^4(\mathrm{e}^{\pi \tau })f_2^4(\mathrm{e}^{\pi \tau })}}\right\}.`$
This gives the same value for $`Q_0`$ as in the uncompactified case, which is expected since the compactification is done in the transverse directions of the D-particle. This also fixes the normalisation factor $`𝒩`$ to be:
$$𝒩=\left(\underset{i=6}{\overset{9}{}}2\pi R_i\mathrm{\Phi }_i\right)^{1/2}.$$
(2.42)
It is easy to check that with this normalisation and using $`\mathrm{\Phi }_i=2\pi R_i`$, in the decompactified limit, we can recover the untwisted boundary state for the non-compact case. Thus the net effect of the compactification is a renormalisation of the tension of the D-particle and as expected, the tension depends on the compactification radii. At this point we can recover the vanishing of the amplitude at the critical radius . The open strings on the D-particle have now winding modes:
$$M^2=\underset{i=6}{\overset{9}{}}\left(\frac{w_iR_i}{\alpha ^{}}\right)^2+\frac{1}{\alpha ^{}}\left(N\frac{1}{2}\right).$$
(2.43)
The groundstate with zero winding is the tachyon and is projected out by the orbifold symmetry. At a particular critical value of the radii
$$R_i=\sqrt{\frac{\alpha ^{}}{2}},i=6,7,8,9,$$
(2.44)
there are four states, for which only one $`w_i0`$, which become massless . These are the modes of the tachyon field that at the critical radius correspond to the marginal deformation which takes the D-particle to a bound state of a D-string and an anti-D-string . At this critical radius a Bose-Fermi degeneracy takes place, which translates to the fact that the 1-loop amplitude vanishes . In the closed string channel this can be seen by using the relations:
$`{\displaystyle \underset{n\mathrm{Z}\mathrm{Z}}{}}\mathrm{e}^{\pi \mathrm{}n^2}`$ $`=`$ $`f_1(\mathrm{e}^\pi \mathrm{})f_3^2(\mathrm{e}^\pi \mathrm{})`$ (2.45)
$`f_4(\mathrm{e}^\pi \mathrm{})f_2(\mathrm{e}^\pi \mathrm{})f_3(\mathrm{e}^\pi \mathrm{})`$ $`=`$ $`\sqrt{2}.`$
Moreover, the normalisation factor $`𝒩`$ at the critical radius becomes
$$𝒩=\frac{1}{2\pi ^2\alpha ^{}}\left(\underset{i=6}{\overset{9}{}}\mathrm{\Phi }_i\right)^{1/2}.$$
(2.46)
Hence we can write the amplitude as
$`D0|𝒟|D0`$ $`=`$ $`{\displaystyle \frac{V_1\alpha ^{}\pi }{16}}(2\pi ^2\alpha ^{})^{5/2}{\displaystyle _0^{\mathrm{}}}d\mathrm{}\mathrm{}^{5/2}\mathrm{e}^{\frac{y^2}{2\pi \alpha ^{}\mathrm{}}}\times `$ (2.47)
$`\times {\displaystyle \frac{f_3^4(\mathrm{e}^\pi \mathrm{})}{f_1^4(\mathrm{e}^\pi \mathrm{})f_2^4(\mathrm{e}^\pi \mathrm{})f_4^4(\mathrm{e}^\pi \mathrm{})}}\{\left({\displaystyle \frac{T_0}{\pi ^2\alpha ^{}}}\right)^2(f_3^8(\mathrm{e}^\pi \mathrm{})f_4^8(\mathrm{e}^\pi \mathrm{}))`$
$`(Q_0)^2f_2^8(\mathrm{e}^\pi \mathrm{})\}.`$
which vanishes, at any distance $`y`$, by the abstruse identity and using the expressions for $`T_0`$ and $`Q_0`$. In the open string channel it is simpler to demonstrate this by using the above identities.
### 2.3 Long and Short Distance Interactions
In this Section we consider the interaction amplitude for non-BPS D-particles in relative motion. This will give a velocity dependent potential which we will compare with the usual BPS case. This amplitude can be obtained by using a boosted boundary state for the non-BPS D-particle. For definiteness we consider a D-particle moving in the direction $`X^1`$ with a velocity $`v`$, interacting with another non-BPS D-particle at rest. We consider first the non-compact orbifold. The boosted boundary state include some $`v`$-dependent modifications that we write explicitly below. In the NS-NS untwisted sector, the bosonic part of the boundary state becomes
$$|D0_X,v=\sqrt{1v^2}\left(\underset{i=2}{\overset{9}{}}\delta (\widehat{q}^p)\right)\delta (\widehat{q}^1+v\widehat{q}^0)\mathrm{exp}\left[\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}\alpha _nS(v)\stackrel{~}{\alpha }_n\right]|k=0,$$
(2.48)
where the matrix $`S(v)`$ is given by
$$S_{00}(v)=S_{11}(v)=\frac{1+v^2}{1v^2},S_{10}(v)=S_{01}(v)=\frac{2v}{1v^2},$$
(2.49)
and $`S(v)=S`$ for the other components. The other pieces in the NS-NS untwisted boundary state have the usual form except for a substitution of the matrix $`S`$ by $`S(v)`$. Similarly in the R-R twisted sector, the boosted bosonic boundary state takes the form:
$$|D0_X,v_\mathrm{T}=\sqrt{1v^2}\underset{p=2}{\overset{5}{}}\delta (\widehat{q}^p)\delta (\widehat{q}^1+v\widehat{q}^0)\mathrm{exp}\left[\underset{t.m.}{}\frac{1}{n}\alpha _nS(v)\stackrel{~}{\alpha }_n\right]|k=0,n=0,$$
where with $`t.m.`$ we indicate that one takes the corresponding twisted moddings of the twisted sector. For the other pieces we do the same substitution $`SS(v)`$. Finally, for the R-R zero-mode part there is as well a $`v`$-dependent modification:
$$|D0,\eta ,v_{\mathrm{R},\mathrm{T}}^{(0)}=M_{ab}(v)|a_\mathrm{T}\stackrel{~}{|b}_\mathrm{T},$$
(2.50)
and the matrix $`M_{ab}(v)`$ is given in Appendix A. We define the distance between the particles as $`r^2=b^2+t^2v^2(1v^2)^1`$, where $`t`$ is the proper time of the moving particle along which we also integrate; $`b`$ is the impact parameter: $`b^2=y_2^2+\mathrm{}y_5^2`$. For convenience we also define the following variable
$$u=\frac{1}{2\pi i}\mathrm{ln}\frac{1v}{1+v}.$$
(2.51)
Finally, the cylinder amplitude takes the following form:
$`D0,v|𝒟|D0`$ $`=`$ $`(8\pi ^2\alpha ^{})^{1/2}\mathrm{sin}\pi u{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}d\mathrm{}\mathrm{}^{\frac{9}{2}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}dt\mathrm{e}^{\frac{r^2}{2\pi \alpha ^{}\mathrm{}}}\times `$
$`\left\{{\displaystyle \frac{\mathrm{\Theta }_3(u,i\mathrm{})f_3^6(\mathrm{e}^\pi \mathrm{})\mathrm{\Theta }_4(u,i\mathrm{})f_4^6(\mathrm{e}^\pi \mathrm{})}{\mathrm{\Theta }_1(u,i\mathrm{})f_1^6(\mathrm{e}^\pi \mathrm{})}}\mathrm{\hspace{0.17em}4}\mathrm{}^2{\displaystyle \frac{\mathrm{\Theta }_2(u,i\mathrm{})f_2^2(\mathrm{e}^\pi \mathrm{})f_3^4(\mathrm{e}^\pi \mathrm{})}{\mathrm{\Theta }_1(u,i\mathrm{})f_1^2(\mathrm{e}^\pi \mathrm{})f_4^4(\mathrm{e}^\pi \mathrm{})}}\right\},`$
where $`\mathrm{\Theta }_s(u,i\mathrm{})`$ are the Jacobi $`\mathrm{\Theta }`$-functions, and we have used the value of the tension and charge of the D-particle. For $`v=0`$ we recover the amplitude for the static interaction given in (2.1).
In order to study the interaction we define the interacting potential of the scattering $`𝒰`$, in the following way:
$$D0,v|𝒟|D0=\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑t𝒰(v,r(t)).$$
(2.53)
We can extract the long range interaction potential taking the limit $`\mathrm{}\mathrm{}`$ in the integrand of (2.3), and then performing the integral in $`\mathrm{}`$. For slow velocities, we find the following expansion in powers of $`v`$:
$$𝒰_{closed}\frac{(2\pi \alpha ^{})^3}{(4\pi )^{1/2}}\mathrm{\hspace{0.17em}8}\left\{(1+\frac{1}{2}v^2+\frac{1}{24}v^4)\left(\frac{\mathrm{\Gamma }(\frac{7}{2})}{r^7}\frac{1}{(2\pi \alpha ^{})^2}\frac{\mathrm{\Gamma }(\frac{3}{2})}{r^3}\right)+\frac{1}{8}v^4\frac{\mathrm{\Gamma }(\frac{7}{2})}{r^7}\right\},$$
(2.54)
to order $`v^4`$. There is a static force, as expected since the D-particle is non-BPS. Moreover, the velocity dependent corrections start at order $`v^2`$, as for the potential between two BPS D-branes of different dimensionality , and as for the non-BPS D-particle of type I . This is something we could expect from the fact that the particle breaks all supersymmetries.
We can analyse this amplitude for very short distances. In this region the open string description dominates. Using the modular properties of the $`\mathrm{\Theta }`$-functions we can write (2.3) in the open string channel:
$`D0,v|𝒟|D0`$ $`=`$ $`i(8\pi ^2\alpha ^{})^{1/2}\mathrm{sin}\pi u{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}d\tau \tau ^{\frac{1}{2}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}dt\mathrm{e}^{\frac{r^2}{2\pi \alpha ^{}}\tau }\times `$
$`\left\{{\displaystyle \frac{\mathrm{\Theta }_3(iu\tau ,i\tau )f_3^6(\mathrm{e}^{\pi \tau })\mathrm{\Theta }_2(iu\tau ,i\tau )f_2^6(\mathrm{e}^{\pi \tau })}{\mathrm{\Theta }_1(iu\tau ,i\tau )f_1^6(\mathrm{e}^{\pi \tau })}}\mathrm{\hspace{0.17em}4}{\displaystyle \frac{\mathrm{\Theta }_4(iu\tau ,i\tau )f_4^2(\mathrm{e}^{\pi \tau })f_3^4(\mathrm{e}^{\pi \tau })}{\mathrm{\Theta }_1(iu\tau ,i\tau )f_1^2(\mathrm{e}^{\pi \tau })f_2^4(\mathrm{e}^{\pi \tau })}}\right\}.`$
We can derive the interaction potential at short distances by taking the limit $`\tau \mathrm{}`$ in the integrand and then performing the integral in $`\tau `$. Note that the pieces originally from the NS-NS and R-R sectors multiply now in (2.3) the same powers of $`\tau `$, hence we can expect that the interaction potential will be very different from the closed string case. This time we obtain:
$$𝒰_{open}=i(8\pi ^2\alpha ^{})^{1/2}\mathrm{sin}\pi u\underset{0}{\overset{\mathrm{}}{}}𝑑\tau \tau ^{1/2}\mathrm{e}^{\frac{r^2}{2\pi \alpha ^{}}\tau }\frac{2+2\mathrm{c}\mathrm{o}\mathrm{s}(2\pi iu\tau )8\mathrm{c}\mathrm{o}\mathrm{s}(i\pi u\tau )}{\mathrm{sin}(i\pi u\tau )},$$
(2.56)
which for slow velocities and after integrating in $`\tau `$ becomes (to order $`v^4`$):
$$𝒰_{open}=\frac{4}{2\pi \alpha ^{}}r+\frac{(2\pi \alpha ^{})^3\mathrm{\Gamma }(\frac{7}{2})}{(4\pi )^{1/2}}\frac{v^4}{r^7}.$$
(2.57)
Surprisingly, there is no $`v^2`$ correction to the potential, as for BPS D-branes. The first term is the linear repulsive force coming from the stretched strings. Moreover, the $`v^4`$ correction has precisely the same form as the $`v^4`$ term for a scattering of two BPS D-particles. For BPS branes the $`v^4`$ term in the open string channel coincides with the $`v^4`$ term from the closed string description . This does not occur in the non-compact orbifold. As we will see below for the compact case, even though at the critical radius the static potential and $`v^2`$ terms will be suppressed, the matching of the $`v^4`$ terms will not occur either. Hence the BPS-like behaviour will not extend beyond order $`v^2`$.
We extend next the analysis of the cylinder amplitude for the case of relative motion of the D-particles in the compact orbifold. We consider again one non-BPS D-particle moving in the direction $`X^1`$ with a velocity $`v`$, interacting with another non-BPS D-particle at rest. The boosted boundary state include the $`v`$-dependent modifications on top of the modifications due to compactification. In the NS-NS untwisted sector, the bosonic part of the boundary state becomes
$`|D0_X,v`$ $`=`$ $`\sqrt{1v^2}\left({\displaystyle \underset{p=2}{\overset{5}{}}}\delta (\widehat{q}^p)\right)\delta (\widehat{q}^1+v\widehat{q}^0)\times `$ (2.59)
$`\times \left({\displaystyle \underset{i=6}{\overset{9}{}}}{\displaystyle \underset{n_i\mathrm{Z}\mathrm{Z}}{}}\mathrm{e}^{i\widehat{q}_n^i\frac{n^i}{R_i}}\right)\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}\alpha _nS(v)\stackrel{~}{\alpha }_n\right]|k=0,n=0,`$
where the matrix $`S(v)`$ is as given before in (2.49). In the R-R twisted sector, the boosted bosonic boundary state and fermionic zero-mode have the same form as for the non-compact case, and the other pieces have the usual form except for a substitution of the matrix $`S`$ by $`S(v)`$. Using the same notation as before, the cylinder amplitude takes the following form:
$`D0,v|𝒟|D0`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi ^2\alpha ^{}}}}\mathrm{sin}\pi u{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}d\mathrm{}\mathrm{}^{\frac{5}{2}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}dt\mathrm{e}^{\frac{r^2}{2\pi \alpha ^{}\mathrm{}}}\times `$
$`\{\left({\displaystyle \frac{\alpha ^{}}{4}}\right)^2{\displaystyle \underset{i=6}{\overset{9}{}}}R_i^1\left({\displaystyle \underset{i=6}{\overset{9}{}}}{\displaystyle \underset{n_i\mathrm{Z}\mathrm{Z}}{}}\mathrm{e}^{\frac{\pi \alpha ^{}\mathrm{}}{2}(\frac{n_i}{R_i})^2}\right)\left({\displaystyle \frac{\mathrm{\Theta }_3(u,i\mathrm{})f_3^6(\mathrm{e}^\pi \mathrm{})\mathrm{\Theta }_4(u,i\mathrm{})f_4^6(\mathrm{e}^\pi \mathrm{})}{\mathrm{\Theta }_1(u,i\mathrm{})f_1^6(\mathrm{e}^\pi \mathrm{})}}\right)`$
$`{\displaystyle \frac{\mathrm{\Theta }_2(u,i\mathrm{})f_2^2(\mathrm{e}^\pi \mathrm{})f_3^4(\mathrm{e}^\pi \mathrm{})}{\mathrm{\Theta }_1(u,i\mathrm{})f_1^2(\mathrm{e}^\pi \mathrm{})f_4^4(\mathrm{e}^\pi \mathrm{})}}\}.`$
For $`v=0`$ we recover the amplitude for the static interaction given in (2.2). For slow velocities ($`u(i\pi )^1v`$) and after integration in $`\mathrm{}`$, the long range interaction potential to order $`v^4`$ is given by
$$𝒰_{closed}\frac{4\pi \alpha ^{}}{r^3}\left\{(1+\frac{1}{2}v^2+\frac{1}{6}v^4)\left(\frac{\alpha ^2}{4}\underset{i=6}{\overset{9}{}}R_i^11\right)\frac{v^4}{8}\right\}.$$
(2.61)
We see that for generic radii the potential has $`v^2`$ corrections, as before, which is typical for potentials between two BPS D-branes of different dimensionality , and occurs also for the non-BPS D-particle of type I . On the other hand, at the critical radius they start at $`v^4`$, as for BPS D-branes . Notice that the static and $`v^2`$ terms of the potential also vanishes for other values of the radii. However, those radii do not make the amplitude (2.2) vanish. Moreover, this would require some of the radii to be below the critical radius, and the D-particle would not be stable.
In order to study the short-distance behaviour we write the scattering amplitude in the open channel:
$`D0,v|𝒟|D0`$ $`=`$ $`i\sqrt{{\displaystyle \frac{2}{\pi ^2\alpha ^{}}}}\mathrm{sin}\pi u{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}d\tau \tau ^{\frac{1}{2}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}dt\mathrm{e}^{\frac{r^2}{2\pi \alpha ^{}}\tau }\times `$
$`\{{\displaystyle \frac{1}{4}}\left({\displaystyle \underset{i=6}{\overset{9}{}}}{\displaystyle \underset{m_i\mathrm{Z}\mathrm{Z}}{}}\mathrm{e}^{\frac{2\pi \tau }{\alpha ^{}}(m_iR_i)^2}\right)\left({\displaystyle \frac{\mathrm{\Theta }_3(iu\tau ,i\tau )f_3^6(\mathrm{e}^{\pi \tau })\mathrm{\Theta }_2(iu\tau ,i\tau )f_2^6(\mathrm{e}^{\pi \tau })}{\mathrm{\Theta }_1(iu\tau ,i\tau )f_1^6(\mathrm{e}^{\pi \tau })}}\right)`$
$`{\displaystyle \frac{\mathrm{\Theta }_4(iu\tau ,i\tau )f_4^2(\mathrm{e}^{\pi \tau })f_3^4(\mathrm{e}^{\pi \tau })}{\mathrm{\Theta }_1(iu\tau ,i\tau )f_1^2(\mathrm{e}^{\pi \tau })f_2^4(\mathrm{e}^{\pi \tau })}}\}.`$
To extract the leading contribution at short distances and low velocities we take the limit $`\tau \mathrm{}`$ in the integrand. This time one must also include contributions with one unit of winding number $`m_i=1`$, which are relevant very close to the critical radius $`R_iR_c`$:
$`𝒰_{open}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi ^2\alpha ^{}}}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}d\tau \tau ^{1/2}\mathrm{e}^{\frac{r^2}{2\pi \alpha ^{}}\tau }{\displaystyle \frac{\mathrm{sin}\pi u}{2i\mathrm{sin}(i\pi u\tau )}}\times `$
$`\times \left(\mathrm{cos}(2\pi iu\tau )4\mathrm{c}\mathrm{o}\mathrm{s}(i\pi u\tau )+1+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=6}{\overset{9}{}}}\mathrm{e}^{\frac{2\pi R_i^2}{\alpha ^{}}\tau +\pi \tau }\right).`$
At low velocities the expression simplify. Carrying out the integral in $`\tau `$ we find:
$$𝒰_{open}=\frac{4}{2\pi \alpha ^{}}\left(r\frac{1}{4}\underset{i=6}{\overset{9}{}}\sqrt{r^2+4\pi ^2(R_i^2\frac{\alpha ^{}}{2})}\right)+\frac{(2\pi \alpha ^{})^3\mathrm{\Gamma }(\frac{7}{2})}{(4\pi )^{1/2}}\frac{v^4}{r^7}.$$
(2.64)
As for the non-compact case we find no $`v^2`$ corrections, and the $`v^4`$ corrections are again the same as for the BPS D-particle in ten dimensions. Notice also that the $`r`$-dependence of the $`v^4`$ term is as for the uncompactified case. The static part of the potential vanishes when all the radii are equal to the critical radius. Remarkably, at the critical radius $`𝒰_{open}`$ takes a very BPS-like form. However, as we announced before, the $`v^4`$ terms does not match with the closed string result.
From this analysis we conclude that the long and short range interactions of the non-BPS D-particle are quite different for generic radii of the orbifold. In the particular case of the critical radius the static and $`v^2`$ terms of the interaction potential are absent in the open and closed string description. At the critical radius, to order $`v^2`$, we expect an equivalence between the (super)gravity and worldvolume description of the non-BPS D-particle. At order $`v^4`$ the open strings begin to describe the geometry of non-BPS D-particle very differently from the closed strings.
## 3 Spacetime Description
As shown in , the boundary state of a D-brane encodes the long distance behaviour of the corresponding classical solution of supergravity. At large distances, this classical solution tends to a flat background configuration. The fluctuations around this background is exactly reproduced by the boundary state. From the NS-NS part one obtains the asymptotic behaviour of the dilaton and the metric, and from the R-R part one obtains the asymptotic behaviour of the R-R form potential under which the brane is charged. In this section we implement this technique on the stable non-BPS D-particle to obtain the asymptotic form of the solution. We obtain a metric and a dilaton in the bulk, which depend on all the coordinates transverse to the D-particle, and whose dependence is the expected one for a particle in 10 dimensions. On the other hand, we also find a twisted R-R 1-form which is restricted to the fixed plane and whose dependence on the spatial coordinates on the fixed plane is the correct one for a particle in 6 dimensions. We describe first the non-compact case.
### 3.1 The Non-compact Case
Given a certain massless closed string state $`|\phi `$, normalised as $`\phi |\phi =\phi ^2`$, one can define a projection operator $`P_{(\phi )}|`$ associated to this state, such that $`P_{(\phi )}|\phi =\phi `$ . The asymptotic behaviour of the classical field $`\phi `$, generated by a D-brane is determined by computing the overlap of the boundary state with the state $`|P_{(\phi )}`$ with the insertion of a closed string propagator :
$$P_{(\phi )}|𝒟|D.$$
(3.1)
Since the D-particle has an untwisted NS-NS part and a twisted R-R part, the D-particle will excite gravitons $`g_{\mu \nu }`$ and dilatons $`\varphi `$ in the untwisted sector and a R-R 1-form field $`𝒞^{(1)}`$ in the twisted sector. At large distances, these fields will be of the form of a fluctuation around a flat background<sup>10</sup><sup>10</sup>10We use the same normalisation as in .:
$$g_{\mu \nu }\eta _{\mu \nu }+2\kappa _{10}\delta h_{\mu \nu },\varphi \kappa _{10}\sqrt{2}\delta \varphi ,𝒞^{(1)}\kappa _6\sqrt{2}\delta C^{(1)}.$$
(3.2)
The fluctuations $`\delta h_{\mu \nu }`$, $`\delta \varphi `$ and $`\delta C^{(1)}`$, are given by (3.1) using the state $`|D0`$ constructed previously. We consider first the case of the uncompactified orbifold. In the untwisted sector, only those which are invariant under the orbifold symmetry will appear in the theory. This implies that string states which depend on the momentum in the orbifolded directions are only invariant if they are symmetric with respect to $`_4:`$ $`k^ik^i`$, $`i=6,7,8,9`$. In particular, these will include gravitons that propagate off the orbifold fixed plane which are symmetric in $`k^i`$.
The projection operators in the untwisted NS-NS sector are given in the $`(1,1)`$ picture. For simplicity, we use the following notation for the NS-NS groundstate in the $`(1,1)`$ picture:
$$|k|k/2_1\stackrel{~}{|k/2}_1.$$
(3.3)
The projectors for the dilaton and the metric are given by:
$`P_{(\varphi )}|`$ $`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}}}(k_{},k^i|+k_{},k^i|)\psi _{1/2}^\nu \stackrel{~}{\psi }_{1/2}^\mu (\eta _{\mu \nu }k_\mu \mathrm{}_\nu k_\nu \mathrm{}_\mu ),`$
$`P_{(h)}^{\mu \nu }|`$ $`=`$ $`{\displaystyle \frac{1}{4}}(k_{},k^i|+k_{},k^i|)(\psi _{1/2}^\nu \stackrel{~}{\psi }_{1/2}^\mu +\psi _{1/2}^\mu \stackrel{~}{\psi }_{1/2}^\nu )`$
$`P_{(\varphi )}|{\displaystyle \frac{1}{2\sqrt{2}}}\left(\eta ^{\mu \nu }k^\mu \mathrm{}^\nu k^\nu \mathrm{}^\mu \right),`$
where $`k=(k_{},k^i)`$, with $`k_{}`$ the spatial components of the momentum along the fixed plane directions. Moreover, $`\mathrm{}`$ is such that $`\mathrm{}^2=0`$ and $`k\mathrm{}=1`$.
In the R-R sector, the most natural picture is the asymmetric<sup>11</sup><sup>11</sup>11See for an explanation about this. picture $`(1/2,3/2)`$, where the R-R vertex operator couples to the potential instead of the field strength . In analogy to the type II case described in and using $`a,b`$ as the SO(1,5) spinor indices, we write the following quantum state in the $`(1/2,3/2)`$ picture associated to the twisted R-R 1-form field<sup>12</sup><sup>12</sup>12For simplicity we omit the subindex T of the twisted groundstate.
$$|C^{(1)}=\frac{1}{\sqrt{2}}C_\alpha ^{(1)}\left\{\left(𝒞\gamma ^\alpha \mathrm{\Pi }_+\right)_{ab}\mathrm{cos}\gamma _0\stackrel{~}{\beta }_0+\left(𝒞\gamma ^\alpha \mathrm{\Pi }_{}\right)_{ab}\mathrm{sin}\gamma _0\stackrel{~}{\beta }_0\right\}|a,k/2_{1/2}\stackrel{~}{|b,k/2}_{3/2},$$
whose conjugate state is given by:
$$C^{(1)}|=C_\alpha ^{(1)}\frac{1}{\sqrt{2}}{}_{1/2}{}^{}\stackrel{~}{b,k/2|}{}_{3/2}{}^{}a,k/2|\{\left(𝒞\gamma ^\alpha \mathrm{\Pi }_{}\right)_{ab}\mathrm{cos}\beta _0\stackrel{~}{\gamma }_0+\left(𝒞\gamma ^\alpha \mathrm{\Pi }_+\right)_{ab}\mathrm{sin}\beta _0\stackrel{~}{\gamma }_0\},$$
with
$$\mathrm{\Pi }_\pm =\frac{1}{2}\left(\text{1l}_8\pm \gamma \right),$$
(3.5)
the chirality projector in 6 dimensions, where the gamma-matrix $`\gamma `$ is defined in Appendix A. This state is normalised as $`C^{(1)}|C^{(1)}=C_\alpha ^{(1)}C^{(1)\alpha }`$. The corresponding projector is given by
$$P_{(C)}{}_{}{}^{\alpha }|=\frac{1}{\sqrt{2}}{}_{1/2}{}^{}\stackrel{~}{b,k_{}/2|}{}_{3/2}{}^{}a,k_{}/2|\{\left(𝒞\gamma ^\alpha \mathrm{\Pi }_{}\right)_{ab}\mathrm{cos}\beta _0\stackrel{~}{\gamma }_0+\left(𝒞\gamma ^\alpha \mathrm{\Pi }_+\right)_{ab}\mathrm{sin}\beta _0\stackrel{~}{\gamma }_0\}.$$
such that $`P_{(C)}{}_{\alpha }{}^{}|C^{(1)}=C_\alpha ^{(1)}`$. The form of this projector is similar to the case of the type IIA BPS D-particle, with the difference that the twisted R-R groundstate carries spinor indices of $`SO(1,5)`$.
The NS-NS (R-R) fields have non-zero overlap with the NS-NS (R-R) boundary state only. We find the following results for the asymptotic fields in momentum space:
$`\delta \varphi (k)`$ $`=`$ $`P_{(\varphi )}|𝒟_{a=1/2}|D0_{\mathrm{NS},\mathrm{U}}=T_0{\displaystyle \frac{3}{\sqrt{8}}}{\displaystyle \frac{V_1}{k^2}},`$ (3.6)
$`\delta h_{\mu \nu }(k)`$ $`=`$ $`P_{(h)\mu \nu }|𝒟_{a=1/2}|D0_{\mathrm{NS},\mathrm{U}}={\displaystyle \frac{T_0}{2}}{\displaystyle \frac{V_1}{k^2}}\mathrm{diag}({\displaystyle \frac{7}{4}},{\displaystyle \frac{1}{4}},\mathrm{},{\displaystyle \frac{1}{4}}).`$ (3.7)
For the twisted R-R field we only have a contribution along the worldline direction of the D-particle, namely:
$$\delta C_0^{(1)}(k_{})=P_{(C)\mathrm{\hspace{0.17em}0}}|𝒟_{a=0}|D0_{\mathrm{R},\mathrm{T}}=\frac{Q_0}{\sqrt{2}}\frac{V_1}{k_{}^2}.$$
(3.8)
Note that $`\delta C_0^{(1)}`$ depends only on the momenta transverse to the D-particle, longitudinal to the fixed plane. We now make use of the following Fourier transformation in order to translate the results into position space. For a generic momentum $`K`$, with $`d1`$ non-zero components, we have:
$$𝑑td^{(d1)}x\frac{e^{\mathrm{i}Kx}}{(d3)\mathrm{\Omega }_{d2}|x|^{(d3)}}=\frac{V_1}{K^2},$$
(3.9)
where $`x`$ are $`d1`$ spatial coordinates transverse to the D-particle and $`\mathrm{\Omega }_{d2}`$ is the area of a unit sphere surrounding the D-particle. For the NS-NS sector we have $`d=9`$, hence we find the following asymptotic behaviour for the metric and dilaton:
$`\delta \varphi (x)`$ $`=`$ $`{\displaystyle \frac{3T_0}{14\sqrt{2}\mathrm{\Omega }_8}}{\displaystyle \frac{1}{|x|^7}},`$
$`\delta h_{\mu \nu }(x)`$ $`=`$ $`{\displaystyle \frac{T_0}{14\mathrm{\Omega }_8}}{\displaystyle \frac{1}{|x|^7}}\mathrm{diag}({\displaystyle \frac{7}{4}},{\displaystyle \frac{1}{4}},\mathrm{},{\displaystyle \frac{1}{4}}).`$ (3.10)
In the R-R sector, since the momentum dependence is restricted to the fixed plane spatial directions, we have $`d=6`$, and the asymptotic form of the R-R 1-form is found to be:
$$\delta C^{(1)}(y)=\frac{Q_0}{3\sqrt{2}\mathrm{\Omega }_4}\frac{1}{|y|^3},$$
(3.11)
where now $`y`$ indicates the spatial directions in the fixed plane only, and $`\mathrm{\Omega }_4`$ is the volume of a unit 4-sphere surrounding the particle inside the fixed plane. Note that the power of $`|y|`$ is exactly the power expected for a particle in 6 dimensions. We see, however, that the untwisted fields do see the entire spacetime. Thus, although the D-particle is stuck on the fixed plane, the associated metric and dilaton background solution in the uncompactified case will extend also in the orbifolded directions. This is a consequence of the fact that the D-particle can emit massless untwisted fields off the fixed plane which are symmetrised in the momenta along the orbifolded directions.
### 3.2 The Compact Case
In order to consider a spacetime description for a stable non-BPS D-brane, in a regime where it remains valid classically, one should make sense of a superposition of them, which turns out to be possible only in the compact orbifold at the critical radius. In this section we extend the analysis of the metric and dilaton of the D-particle to the case in which the space transverse to the orbifold fixed plane is a 4-torus, i.e. type IIB on $`T^4/`$ $`(1)^{F_L}_4`$.
We derive first the asymptotic form of the metric and dilaton in the compactified case using an approximation in the background fields, that we will compare with the boundary state calculation. Consider one of the compact orbifolded directions: $`x_9x_9+2\pi R_9`$. A D-particle sitting at the origin of this $`S^1`$ can be seen from the point of view of the covering space as an infinite array of equally spaced D-particles. Thus we can write:
$$\frac{1}{|x|^7}\underset{n\mathrm{Z}\mathrm{Z}}{}\frac{1}{\left(r^2+(x_92\pi nR_9)^2\right)^{7/2}},$$
(3.12)
with $`|x|^2=r^2+x_9^2`$, and $`r^2=x_1^2+\mathrm{}+x_8^2`$. We can approximate the sum by an integral assuming that the distance in the non-compact directions is much larger that the size of the compact one, i.e. $`rR_9`$. Changing variables we can write
$$\underset{n\mathrm{Z}\mathrm{Z}}{}\frac{1}{\left(r^2+(x_92\pi nR_9)^2\right)^{7/2}}\frac{1}{2\pi R_9r^6}\underset{\mathrm{}}{\overset{+\mathrm{}}{}}𝑑u\frac{1}{(1+u^2)^{7/2}}=\frac{I_5}{2\pi R_9}\frac{1}{r^6},$$
(3.13)
where we use the notation:
$$I_n=_0^\pi 𝑑\theta \mathrm{sin}^n\theta .$$
(3.14)
Moreover, these integrals satisfy the properties $`nI_n=(n1)I_{n2}`$, for $`n2`$, and
$$\mathrm{\Omega }_d=2I_{d1}I_{d2}\mathrm{}I_1I_0.$$
(3.15)
If we do this approximation for each of the compactified directions of the orbifold we obtain
$$\frac{1}{|x|^7}I_2I_3I_4I_5\underset{i=6}{\overset{9}{}}(2\pi R_i)^1\frac{1}{|y|^3},$$
(3.16)
with $`y`$ now indicating the spatial directions in the fixed plane. Within this approximation we obtain the following asymptotic behaviour for the dilaton and metric for the compactified case:
$`\delta \varphi (y)`$ $``$ $`{\displaystyle \frac{T_0}{2\sqrt{2}\mathrm{\Omega }_4}}{\displaystyle \underset{i=6}{\overset{9}{}}}(2\pi R_i)^1{\displaystyle \frac{1}{|y|^3}},`$
$`\delta h_{\mu \nu }(y)`$ $``$ $`{\displaystyle \frac{T_0}{6\mathrm{\Omega }_4}}{\displaystyle \underset{i=6}{\overset{9}{}}}(2\pi R_i)^1{\displaystyle \frac{1}{|y|^3}}\mathrm{diag}({\displaystyle \frac{7}{4}},{\displaystyle \frac{1}{4}},\mathrm{},{\displaystyle \frac{1}{4}}).`$ (3.17)
We rewrite the fluctuations using the 6-dimensional gravitational constant
$`g_{\mu \nu }\eta _{\mu \nu }+2\kappa _{10}\delta h_{\mu \nu }`$ $``$ $`\eta _{\mu \nu }+2\kappa _6\delta \overline{h}_{\mu \nu },`$
$`\varphi \kappa _{10}\sqrt{2}\delta \varphi `$ $``$ $`\kappa _6\sqrt{2}\delta \overline{\varphi },`$ (3.18)
$`\delta C^{(1)}`$ $``$ $`\delta \overline{C}^{(1)}.`$
Using the relation<sup>13</sup><sup>13</sup>13We use that for any dimension $`D`$, $`2\kappa _D^2=16\pi G_D`$, and for any lower dimension $`Dd`$, $`G_{Dd}=G_D/V_d`$, with $`V_d`$ the volume of the compact space. between $`\kappa _{10}`$ and $`\kappa _6`$: $`\kappa _{10}=\kappa _6\underset{i=6}{\overset{9}{}}(2\pi R_i)^{1/2}`$, we obtain
$`\delta \overline{\varphi }(y)`$ $``$ $`{\displaystyle \frac{T_0}{2\sqrt{2}\mathrm{\Omega }_4}}{\displaystyle \underset{i=6}{\overset{9}{}}}(2\pi R_i)^{1/2}{\displaystyle \frac{1}{|y|^3}},`$
$`\delta \overline{h}_{\mu \nu }(y)`$ $``$ $`{\displaystyle \frac{T_0}{6\mathrm{\Omega }_4}}{\displaystyle \underset{i=6}{\overset{9}{}}}(2\pi R_i)^{1/2}{\displaystyle \frac{1}{|y|^3}}\mathrm{diag}({\displaystyle \frac{7}{4}},{\displaystyle \frac{1}{4}},\mathrm{},{\displaystyle \frac{1}{4}}).`$ (3.19)
The asymptotic behaviour of the twisted R-R 1-form stays the same as for the non-compact case (3.11).
We compare now this approximation with the asymptotic behaviour derived from the boundary state. The closed string fields have an infinite number of Kaluza-Klein modes, since the momentum is quantised in the orbifolded directions. The massless closed string states excited by the D-particle correspond to the massless Kaluza-Klein modes. Using the same notation as before and using bars for the objects in the compact case, we can write the projection operators for the massless states as follows:
$`\overline{P}_{(\varphi )}|`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}{\displaystyle \underset{i=6}{\overset{9}{}}}\mathrm{\Phi }_i^{1/2}k_{},n_i=0|\psi _{1/2}^\nu \stackrel{~}{\psi }_{1/2}^\mu (\eta _{\mu \nu }(\delta _\mu {}_{}{}^{\alpha }\delta _{\nu }^{}{}_{}{}^{\beta }+\delta _\mu {}_{}{}^{\beta }\delta _{\nu }^{}{}_{}{}^{\alpha })k_\alpha \mathrm{}_\beta ),`$
$`\overline{P}_{(h)}^{\mu \nu }|`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=6}{\overset{9}{}}}\mathrm{\Phi }_i^{1/2}k_{},n^i=0|\left(\psi _{1/2}^\nu \stackrel{~}{\psi }_{1/2}^\mu +\psi _{1/2}^\mu \stackrel{~}{\psi }_{1/2}^\nu \right)`$
$`\overline{P}_{(\varphi )}|{\displaystyle \frac{1}{2\sqrt{2}}}(\eta ^{\mu \nu }(\delta ^\mu {}_{\alpha }{}^{}\delta _{}^{\nu }{}_{\beta }{}^{}+\delta ^\mu {}_{\beta }{}^{}\delta _{}^{\nu }{}_{\alpha }{}^{})k^\alpha \mathrm{}^\beta ),`$
where now $`k^\alpha \eta _{\alpha \beta }\mathrm{}^\beta =1`$ and $`\mathrm{}^\alpha \mathrm{}^\beta \eta _{\alpha \beta }=0`$. Instead of splitting the projectors in terms of 6-dimensional fields, we keep the ten-dimensional notation in order to compare it with the non-compact results, and with the approximation made previously. As before, the NS-NS fields will have non-zero overlap with the NS-NS boundary state. The type of calculation we perform in this case is similar to the non-compact case, with the difference that there appear extra normalisation coefficients, and the momentum dependence is only in the spatial directions of the fixed plane. Note that the factors of $`\underset{i}{}\mathrm{\Phi }_i^{1/2}`$ in the normalisation of the boundary state and the massless states cancel with a factor $`\underset{i}{}\mathrm{\Phi }_i`$ coming from the normalisation (2.35) in the amplitude. We find the following results:
$`\delta \overline{\varphi }(k)`$ $`=`$ $`\overline{P}_{(\varphi )}|𝒟_{a=1/2}|D0_{\mathrm{NS},\mathrm{U}}={\displaystyle \frac{3T_0}{2\sqrt{2}}}{\displaystyle \underset{i}{}}(2\pi R_i)^{1/2}{\displaystyle \frac{V_1}{k_{}^2}},`$
$`\delta \overline{h}_{\mu \nu }(k)`$ $`=`$ $`\overline{P}_{(h)\mu \nu }|𝒟_{a=1/2}|D0_{\mathrm{NS},\mathrm{U}}`$
$`=`$ $`{\displaystyle \frac{T_0}{2}}{\displaystyle \underset{i}{}}(2\pi R_i)^{1/2}{\displaystyle \frac{V_1}{k_{}^2}}\mathrm{diag}({\displaystyle \frac{7}{4}},{\displaystyle \frac{1}{4}},\mathrm{},{\displaystyle \frac{1}{4}}).`$
We can make use of the Fourier transformation (3.9) in the NS-NS sector to derive the spacetime behaviour:
$`\delta \overline{\varphi }(y)`$ $`=`$ $`{\displaystyle \frac{T_0}{2\sqrt{2}\mathrm{\Omega }_4}}{\displaystyle \underset{i}{}}(2\pi R_i)^{1/2}{\displaystyle \frac{1}{|y|^3}},`$
$`\delta \overline{h}_{\mu \nu }(y)`$ $`=`$ $`{\displaystyle \frac{T_0}{6\mathrm{\Omega }_4}}{\displaystyle \underset{i}{}}(2\pi R_i)^{1/2}{\displaystyle \frac{1}{|y|^3}}\mathrm{diag}({\displaystyle \frac{7}{4}},{\displaystyle \frac{1}{4}},\mathrm{},{\displaystyle \frac{1}{4}}).`$ (3.22)
The twisted R-R sector remains the same as for the uncompactified case (3.11). Note that the asymptotic behaviour derived from the compactified boundary state (3.22) exactly reproduces the behaviour of the 6-dimensional fluctuations obtained before (3.19) using the approximation given in (3.16). The asymptotics of the dilaton and the graviton are now certainly comparable with that of the twisted R-R 1-form. In particular, at the critical radius there is an accidental Bose-Fermi degeneracy, which as we will see in the next Section, implies a no-force in a brane probe. At the critical radius we can then consider a large number of branes $`N`$, hence the fluctuations take the following form:
$`\delta \overline{\varphi }(y)`$ $`=`$ $`{\displaystyle \frac{T_0N}{2\sqrt{2}\mathrm{\Omega }_4(2\pi ^2\alpha ^{})}}{\displaystyle \frac{1}{|y|^3}},`$
$`\delta \overline{h}_{\mu \nu }(y)`$ $`=`$ $`{\displaystyle \frac{T_0N}{6\mathrm{\Omega }_4(2\pi ^2\alpha ^{})}}{\displaystyle \frac{1}{|y|^3}}\mathrm{diag}({\displaystyle \frac{7}{4}},{\displaystyle \frac{1}{4}},\mathrm{},{\displaystyle \frac{1}{4}}),`$ (3.23)
$`\delta C^{(1)}(y)`$ $`=`$ $`{\displaystyle \frac{Q_0N}{3\sqrt{2}\mathrm{\Omega }_4}}{\displaystyle \frac{1}{|y|^3}}.`$
## 4 No-force Condition at the Critical Radius
A composite of branes preserving the BPS property satisfies a no-force condition between the constituents. This was verified at the level of the effective action and background geometries for BPS branes in . The no-force property is a consequence of the vanishing of the one-loop amplitude between the D-branes. On the other hand, the Bose-Fermi degeneracy for non-BPS D-branes at the critical radius described in represents a remarkable example in which the no-force property takes place without being BPS. In this section we recover the no-force condition using the non-BPS D-particle as a probe of the geometry of another non-BPS D-particle. Firstly, we consider the effective action for the non-BPS D-particle.
### 4.1 The Action of the D-particle
Following the prescription given in one can construct the action for the D-particle in the orbifold of type IIB we are considering. One starts with the action of a non-BPS D-particle in type IIB, and set to zero all fields which are odd under $`(1)^{F_L}_4`$. The worldvolume theory of this D-particle is given by scalars $`b_{1/2}^\mu |0𝒾_{\mathrm{NS}}\text{1l}`$, $`\mu =1,\mathrm{},9`$, a tachyon $`|0𝒾_{\mathrm{NS}}\sigma _1`$, and 16 fermions in the Ramond groundstate of each of the Chan-Paton sectors 1l and $`\sigma _1`$. The orbifold $`(1)^{F_L}_4`$ eliminates the scalars in the directions $`i=6,7,8,9`$, the tachyon and the fermions with Chan-Paton factor $`\sigma _1`$ . Accordingly, the non-BPS D-particle is stable and contains 5 physical bosonic degrees of freedom, describing the motion of the D-particle on the orbifold fixed plane, plus 16 fermionic d.o.f. If we label the coordinates as before $`X^\mu =(X^\alpha ,X^i)`$, where $`\alpha =0,1,\mathrm{},5`$ labels the orbifold fixed plane directions, we can write the following action for the D-particle:
$$S_{kin}=\frac{T_0}{\kappa _{10}}𝑑\tau \mathrm{e}^\varphi \sqrt{|_\tau X^\alpha _\tau X^\beta G_{\alpha \beta }|},$$
(4.1)
where $`G_{\mu \nu }`$ is the background string metric. Furthermore, since the D-particle carries charge under the twisted R-R 1-form $`𝒞^{(1)}`$, we can include a WZ coupling<sup>14</sup><sup>14</sup>14Wess-Zumino couplings for stable non-BPS D-branes have been recently considered in ., which involves the 1-form potential on the fixed plane:
$$S_{WZ}=\frac{Q_0}{2\kappa _6}_\tau X^\alpha 𝒞_\alpha ^{(1)}.$$
(4.2)
The couplings of the stable non-BPS D-particle to the background fields given by equations (4.1) and (4.2) can be also derived by projecting the boundary state of the D-particle onto the projectors defined previously, analogously to the BPS case considered in . The coupling to the dilaton is given by
$$𝒥_{(\phi )}=P_{(\varphi )}|D0_{\mathrm{NS},\mathrm{U}}=\frac{3T_0}{2\sqrt{2}}V_1,$$
(4.3)
whereas the coupling to the metric is given by
$$𝒥_{(h)}=\eta _{\mu \nu }P_{(h)}^{\mu \nu }|D0_{\mathrm{NS},\mathrm{U}}=T_0V_1h_{00},$$
(4.4)
where $`h_{\mu \nu }`$ is the symmetric and traceless helicity tensor of the graviton. These two results reproduce the couplings of (4.1) in Einstein frame ($`g_{\mu \nu }=\mathrm{e}^{\varphi /2}G_{\mu \nu }`$) after rescaling the dilaton as $`\varphi =\kappa _{10}\sqrt{2}\phi `$. Finally, the coupling to the twisted R-R 1-form is given by:
$$𝒥_{(C)}=C_\alpha ^{(1)}P_{(C)}^\alpha |D0_{\mathrm{R},\mathrm{T}}=\frac{Q_0}{\sqrt{2}}V_1C_0^{(1)},$$
(4.5)
which reproduces (4.2), after the rescaling $`𝒞^{(1)}=\kappa _6\sqrt{2}C^{(1)}`$.
In the compact case, the coupling of the D-particle to the massless closed string fields can be obtained as above, this time using the compactified boundary state constructed in Section 2.2. The net result is a change of the factor in front of the kinetic term, which amounts to the relation of the 6-dimensional gravitational constant with the 10-dimensional one. The boundary state reproduces the couplings of the following action:
$$S_{kin}=\frac{T_0}{\kappa _6}\underset{i=6}{\overset{9}{}}(2\pi R_i)^{1/2}𝑑\tau \mathrm{e}^\varphi \sqrt{|_\tau X^\alpha _\tau X^\beta G_{\alpha \beta }|}.$$
(4.6)
Thus although the D-particle does not couple to the components $`G_{ij}`$ (4.1), it does feel the transverse geometry through the renormalised tension. Moreover, the open strings on the D-particle carry winding modes, and at the critical radius, the winding modes of the tachyon become massless. These extra modes appear as new degrees of freedom, $`\chi ^i`$, in the effective action :
$$S_{kin}=\frac{T_0}{2\pi ^2\alpha ^{}\kappa _6}𝑑\tau \mathrm{e}^\varphi \sqrt{|_\tau X^\alpha _\tau X^\beta G_{\alpha \beta }|}\left(1G^{\tau \tau }_\tau \chi ^i_\tau \chi ^i\right).$$
(4.7)
However, these extra modes will not play any role in our discussion below.
### 4.2 D-particle at the Critical Radius
We consider a non-BPS D-particle probe moving in the background of another non-BPS D-particle in the compactified orbifold. We split the embedding coordinates as $`X^\alpha =(X^0,X^m)`$, $`m=1,\mathrm{},5`$, and choose the static gauge:
$$X^0=\tau .$$
(4.8)
We identify the embedding scalars with the transverse coordinates of the background geometry and expand the action in derivatives<sup>15</sup><sup>15</sup>15We use the Einstein frame since we want to use in the action the results given by the boundary state computation.:
$$\sqrt{|_\tau y^\alpha _\tau y^\beta g_{\alpha \beta }|}=\sqrt{|g_{00}|}\left(1\frac{1}{2g_{00}}_0y^m_0y^ng_{mn}+\mathrm{}\right).$$
(4.9)
Moreover, for a D-particle background, only the zeroth component of the twisted R-R 1-form is turned on:
$$_\tau X^\alpha 𝒞_\alpha ^{(1)}=𝒞_0^{(1)}.$$
(4.10)
The action can then be approximated by
$$S𝑑\tau 𝒱(y)+\frac{T_0}{\kappa _6}\underset{i=6}{\overset{9}{}}(2\pi R_i)^{1/2}𝑑\tau \frac{\mathrm{e}^{\frac{3}{4}\varphi }}{2\sqrt{|g_{00}|}}_0y^m_0y^ng_{mn},$$
(4.11)
where we have defined the effective static potential $`𝒱`$:
$$𝒱(y)=\frac{T_0}{\kappa _6}\underset{i=6}{\overset{9}{}}(2\pi R_i)^{1/2}\mathrm{e}^{\frac{3}{4}\varphi }\sqrt{|g_{00}|}+\frac{Q_0}{2\kappa _6}𝒞_0^{(1)}.$$
(4.12)
If this potential is not constant or zero for a given background $`g_{00}(y)`$, $`\varphi (y)`$, $`𝒞_0^{(1)}(y)`$, there will be a force term in the field equations for $`y^m`$. Plugging into this potential the background geometry generated by another D-particle, one can check whether the Bose-Fermi degeneracy takes place. Since from the boundary state we obtain the asymptotic form of the D-particle solution, we rewrite the potential $`𝒱`$ in terms of the fluctuations (3.2)<sup>16</sup><sup>16</sup>16We neglect the constant parts of the potential, since these would not generate any force.:
$$𝒱(y)T_0\underset{i=6}{\overset{9}{}}(2\pi R_i)^{1/2}\left(\frac{3}{2\sqrt{2}}\delta \overline{\varphi }+\delta \overline{h}_{00}\right)+\frac{Q_0}{\sqrt{2}}\delta C_0^{(1)}.$$
(4.13)
We can insert now the asymptotic form of the solution for the compactified case given by (3.22) and (3.11). The potential then takes the form:
$`𝒱(y)`$ $`=`$ $`{\displaystyle \frac{2}{3\mathrm{\Omega }_4|y|^3}}\left(T_0^2{\displaystyle \underset{i=6}{\overset{9}{}}}(2\pi R_i)^1{\displaystyle \frac{Q_0^2}{4}}\right)`$ (4.14)
$`=`$ $`{\displaystyle \frac{4\pi \alpha ^{}}{|y|^3}}\left(\left({\displaystyle \frac{\alpha ^{}}{2}}\right)^2{\displaystyle \underset{i=6}{\overset{9}{}}}R_i^11\right),`$
which coincides with the static potential derived from the cylinder amplitude in (2.61) and vanishes for the critical radius ($`R_i=\sqrt{\alpha ^{}/2}`$), hence we recover the no-force at the critical radius using the background geometry generated by the non-BPS D-particles.
According to the interaction potential $`𝒰_{open}`$ derived in Section 2.3, the velocity corrections start at order $`v^4`$ and there is no $`v^2`$ corrections to the potential in for any radius, as happens for BPS branes. From the probe point of view, this translates into the fact that the metric on the moduli space, multiplying the velocity dependent piece of the action (4.11), is flat . In the asymptotic limit, for the non-compact orbifold this metric takes the form:
$$\frac{e^{\frac{3}{4}\varphi }}{\sqrt{|g_{00}|}}g_{mn}\delta _{mn}\left(1\frac{3}{2\sqrt{2}}\kappa _{10}\delta \varphi (x)+\kappa _{10}\delta h_{00}(x)\right)+2\kappa _{10}\delta h_{mn}(x).$$
(4.15)
Using the asymptotic fluctuations given in (3.10), it is easy to see that the metric is flat ($`=\delta _{mn}`$). In the case of the compact orbifold, this metric factor takes the form:
$$\frac{e^{\frac{3}{4}\varphi }}{\sqrt{|g_{00}|}}g_{mn}\delta _{mn}\left(1\frac{3}{2\sqrt{2}}\kappa _6\delta \overline{\varphi }(y)+\kappa _6\delta \overline{h}_{00}(y)\right)+2\kappa _6\delta \overline{h}_{mn}(y).$$
(4.16)
Substituting the asymptotic fluctuations (3.22), we obtain the same flat metric for any value of the radii. Thus we recover the behaviour given by the open strings in (2.57) and (2.64).
### 4.3 The Classical Geometry of the non-BPS D-particle
In this Section we assume that the no-force condition will persist at the full non-linear level of the field equations. We can then restrict the possible classical geometries by imposing the no-force at the level of equation (4.12). No-force can occur for $`𝒱`$ either constant or zero . However, in our case, the no-force condition at the critical radius must enforce the relation $`T_0=(\pi ^2\alpha ^{})Q_0`$, as seen in (2.47). These are in fact the factors of the NS-NS and R-R contributions of the potential (4.12), respectively. Hence we conclude that the classical solution must be such that at the critical radius the following equality holds
$$e^{\frac{3}{4}\varphi }\sqrt{|g_{00}|}=𝒞_0^{(1)}.$$
(4.17)
We can then deduce part of the form for $`g_{00}`$, $`\varphi `$, and $`𝒞^{(1)}`$:
$`g_{00}(y)`$ $`=`$ $`\left(1+{\displaystyle \frac{\kappa _6T_0}{2a\mathrm{\Omega }_4}}(2\pi ^2\alpha ^{})^1{\displaystyle \frac{1}{|y|^3}}+\mathrm{}\right)^{\frac{7}{6}a},`$
$`\mathrm{e}^\varphi (y)`$ $`=`$ $`\left(1+{\displaystyle \frac{\kappa _6T_0}{2a\mathrm{\Omega }_4}}(2\pi \alpha ^{})^1{\displaystyle \frac{1}{|y|^3}}+\mathrm{}\right)^a,`$ (4.18)
$`𝒞_0^{(1)}(y)`$ $`=`$ $`\left(1+{\displaystyle \frac{\kappa _6Q_0}{4a\mathrm{\Omega }_4}}{\displaystyle \frac{1}{|y|^3}}+\mathrm{}\right)^{\frac{4}{3}a}1.`$
These functions are given in terms of a parameter $`a`$ yet to be determined, and are such that asymptotically they become the fluctuations in (3.19) and (3.11), and at the critical radius equation (4.17) holds. The dots indicate other possible contributions with dependence $`|y|^n`$, $`n<3`$, which are subleading in the asymptotic limit ($`|y|\mathrm{}`$), but which are relevant when we come closer to the brane. We expect these extra terms to appear since harmonicity, which is a direct consequence of supersymmetry, may not be present in the complete solution.
We can make one further assumption if we consider that the metric factor in the velocity dependent piece of the action (4.11) remains flat for the complete geometry. We impose then the following relation:
$$\frac{\mathrm{e}^{\frac{3}{4}\varphi }}{\sqrt{|g_{00}|}}g_{mn}=\delta _{mn}.$$
(4.19)
Note that the functions describing the dilaton and the metric must be the same at the critical radius. Moreover, it is only at this radius where we actually expect to find a consistent classical geometry. Therefore, we make use of the expressions for $`g_{00}(y)`$ and $`\varphi (y)`$ given above to obtain:
$$g_{mn}(y)=\left(1+\frac{\kappa _6T_0}{2a\mathrm{\Omega }_4}(2\pi ^2\alpha ^{})^1\frac{1}{|y|^3}+\mathrm{}\right)^{\frac{1}{6}a}\delta _{mn},$$
(4.20)
where the dots represent the same contributions as in (4.18) at the critical radius. Finally, it seems that $`g_{ij}`$ is out of the reach of the present analysis. A possible way to include it in the analysis is by a coupling to the extra massless states $`\chi ^i`$ in (4.7). On the other hand, although the asymptotic behaviour of the metric in the non-compact case (3.10) presents an $`SO(9)`$ symmetry, this is certainly broken to $`SO(5)\times SO(4)`$ by the orbifold when we get closer to the position of the brane. This is in fact already suggested by the asymptotics in the compact case (3.19). Accordingly, the form of $`g_{ij}(y)`$ is expected to be different from $`g_{mn}(y)`$.
## 5 Comments
In this paper we have investigated the description of a stable non-BPS D-particle in terms of a classical solution. We have used the technique of the boundary state to obtain the asymptotic form of the classical solution for the non-compact and compact versions of the orbifold. We find a metric and a dilaton propagating in the bulk, and a twisted R-R 1-form propagating in the fixed plane. In the non-compact case the bulk fields have the dependence expected for a particle in ten dimensions, whereas in the compact case they have a dependence typical of a particle in six dimensions. The twisted R-R 1-form has in both cases the same asymptotic form with the usual dependence for a particle in six-dimensions. Using the non-BPS D-particle as a probe in the background of another non-BPS D-particle, we have recovered the no-force property at the critical radius using the asymptotic behaviour. Moreover, we have calculated the cylinder amplitude for non-BPS D-particles in relative motion. From it we have extracted the long and short distance interactions. For generic radii these contain $`v^2`$ corrections in the closed string description, but the open string description presents no $`v^2`$ terms for all radii, like for BPS D-branes. Moreover, the $`v^4`$ corrections do not match, unlike for BPS branes. On the other hand, at the critical radius they present a BPS-like behaviour, up to the $`v^4`$ corrections, which do not match in the open and closed descriptions.
We have assumed that the no-force property of a brane probe holds for the full background geometry. This is acceptable for distances much larger than the string scale. This assumption allows us to derive part of the classical solution, which reproduces the asymptotic behaviour and the no-force property. On the other hand, we expect that extra terms may appear in the solution at the classical level. This is due to the fact that the boundary state only gives information about the next-to-leading term in the asymptotic limit, hence subleading terms which become relevant at short distances escape from this analysis. Moreover, the fact that there is a coordinate system in which the metric can be written in terms of harmonic functions is very much related to residual supersymmetry, which does not occur in our case. This fact does not permit to find information about the geometry near the core of the non-BPS D-particle. Moreover, although the asymptotic behaviour of the metric (3.10) presents an $`SO(9)`$ symmetry, this is expected to be broken to $`SO(5)\times SO(4)`$ by the orbifold when we get closer to the position of the brane.
The form of the asymptotic behaviour we have found suggests that the classical solution for the non-BPS D-particle will be a solution to the field equations derived from an action involving 10-dimensional bulk fields and 6-dimensional matter fields constrained on an orbifold fixed plane:
$$𝒮_{total}=𝒮_{bulk}+𝒮_{plane}.$$
(5.1)
The bulk action involves the metric and the dilaton, and other fields of the massless untwisted sector. For the case of the D-particle, only the metric and dilaton are relevant, hence in Einstein frame we have:
$$𝒮_{bulk}=\frac{1}{2\kappa _{10}^2}d^{10}x\sqrt{|\mathrm{det}(g_{\mu \nu })|}\left(\frac{1}{2}(\varphi )^2\right).$$
(5.2)
Since the twisted sector is provided by a type IIB NS-5 brane hidden in the orbifold fixed plane, we can derive the fixed plane action from the effective action of the type IIB NS-5 brane in Einstein frame:
$$𝒮_{plane}=m\underset{x^i=0}{}d^6y\sqrt{|\mathrm{det}(\stackrel{~}{g}_{\alpha \beta })|}\frac{\mathrm{e}^{\frac{3}{2}\stackrel{~}{\varphi }}}{4}F_{\alpha \beta }F^{\alpha \beta },$$
(5.3)
where the tilde denotes the restriction of the bulk fields to the position of the fixed plane:
$$\stackrel{~}{g}_{\alpha \beta }(y^\alpha )=g_{\alpha \beta }(y^\alpha ,x^i=0),\stackrel{~}{\varphi }(y^\alpha )=\varphi (y^\alpha ,x^i=0),$$
(5.4)
and $`\alpha ,\beta =0,1,\mathrm{},5`$ are the indices along the fixed plane. This particular coupling of the twisted fields with the bulk metric is obtained by expanding the NS-5 brane kinetic term in powers of the worldvolume fields. The twisted R-R 1-form has been identified with the $`U(1)`$ gauge field on the NS-5 brane worldvolume $`F`$. Here $`m`$ is a factor related to the tension of the NS-5 brane. Moreover, the embedding scalars of the NS-5 brane has not been included, since they correspond to the twisted NS-NS sector, to which the non-BPS D-particle does not couple. It is straightforward to see that the asymptotic fields (3.10) and (3.11) are a solution to the weak field limit of the equations of motion of the action $`𝒮_{total}`$. Note that the non-BPS D-particle is not stable below the critical radius, therefore we do not expect it would appear as a solution of this action reduced to six dimensions.
Finally, the behaviour of the stable non-BPS D-particle at the critical radius suggests that it probably saturates a BPS type of bound in the effective theory, without being supersymmetric. This latter property comes from the fact that the particle couples to a 1-form in the wrong supersymmetric multiplet. The analysis carried out here can be extended to other stable non-BPS branes. A study of the classical geometry of other stable non-BPS D-branes will be presented in a future publication .
## Acknowledgements
We would like to thank E. Bergshoeff and M. de Roo for useful discussions. E.E. is greatful to M. Gaberdiel, T. Dasgupta, N. Lambert, A. Liccardo, M.J. Perry, B. Stefanski, P. Townsend and A. Uranga for useful discussions. E.E. has also enjoyed discussions with C. van der Bruck, P. Vanhove and N. Wyllard. The work of E.E. is supported by the European Community program ”Human Potential” under the contract HPMF-CT-1999-00018. This work is also partially supported by the PPARC grant PPA/G/S/1998/00613.
## Appendix A Appendix
In this Appendix we present some details about the zero-mode part of the twisted R-R boundary state and its GSO projection. In order to describe the twisted R-R groundstate we make use of the $`8\times 8`$ gamma matrices of $`SO(1,5)`$:
$$\{\gamma ^\alpha ,\gamma ^\beta \}=2\text{1l}_8\eta ^{\alpha \beta }.$$
(A.1)
We define
$$\gamma =\gamma ^0\gamma ^1\gamma ^2\gamma ^3\gamma ^4\gamma ^5,$$
(A.2)
such that $`\{\gamma ^\alpha ,\gamma \}=0`$ and $`(\gamma )^2=\text{1l}_8`$. Furthermore, there is a conjugation matrix
$$𝒞=\gamma ^3\gamma ^5\gamma ^0,$$
(A.3)
such that
$$\{\gamma ,𝒞\}=0,(\gamma ^\alpha )^T=𝒞\gamma ^\alpha 𝒞^1,𝒞^1=𝒞,𝒞^T=𝒞.$$
(A.4)
The twisted R-R groundstate is characterised by left and right spinor indices of $`SO(1,5)`$, and can be constructed from the NS-vacuum by means of spin and twist fields as follows:
$$\underset{z,\overline{z}0}{lim}S^a(z)\mathrm{\Sigma }(z)\stackrel{~}{S}^b(\overline{z})\stackrel{~}{\mathrm{\Sigma }}(\overline{z})|0=|a_\mathrm{T}\stackrel{~}{|b}_\mathrm{T}.$$
(A.5)
The zero-mode part of the twisted R-R boundary state satisfies the following overlap equations in the twisted sector:
$$\left(\psi _0^\alpha \mathrm{i}\eta S^\alpha {}_{\beta }{}^{}\stackrel{~}{\psi }_{0}^{\beta }\right)|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}^{(0)}=0,$$
(A.6)
where $`\alpha =0,1,\mathrm{}5`$ and $`S^\alpha _\beta `$ is the restriction of the matrix $`S^\mu _\nu `$ of the D-particle to the 6-dimensional orbifold fixed plane: $`S^\alpha {}_{\beta }{}^{}=\mathrm{diag}(+1,1,1,1,1,1)`$. We define
$$|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}^{(0)}=M_{ab}|a_\mathrm{T}\stackrel{~}{|b}_\mathrm{T}.$$
(A.7)
If we define the action of the fermionic zero-modes in the twisted Ramond sector as
$`\psi _0^\alpha |a_\mathrm{T}\stackrel{~}{|b}_\mathrm{T}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\gamma ^\alpha )^a{}_{c}{}^{}(\text{1l}_8)_{}^{b}{}_{d}{}^{}|c_\mathrm{T}\stackrel{~}{|d}_\mathrm{T},`$
$`\stackrel{~}{\psi }_0^\beta |a_\mathrm{T}\stackrel{~}{|b}_\mathrm{T}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\gamma )^a{}_{c}{}^{}(\gamma ^\beta )_{}^{b}{}_{d}{}^{}|c_\mathrm{T}\stackrel{~}{|d}_\mathrm{T},`$ (A.8)
the matrix $`M`$ must satisfy:
$$(\gamma ^\alpha )^TM\mathrm{i}\eta S^\alpha {}_{\beta }{}^{}\left(\gamma M\gamma ^\beta \right)=0.$$
(A.9)
The solution to this equation is given by:
$$M=𝒞\gamma ^0\frac{1+i\eta \gamma }{1+i\eta }.$$
(A.10)
The zero-mode R-R boundary state for a non-BPS D-particle moving in a direction $`m`$ is given by
$$|D0_\psi ,\eta ,v_{\mathrm{R},\mathrm{T}}^{(0)}=M_{ab}(v)|a_\mathrm{T}\stackrel{~}{|b}_\mathrm{T},$$
(A.11)
where the matrix $`M_{ab}(v)`$ is given by
$$M(v)=\frac{1}{\sqrt{1v^2}}𝒞\left(\gamma ^0+v\gamma ^m\right)\frac{1+i\eta \gamma }{1+i\eta }$$
(A.12)
On the other hand, defining
$${}_{\mathrm{R},\mathrm{T}}{}^{(0)}D0_\psi ,\eta |={}_{\mathrm{T}}{}^{}a|{}_{\mathrm{T}}{}^{}\stackrel{~}{b|}N_{ab}.$$
(A.13)
and using
$`{}_{\mathrm{T}}{}^{}a|{}_{\mathrm{T}}{}^{}\stackrel{~}{b|}\psi _0^\alpha `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}{}_{\mathrm{T}}{}^{}c|{}_{\mathrm{T}}{}^{}\stackrel{~}{d|}(\gamma ^\alpha )^a{}_{c}{}^{}(\text{1l}_8)_{}^{b}_d`$
$`{}_{\mathrm{T}}{}^{}a|{}_{\mathrm{T}}{}^{}\stackrel{~}{b|}\stackrel{~}{\psi }_0^\beta `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}{}_{\mathrm{T}}{}^{}c|{}_{\mathrm{T}}{}^{}\stackrel{~}{d|}(\gamma )^a{}_{c}{}^{}(\gamma ^\beta )_{}^{b}_d`$ (A.14)
we can write the matrix $`N`$ as follows:
$$N=𝒞\gamma ^0\frac{1+i\eta \gamma }{1i\eta }.$$
(A.15)
The boosted version can be written similarly to (A.12):
$$N(v)=\frac{1}{\sqrt{1v^2}}𝒞\left(\gamma ^0+v\gamma ^m\right)\frac{1+i\eta \gamma }{1i\eta }.$$
(A.16)
In order to implement the GSO projection on the zero-mode of the boundary state, we define the zero-mode part of the GSO-projector on the twisted R-R sector as:
$$\mathrm{\Psi }=8\psi _0^5\mathrm{}\psi _0^0,\stackrel{~}{\mathrm{\Psi }}=8\stackrel{~}{\psi }_0^5\mathrm{}\stackrel{~}{\psi }_0^0.$$
(A.17)
Using this definition and (A.8), the action of $`\mathrm{\Psi }`$ and $`\stackrel{~}{\mathrm{\Psi }}`$ on the twisted R-R groundstate is found to be:
$$\mathrm{\Psi }|a_\mathrm{T}\stackrel{~}{|b}_\mathrm{T}=(\gamma )^a{}_{c}{}^{}|c_\mathrm{T}\stackrel{~}{|b}_\mathrm{T},\stackrel{~}{\mathrm{\Psi }}|a_\mathrm{T}\stackrel{~}{|b}_\mathrm{T}=(\gamma )^b{}_{d}{}^{}|a_\mathrm{T}\stackrel{~}{|d}_\mathrm{T}.$$
(A.18)
Finally, combining (A.10) and (A.18) we find:
$$\mathrm{\Psi }|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}^{(0)}=|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}^{(0)},\stackrel{~}{\mathrm{\Psi }}|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}^{(0)}=|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}^{(0)}.$$
(A.19)
## Appendix B Appendix
For the sake of completeness we include in this Appendix the explicit definitions of the different parts of the boundary state for the non-BPS D-particle. We include as well the conjugate states. Before GSO-projection the NS-NS boundary state is given by:
$$|D0,\eta _{\mathrm{NS},\mathrm{U}}=\frac{T_0}{2}|D0_X|D0_{\mathrm{gh}}|D0_\psi ,\eta _{\mathrm{NS}}|D0_{\mathrm{sgh}},\eta _{\mathrm{NS}}.$$
(B.1)
The bosonic part is:
$$|D0_X=\delta ^{(5)}(\widehat{q}^py^p)\delta ^{(4)}(\widehat{q}^i)\mathrm{exp}\left(\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}\alpha _nS\stackrel{~}{\alpha }_n\right)|k=0.$$
(B.2)
The ghost part:
$$|D0_{\mathrm{gh}}=\mathrm{exp}\left(\underset{n=1}{\overset{\mathrm{}}{}}(c_n\stackrel{~}{b}_nb_n\stackrel{~}{c}_n)\right)\left(\frac{c_0+\stackrel{~}{c}_0}{2}\right)|1\stackrel{~}{|1}.$$
(B.3)
The fermionic part:
$$|D0_\psi ,\eta _{\mathrm{NS}}=i\mathrm{exp}\left(i\eta \underset{r=1/2}{\overset{\mathrm{}}{}}\psi _rS\stackrel{~}{\psi }_r\right)|0,$$
(B.4)
The superghost part:
$$|D0_{\mathrm{sgh}},\eta _{\mathrm{NS}}=\mathrm{exp}\left(i\eta \underset{r=1/2}{\overset{\mathrm{}}{}}(\gamma _r\stackrel{~}{\beta }_r\beta _r\stackrel{~}{\gamma }_r)\right)|1\stackrel{~}{|1}.$$
(B.5)
The conjugate states are:
$`D0_X|`$ $`=`$ $`k=0|\delta ^{(5)}(\widehat{q}^py^p)\delta ^{(4)}(\widehat{q}^i)\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}\alpha _nS\stackrel{~}{\alpha }_n\right),`$
$`D0_{\mathrm{gh}}|`$ $`=`$ $`2|\stackrel{~}{2|}\left(b_0\stackrel{~}{b}_0\right)\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(\stackrel{~}{b}_nc_n\stackrel{~}{c}_nb_n)\right),`$
$`{}_{\mathrm{NS}}{}^{}D0_\psi ,\eta |`$ $`=`$ $`i0|\mathrm{exp}\left(i\eta {\displaystyle \underset{r=1/2}{\overset{\mathrm{}}{}}}\psi _rS\stackrel{~}{\psi }_r\right),`$
$`{}_{\mathrm{NS}}{}^{}D0_{\mathrm{sgh}},\eta |`$ $`=`$ $`1|\stackrel{~}{1|}\mathrm{exp}\left(i\eta {\displaystyle \underset{r=1/2}{\overset{\mathrm{}}{}}}(\beta _r\stackrel{~}{\gamma }_r\gamma _r\stackrel{~}{\beta }_r)\right).`$ (B.6)
For the twisted R-R sector we have (before GSO-projection):
$$|D0,\eta _{\mathrm{R},\mathrm{T}}=\frac{Q_0}{2}|D0_X_\mathrm{T}|D0_{\mathrm{gh}}|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}|D0_{\mathrm{sgh}},\eta _\mathrm{R},$$
(B.7)
The bosonic part is given by
$$|D0_X_\mathrm{T}=\delta ^{(5)}(\widehat{q}^py^p)\mathrm{exp}\left(\underset{t.m.}{}\frac{1}{n}\alpha _nS\stackrel{~}{\alpha }_n\right)|k=0$$
(B.8)
where $`t.m.`$ indicates that the sum is performed according to the twisted moddings of the (twisted) R-R sector given in (2.10). The ghost part is the same as for the NS-NS boundary state. The fermionic part reads:
$$|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}=\mathrm{exp}\left(i\eta \underset{t.m.}{}\psi _nS\stackrel{~}{\psi }_n\right)|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}^{(0)},$$
(B.9)
where $`|D0_\psi ,\eta _{\mathrm{R},\mathrm{T}}^{(0)}`$ and is given in Appendix A. Finally, the superghost part:
$$|D0_{\mathrm{sgh}},\eta _\mathrm{R}=\mathrm{exp}\left(i\eta \underset{n=1}{\overset{\mathrm{}}{}}(\gamma _n\stackrel{~}{\beta }_n\beta _n\stackrel{~}{\gamma }_n)\right)|D0_{\mathrm{sgh}},\eta _\mathrm{R}^{(0)},$$
(B.10)
with the zero-mode given by
$$|D0_{\mathrm{sgh}},\eta _\mathrm{R}^{(0)}=\mathrm{e}^{i\eta \gamma _0\stackrel{~}{\beta }_0}|1/2\stackrel{~}{|3/2}.$$
(B.11)
Finally, the conjugate states are:
$`{}_{\mathrm{T}}{}^{}D0_X|`$ $`=`$ $`k=0|\delta ^{(5)}(\widehat{q}^py^p)\mathrm{exp}\left({\displaystyle \underset{t.m.}{}}{\displaystyle \frac{1}{n}}\alpha _nS\stackrel{~}{\alpha }_n\right),`$
$`{}_{\mathrm{R},\mathrm{T}}{}^{}D0_\psi ,\eta |`$ $`=`$ $`{}_{\mathrm{R},\mathrm{T}}{}^{(0)}D0_\psi ,\eta |\mathrm{exp}(i\eta {\displaystyle \underset{t.m.}{}}\psi _nS\stackrel{~}{\psi }_n),`$
$`{}_{\mathrm{NS}}{}^{}D0_{\mathrm{sgh}},\eta |`$ $`=`$ $`{}_{\mathrm{R}}{}^{(0)}D0_{\mathrm{sgh}},\eta |\mathrm{exp}(i\eta {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(\beta _n\stackrel{~}{\gamma }_n\gamma _n\stackrel{~}{\beta }_n)),`$ (B.12)
where $`{}_{\mathrm{R}}{}^{(0)}D0_\psi ,\eta |`$ is given in Appendix A and
$${}_{\mathrm{R}}{}^{(0)}D0_{\mathrm{sgh}},\eta |=3/2|\stackrel{~}{1/2|}\mathrm{e}^{i\eta \beta _0\stackrel{~}{\gamma }_0}.$$
(B.13)
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# Persistent currents for Coulomb interacting electrons on 2d disordered lattices
## 1 Introduction
The interplay of disorder and interaction in mesoscopic samples has attracted considerable interest in recent years Moriond . A prominent experimental finding in this field is the insulator-metal transition in 2d systems of high mobility kravchenko , occurring when the carrier density is increased. This transition can neither be explained by disorder effects alone nor by interaction effects in clean samples and is the subject of still increasing experimental and theoretical activities (for a review see mit-review and references therein).
One possible mechanism invokes both, interaction and disorder and associates the insulator–metal transition with the melting of a pinned Wigner crystal when the carrier density increases and therewith $`r_\mathrm{s}`$ (the interaction energy in units of the kinetic energy) decreases, allowing for metallic behavior at intermediate $`r_\mathrm{s}`$. Only at a higher density, when $`r_\mathrm{s}`$ is small, the interaction becomes negligible and the now dominating disorder-induced Anderson localization leads back to insulating behavior. Such a scenario is supported by experiments hamilton and numerical investigations of few interacting particles in disordered lattice models benenti ; benenti\_mor ; waintal\_ann .
Another important experimental result in mesoscopic physics whose explanation necessitates to invoke disorder and interactions simultaneously is the value of the persistent current in diffusive rings pc\_exp . These persistent currents are much larger than the theoretical prediction for non-interacting electrons in disordered rings pc\_nonint . While the electron-electron interaction seems to play an essential role, the disorder in the sample is also important: Interactions cannot affect the persistent current in clean rotationally invariant 1d rings shastry ; axel , and the non-interacting result is consistent with the experimental one for a clean semiconductor ring in the ballistic regime mailly .
This has generated a large theoretical activity, dealing with the combined effect of interactions and disorder on the enhancement of persistent currents in mesoscopic rings (for an overview see e.g. Moriond ; Imry\_book ; eckern and references therein). Even though different theoretical approaches suggest an increase of the persistent current in disordered samples due to repulsive Coulomb interactions, a quantitative understanding of the experiments is still lacking.
Within continuous models of fermions in 1d with disorder, repulsive interactions are found to enhance the persistent currents, without axel and with cohen spin. On moderately disordered lattices, repulsive interactions are however found to decrease the persistent current for spinless fermions in 1d berko3 ; bouzerar ; kato . It was concluded from analytical considerations gs , using renormalization group arguments, that 1d lattice models exhibit an interaction-induced enhancement of the persistent current only when the spin degree of freedom is taken into account. Nevertheless, a numerical study of small systems abraham revealed that repulsive interactions can slightly enhance the averaged persistent current even for spinless fermions on a 1d lattice, provided the disorder is very strong.
A recent numerical investigation treating the persistent current for strongly disordered individual 1d chains schmitteckert ; WPSJ leads to the conclusion that, while Anderson localization is dominating the non-interacting case, the persistent current can be strongly enhanced by repulsive local interactions at sample dependent intermediate values of the interaction strength. This delocalization is accompanied by a reorganization of the ground state structure. At half filling, strong interactions induce a regular charge density in the Mott-insulator regime and decrease the persistent current.
This suppression of the persistent current in the limit of strong repulsive interactions is not limited to 1d models of spinless fermions at half-filling with local interactions. It occurs with long-range interactions berko3 at arbitrary filling and also in 2d lattices benenti , when strong interaction (large $`r_\mathrm{s}`$) leads to the formation of a Wigner crystal pinned by the disorder in the ground state.
In the case of the Hubbard model at half filling stafford and for spinless fermions in 1d chains without tsiper1 and with disorder WSJP , this suppression of the persistent current at strong interaction has been understood quantitatively from a perturbation theory starting at the Mott insulator limit. In contrast, numerical Hartree-Fock approaches which allow to treat larger systems than the exact diagonalization used in Refs. benenti ; berko3 and which have been used for weakly interacting particles in 1d and 2d bouzerar2 , are, at least in 1d, not able to quantitatively describe the persistent current in the limit of strong interaction piechon .
The sign of the persistent current in the case of spinless particles in strictly 1d rings is independent on the disorder realization and the interaction strength. According to a well-known theorem by Leggett leggett , the sign is given solely by the parity of the particle number $`N`$ (paramagnetic for $`N`$ even and diamagnetic for $`N`$ odd). This general rule is confirmed by an explicit calculation for a Luttinger liquid ring without disorder Loss .
On the other hand, the sign of the persistent current for particles with spin in disordered 1d rings or spinless fermions in disordered 2d systems differs from sample to sample and no general sign rule exists. Only in some special cases, like non-interacting particles with spin in clean 1d rings loss-gold , and for some interacting situations using the Hubbard model fye ; yu , the sign of the persistent current has been determined. The situation is less clear in the presence of long-range interactions and disorder we address in this work. For electrons with spin in 1d rings having a particular disorder consisting of only one barrier, a tendency towards diamagnetic responses, independent of the particle number, was found at strong repulsive interaction haeusler .
For spinless fermions in strongly disordered 2d lattice models, it has been noticed in numerical studies that the sign of the persistent current becomes realization-independent in the limit of strong interaction benenti ; benenti\_mor ; berko1 ; berko2 . Detailed studies of the local current show that the suppression of its transverse component by the interactions is much stronger than the decrease of the longitudinal current. On such a lattice, closed to a torus, the structure of the ground state at strong interaction is a Wigner crystal pinned by the disorder benenti ; benenti\_mor ; berko2 . While the system exhibits Anderson localization at weak interaction, the regime of intermediate interaction shows indications of a new type of correlated metal waintal .
For spinless fermions in 2d lattice models without disorder, the amplitude of the persistent current has recently been studied analytically and numerically at strong interaction tsiper2 . When $`r_s`$ is large, the hopping matrix elements between neighboring lattice sites being much smaller than the interaction strength, the behavior can be understood from a perturbation theory in terms of the hopping matrix elements.
In this paper we report a study of the persistent current in disordered 2d lattice models at very strong Coulomb interaction, using a perturbation theory expansion around the pinned Wigner crystal. While the understanding of this regime cannot directly explain the high persistent currents observed experimentally in diffusive metal rings, it may be relevant for the insulating side at low carrier density of the insulator-metal transition in 2d. We show how the sign of the persistent current at strong interaction follows systematically from the structure of the Wigner crystal and find simple rules for this sign. For the absolute values, power laws similar to the ones found in tsiper2 are obtained also in the disordered case. We shall show how the presence of disorder and spin can influence the prefactors and the exponents of these power laws.
In the following section we introduce the model for interacting fermions on a disordered lattice and the quantities used to characterize its properties. The perturbation theory is developed in section 3 and applied in section 4 to the persistent current in longitudinal and transverse direction before we conclude the paper.
## 2 Model
### 2.1 Hamiltonian
We consider $`N`$ fermions on a disordered square lattice with Coulomb interaction. The corresponding Hamiltonian reads
$$H=H_\mathrm{K}+H_\mathrm{D}+H_\mathrm{I}.$$
(1)
The kinetic energy term is
$$H_\mathrm{K}=t\underset{<i,i^{}>}{}\underset{\sigma }{}c_{i,\sigma }^+c_{i^{},\sigma }^{}$$
(2)
with the hopping matrix element $`t=1`$ setting the energy scale. We concentrate on rectangular 2d lattice structures with $`L_x\times L_y`$ sites $`i`$. The fermionic on-site operators $`c_{i,\sigma }c_{(x_i,y_i),\sigma }`$ destroy a particle with spin $`\sigma `$ located at $`\stackrel{}{r}_i=(x_i,y_i)`$, where the position coordinates $`x_i\{1,2,\mathrm{},L_x\}`$ and $`y_i\{1,2,\mathrm{},L_y\}`$ are measured in units of the lattice spacing $`a`$.
The disordered potential contribution is
$$H_\mathrm{D}=W\underset{i,\sigma }{}v_i\widehat{n}_{i,\sigma },$$
(3)
where $`W`$ is the disorder strength, with the independent random variables $`v_i`$, drawn from a box distribution within the interval $`[1/2;+1/2]`$. The occupation number operators are as usual given by $`\widehat{n}_{i,\sigma }=c_{i,\sigma }^+c_{i,\sigma }^{}`$. The Coulomb interaction is described by the term
$$H_\mathrm{I}=\frac{U}{2}\underset{\genfrac{}{}{0pt}{}{i,i^{}}{ii^{}}}{}\underset{\sigma ,\sigma ^{}}{}\frac{\widehat{n}_{i,\sigma }\widehat{n}_{i^{},\sigma ^{}}}{|\stackrel{}{r}_i\stackrel{}{r}_i^{}|}+U\underset{i}{}\frac{\widehat{n}_{i,}\widehat{n}_{i,}}{d},$$
(4)
the interaction strength being parametrized by $`U`$. The sum in $`H_\mathrm{K}`$ runs over all pairs of sites which are next neighbors on the lattice $`<i,i^{}>`$, while the interaction term is composed of a sum over all pairs of different sites. With these definitions, one finds $`r_\mathrm{s}=\frac{U}{2t\sqrt{\pi \nu }}`$, with the electronic density or filling factor $`\nu =N/L_xL_y=1/b^2`$, $`b`$ being the average distance between particles. An additional term takes into account double occupancy of a site by two particles of different spin, $`d(<a)`$ being a measure for the size of the on-site orbitals.
In order to study the persistent current, we close the 2d lattice first to a cylinder by imposing generalized periodic boundary conditions
$$c_{(L_x,y);\sigma }=\mathrm{exp}(\mathrm{i}\varphi _x)c_{(0,y);\sigma }$$
(5)
in $`x`$-direction. For $`\varphi _x=0`$, this is equivalent to usual periodic boundary conditions. Finite $`\varphi _x`$ accounts for a magnetic flux $`\mathrm{\Phi }=\varphi _x\mathrm{\Phi }_0/2\pi `$ threading the ring, $`\mathrm{\Phi }_0`$ being the flux quantum. We choose the units such that $`\mathrm{\Phi }_0/2\pi =1`$. In order to reduce finite size effects, we use periodic boundary conditions in $`y`$-direction, and thus obtain a torus topology with the fluxes $`(\varphi _x,\varphi _y)=(\varphi _x,0)`$. In section (4.3) we will use the dependence on the transverse flux $`\varphi _y`$ to study also the transverse current.
### 2.2 Persistent current
The magnetic flux threading the ring can drive a persistent current through the system. At zero temperature, it is given by
$$I(\varphi _x)=\frac{E_0}{\varphi }|_{\varphi =\varphi _x},$$
(6)
where $`E_0`$ is the many-particle ground state energy. Thus, the persistent current at $`T=0`$ is a measure of the dependence of the ground state energy on the magnetic flux. Since the latter can be expressed in the form of a boundary condition, it is at the same time a measure of the ground state sensitivity to the boundary conditions and can be related to the conductance of the sample kohn .
## 3 Theoretical approach
### 3.1 Wigner crystal at strong interaction
In the non-interacting limit, disorder leads to Anderson localization of the one-particle states and the problem can be treated by a perturbative expansion around the on-site localized states in terms of the hopping matrix elements $`t`$ bouch89 . Hopping to distant sites costs disorder energy of the order of $`W`$ such that a series expansion in $`t/W`$ results.
In the many-body case, strong repulsive interaction $`U`$ leads to Wigner crystallization of the ground state with on-site localized charges in the electro-statically most favorable position. One can use a similar perturbative formalism in terms of the hopping $`t`$, but now, the essential cost in energy caused by deplacing one of the many particles is given by the increase of the interaction energy such that one obtains a systematic expansion in terms of $`t/U`$.
We decompose the Hamiltonian (1) as
$$H=H_0+H_\mathrm{K}$$
(7)
with an unperturbed part containing disorder and interaction
$$H_0=H_\mathrm{D}+H_\mathrm{I},$$
(8)
and the perturbation given by the hopping terms of $`H_\mathrm{K}`$.
$`H_0`$ is composed of terms containing only occupation number operators in the one-particle on-site basis. Therefore, its $`N`$-particle eigenstates $`|\psi _\alpha `$ are Slater determinants built from $`N`$ different one-particle functions and are completely characterized by the occupation numbers $`n_{i,\sigma }(\alpha )\{0,1\}`$ of the one-particle states on site $`i`$ with spin $`\sigma `$, fulfilling the condition $`N=_{i,\sigma }n_{i,\sigma }(\alpha )`$. Therewith, the many-body eigenstates of $`H_0`$ can be written in the form
$$|\psi _\alpha =\left(\underset{i,\sigma }{}(c_{i,\sigma }^+)^{n_{i,\sigma }(\alpha )}\right)|0$$
(9)
($`|0`$ is the vacuum state), and the corresponding eigenenergies are given by $`E_\alpha =E_\alpha ^\mathrm{D}+E_\alpha ^\mathrm{I}`$ with
$`E_\alpha ^\mathrm{D}`$ $`=`$ $`W{\displaystyle \underset{i,\sigma }{}}v_in_{i,\sigma }(\alpha )\text{and}`$ (10)
$`E_\alpha ^\mathrm{I}`$ $`=`$ $`{\displaystyle \frac{U}{2}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{i,i^{}}{ii^{}}}{}}{\displaystyle \underset{\sigma ,\sigma ^{}}{}}{\displaystyle \frac{n_{i,\sigma }(\alpha )n_{i^{},\sigma ^{}}(\alpha )}{|\stackrel{}{r}_i\stackrel{}{r}_i^{}|}}+U{\displaystyle \underset{i}{}}{\displaystyle \frac{n_{i,}n_{i,}}{d}}.`$ (11)
The ground state $`|\psi _0`$ of this Coulomb glass problem is given by purely classical considerations, minimizing disorder and interaction energy. Its charge configuration depends in general on the specific disorder realization of the sample. At strong enough interaction, when the disorder effects are dominated by the interaction, the structure of $`|\psi _0`$ is the Wigner crystal of minimal interaction energy, independent of the disorder realization. The rigid array of charges can be translated as a whole through the system without changing the interaction energy. Nevertheless, the contribution of the disordered potential to the energy depends on the realization and pins the Wigner crystal in a realization dependent position. It is important to realize that the structure of the Wigner crystal is entirely given by the lattice geometry and the Coulomb interaction of the $`N`$ particles.
In contrast to $`H_0`$, the perturbing part $`H_\mathrm{K}`$ of the Hamiltonian depends on the magnetic flux through the ring since the latter appears in the boundary condition in $`x`$-direction and therefore influences some hopping matrix elements. Writing all the hopping terms explicitly, one obtains
$`H_\mathrm{K}(\varphi _x,\varphi _y)=t{\displaystyle \underset{\sigma }{}}({\displaystyle \underset{x=1}{\overset{L_x}{}}}{\displaystyle \underset{y=2}{\overset{L_y}{}}}c_{(x,y);\sigma }^+c_{(x,y1);\sigma }^{}`$
$`+{\displaystyle \underset{x=2}{\overset{L_x}{}}}{\displaystyle \underset{y=1}{\overset{L_y}{}}}c_{(x,y);\sigma }^+c_{(x1,y);\sigma }^{}+e^{\mathrm{i}\varphi _y}{\displaystyle \underset{x=1}{\overset{L_x}{}}}c_{(x,1);\sigma }^+c_{(x,L_y);\sigma }^{}`$
$`+e^{\mathrm{i}\varphi _x}{\displaystyle \underset{y=1}{\overset{L_y}{}}}c_{(1,y);\sigma }^+c_{(L_x,y);\sigma }^{}+\text{H.C.}).`$ (12)
### 3.2 Perturbation expansion for spinless fermions
For electrons with spin, the unperturbed ground state is $`2^N`$-fold degenerate since all the spin configurations yield the same energy when hopping is suppressed. Further degeneracies appear in the case of clean systems when $`W=0`$ allows for translational symmetry.
For simplicity we first treat the case of completely spin polarized systems (all spins up, equivalent to spinless fermions) with disorder where the ground state of $`H_0`$ is not degenerate and the expansion in $`H_\mathrm{K}`$ is straightforward.
Using standard perturbation theory, the correction to the ground state energy in $`n^{\mathrm{th}}`$ order is given by
$`E_0^{(n)}={\displaystyle \underset{\alpha _1,\alpha _2,\mathrm{},\alpha _{n1}}{}}`$ (13)
$`{\displaystyle \frac{\psi _0|H_\mathrm{K}|\psi _{\alpha _1}\psi _{\alpha _1}|H_\mathrm{K}|\psi _{\alpha _2}\mathrm{}\psi _{\alpha _{n1}}|H_\mathrm{K}|\psi _0}{(E_0E_{\alpha _1})(E_0E_{\alpha _2})\mathrm{}(E_0E_{\alpha _{n1}})}},`$
with the sums running over all the eigenstates of $`H_0`$ except the ground state itself.
The numerator of the contributions to the sum of equation (13) contains matrix elements $`\psi _{\alpha _k}|H_\mathrm{K}|\psi _{\alpha _{k+1}}`$ of the perturbing Hamiltonian. Since $`H_\mathrm{K}`$ consists only of one-particle hopping terms, non-zero matrix elements can arise only if the two states $`|\psi _{\alpha _k}`$ and $`|\psi _{\alpha _{k+1}}`$ differ by nothing else than a single hop of one of the particles. From (13) one sees that a finite contribution to the sum over different sequences of intermediate states $`\alpha `$ is obtained only if the $`n`$ one-particle hops are such that the final configuration has an overlap with the initial one, corresponding to the ground state. The $`n`$ sums over intermediate states $`\alpha _k`$ in equation (13) can then be rewritten as a sum over all the sequences $`S=(\alpha _1,\alpha _2,\mathrm{},\alpha _{n1})`$ which give a nonzero contribution. Denoting the numerator of the terms by $`\mathrm{Num}(S)`$ and the denominator by $`\mathrm{Den}(S)`$, equation (13) takes the form
$$E_0^{(n)}=\underset{S}{}\frac{\mathrm{Num}(S)}{\mathrm{Den}(S)}.$$
(14)
We will now evaluate the numerator and the denominator separately. The numerator $`\mathrm{Num}(S)`$ can be calculated directly from the $`n`$ hopping matrix elements, thereby taking into account the flux dependent phase for hops crossing the boundary. Since we consider fermions, the corresponding operators anti-commute and their order in the products (9) defining the basis states $`|\psi _\alpha `$ is crucial for the sign. When the starting point (the ground state configuration) is re-established after $`n`$ consecutive hopping processes, the order of the operators can be modified. Then, the sign of the permutation $`P_S`$ of the operators, caused by the sequence of one-particle hops $`S`$, must be incorporated in the result. Altogether, one obtains
$$\mathrm{Num}(S)=\mathrm{sign}(P_S)(t)^n\mathrm{exp}\left[\mathrm{i}\varphi _x(h_\mathrm{f}h_\mathrm{b})\right].$$
(15)
$`h_\mathrm{f}`$ and $`h_\mathrm{b}`$ denote the number of hoppings across the boundary between sites $`(L_x,y)`$ and $`(1,y)`$ in forward and backward direction, respectively. Therefore, only the corrections to the ground state energy due to sequences with $`h_\mathrm{f}h_\mathrm{b}0`$ are flux dependent.
Moving particles creates defects in the Wigner crystal. This increases the interaction energy (see equation 11) of the ground state $`E_0^\mathrm{I}`$ by the amount $`Uϵ_\alpha `$, where
$$ϵ_\alpha =\frac{1}{U}\left(E_\alpha ^\mathrm{I}E_0^\mathrm{I}\right)$$
(16)
is non-negative and independent of $`U`$. It accounts for the difference of the inter-particle distances between the configurations of the state $`|\psi _\alpha `$ and the ground state $`|\psi _0`$. We assume always $`UW`$, such that the difference in potential energy can be neglected in a first step except for the case $`ϵ_\alpha =0`$ which occurs for some special sample geometries and particle numbers. Corrections due to the disorder will be considered in Section (4.2).
In the generic case when $`ϵ_\alpha >0`$ for all intermediate states, the energy differences in the denominator $`\mathrm{Den}(S)`$ are dominated by the difference in interaction energy. In the limit of strong interaction, one therefore gets
$$\mathrm{Den}(S)(U)^{n1}\underset{\alpha =1}{\overset{n1}{}}ϵ_\alpha .$$
(17)
## 4 Persistent current
### 4.1 Longitudinal current
Since the longitudinal persistent current at zero temperature $`I=E_0/\varphi _x`$ is given by the flux dependence of the ground state energy, a perturbative treatment of the latter in terms of the hopping $`t`$ yields a systematic expansion of the persistent current.
Calculating the $`n^{\mathrm{th}}`$ order correction $`E_0^{(n)}(\varphi _x)`$ to the ground state energy, one gets the correction to the persistent current
$$I^{(n)}(\varphi _x)=\frac{E_0^{(n)}}{\varphi _x}$$
(18)
in $`n^{\mathrm{th}}`$ order in the perturbation $`H_\mathrm{K}`$.
#### 4.1.1 Relevant terms of the perturbation theory
When the sequence $`S`$ of $`n`$ hopping elements is chosen such that each of the particles returns to its initial position without completing a tour around the ring, every particle which has crossed the boundary must necessarily cross it a second time in the opposite direction such that $`h_\mathrm{b}=h_\mathrm{f}`$. Therefore the flux dependence of these contributions disappears and they cannot influence the persistent current. The lowest order of the perturbation theory which yields a finite contribution $`I^{(n)}`$ is $`n=L_x`$, corresponding to the sequences in which one particle starting at $`(x_0,y_0)`$ crosses the boundary and returns to its original position after completion of its journey around the ring at constant $`y=y_0`$. If there are more than one particle in the line of the lattice with constant $`y=y_0`$, a contribution of the same order arises from sequences of hops which move each of the particles to the position of its neighbor (see Fig. 1). Since any hopping in $`y`$-direction leads to an increase of the order of the contribution, the lowest order correction to the persistent current is given by the considered processes in which $`y`$ is kept constant and all of the hops are either in forward or in backward direction.
We now address the dependence of the permutation $`P_S`$ corresponding to this kind of process on the number $`N_y`$ of particles in the line at constant $`y`$. If $`N_y=1`$, the particle returns to its initial site, the final configuration is exactly equal to the initial one and $`\mathrm{sign}(P_S)=+1`$. In the general case of arbitrary $`N_y`$, the considered sequences of hops (see Fig. 1) lead to a cyclic rotation of the order of the $`N_y`$ particles and the sign of the corresponding permutation is $`\mathrm{sign}(P_S)=(1)^{N_y1}`$.
This yields the result for the lowest-order $`\varphi _x`$-dependence of the ground state energy
$$E_0^{(L_x)}(\varphi _x)\underset{S}{}\frac{t^{L_x}\mathrm{sign}(P_S)\mathrm{exp}\left[\mathrm{i}\varphi _x(h_\mathrm{f}h_\mathrm{b})\right]}{U^{L_x1}ϵ_{\alpha _1}ϵ_{\alpha _2}\mathrm{}ϵ_{\alpha _3}}.$$
(19)
Each sequence considered in this sum contains either $`L_x`$ forward hops with $`h_\mathrm{f}=1`$; $`h_\mathrm{b}=0`$ or it contains $`L_x`$ backward hops with $`h_\mathrm{f}=0`$; $`h_\mathrm{b}=1`$. Each given backward sequence $`S_\mathrm{b}`$ can now be assigned to the forward sequence $`S_\mathrm{f}`$ with the reversed order of hops whose contribution differs only in the sign of the flux-dependent phase-factor. One can express the result as a sum over the forward hopping sequences
$$E_0^{(L_x)}(\varphi _x)2\frac{t^{L_x}}{U^{L_x1}}\underset{S_\mathrm{f}}{}\frac{\mathrm{sign}(P_S)\mathrm{cos}\varphi _x}{ϵ_{\alpha _1}ϵ_{\alpha _2}\mathrm{}ϵ_{\alpha _{L_x1}}}.$$
(20)
Within the above perturbation theory, when processes corresponding to two loops around the ring (which are at least of order $`2L_x`$) are neglected as compared to the lowest order one-loop processes, the flux dependence of the ground state energy is harmonic and $`2\pi `$-periodic in $`\varphi _x`$.
#### 4.1.2 Lowest order result for the persistent current
From (20), one obtains the persistent current in $`L_x`$-th order perturbation theory
$$I^{(L_x)}(\varphi _x)=\stackrel{~}{I}^{(L_x)}\mathrm{sin}\varphi _x$$
(21)
with the flux–independent amplitude
$$\stackrel{~}{I}^{(L_x)}2\frac{t^{L_x}}{U^{L_x1}}\underset{S_\mathrm{f}}{}\frac{\mathrm{sign}(P_S)}{ϵ_{\alpha _1}ϵ_{\alpha _2}\mathrm{}ϵ_{\alpha _{L_x1}}}.$$
(22)
This result contains several interesting features. First, the absolute value of the persistent current decays proportionally to $`t^{L_x}/U^{L_x1}`$, in the limit of strong interaction. In order to determine the constant prefactor of this power law, it is sufficient to figure out all possible processes which transform the ground state into itself, using $`L_x`$ forward hopping processes, and to calculate the corresponding $`ϵ_\alpha `$ from (16).
Furthermore, the sign of the dominating contributions to the persistent current in the limit of strong interaction is given by $`\mathrm{sign}(P_S)`$, which for spin-polarized electrons is given by $`(1)^{N_y}`$ with the number $`N_y`$ of electrons in the line of the sample at constant $`y`$. This is consistent with the well-known theorem by Leggett for the sign of the persistent current of spinless fermions in 1d leggett (positive $`\stackrel{~}{I}`$ or paramagnetic response for $`N`$ even and negative or diamagnetic for $`N`$ odd). Only if the unperturbed ground state (Wigner crystal) configuration of the particles is such that the particle numbers $`N_y`$ in different occupied lines $`y`$ have different parity, the prefactors of the corresponding terms have to be considered to determine the sign of the persistent current. For $`N_y`$ particles in a line of length $`L_x`$, the number of hopping sequences $`N_{\mathrm{seq}}(N_y)`$ going from the ground state to itself is the number of terms contributing to the sum in equation (22). Therewith, the result for the persistent current can be roughly estimated to be
$$\stackrel{~}{I}^{(L_x)}\frac{t^{L_x}}{U^{L_x1}}\underset{y=1}{\overset{L_y}{}}N_{\mathrm{seq}}(N_y)(1)^{N_y}.$$
(23)
In the limit of low filling $`N_y/L_x0`$, $`N_{\mathrm{seq}}(N_y)`$ is approximatively given by the number of possibilities $`N_{\mathrm{seq}}(N_y)`$ for $`N_y`$ particles to each make $`L_x/N_y`$ (here we assume for simplicity that $`L_x/N_y`$ is an integer) forward jumps to reach the position of its neighbor, leading to the estimate $`N_{\mathrm{seq}}(N_y)L_x!/[(L_x/N_y)!]^{N_y}`$. At finite filling, one must consider that the neighbor particle must have left its starting site before the arriving particle can do its last hop. This correlation of the order of the hops of different particles reduces $`N_{\mathrm{seq}}(N_y)`$. With increasing filling, $`N_{\mathrm{seq}}(N_y)`$ starts to exponentially increase with $`N_y`$ and continues to increase more slowly until $`N_y=L_x/2`$ (half filling). For larger filling it decreases, thereby obeying a symmetry with respect to half filling which is a consequence of the symmetry between particles and holes. One can expect that the contribution of the line $`N_y=N_y^{\mathrm{max}}`$ with the largest number $`N_{\mathrm{seq}}(N_y^{\mathrm{max}})`$ of sequences (at low filling this is the one with the maximum number of particles) may dominate over the contributions of the ones with fewer sequences. The sign of the persistent current is then likely to be $`(1)^{N_y^{\mathrm{max}}}`$.
#### 4.1.3 Examples
As an example, we calculate explicitly the lowest order term in $`\frac{1}{U}`$ of the persistent current for 4 spinless fermions on a few small $`L_x\times L_y`$ rectangular lattices, using the formula (22).
##### $`4\times 2`$ sites
We start with the simple case of a $`4\times 2`$ lattice. The electro-statically lowest energy configuration (the Wigner crystal) is shown in Fig. 2.
The number of particles in the two lines at $`y=1`$ and $`y=2`$ of the system is $`N_y=2`$, and the lowest order of the perturbation theory which yields a contribution to the persistent current is $`n=L_x=4`$. From these two ingredients, we can immediately determine the sign ($`\mathrm{sign}(P_S)=1`$ for $`N_y`$ even) and the power law of the decrease of the persistent current at strong interaction strength. The leading term of the amplitude (22) is given by
$$\stackrel{~}{I}^{(4)}2\frac{t^4}{U^3}\underset{S_\mathrm{f}}{}\frac{1}{ϵ_{\alpha _1}ϵ_{\alpha _2}\mathrm{}ϵ_{\alpha _{L_x1}}}.$$
(24)
This result is always positive and $`1/U^3`$. Therefore, the response of the system to the applied flux is always paramagnetic at strong interaction.
For this example, it is an easy exercise to figure out all hopping sequences which contribute in lowest order. For each line $`y=1,2`$, one finds 4 different sequences and can explicitly calculate the interaction energies and $`ϵ_\alpha `$ of the intermediate states. This allows to evaluate the sum over all sequences in (24) with the result
$$\stackrel{~}{I}\stackrel{~}{I}^{(4)}853\frac{t^4}{U^3}\text{for}\frac{t}{U}0.$$
(25)
In the case of 8 spinless fermions on $`4\times 4`$ sites, the same analysis can be carried out except for the fact that now four lines $`y=1,2,3,4`$ must be considered, yielding an additional factor of 2 in (25). In particular, the persistent current is always paramagnetic in the strong interaction limit, as found numerically in Ref. berko2 .
##### $`2\times 4`$ sites
The situation for 4 particles on $`2\times 4`$ sites is even simpler since there are four lines $`y=1,2,3,4`$, each of them containing $`N_y=1`$ particle. Since $`L_x=2`$, the leading order of the perturbation theory at strong interaction is $`n=2`$, and the sequences of hops contain only one intermediate state. The evaluation of (22) yields
$$\stackrel{~}{I}\stackrel{~}{I}^{(2)}15\frac{t^2}{U}\text{for}\frac{t}{U}0,$$
(26)
and the persistent current is diamagnetic at strong interaction.
##### $`6\times 6`$ sites
4 particles on $`6\times 6`$ sites is the situation investigated numerically in Ref. benenti , at strong disorder. As a function of the interaction strength, an increase of the average persistent current at intermediate strength and a decrease at strong interaction was found. An exponential dependence on the interaction strength was fitted to the data for not too strong interaction. Furthermore, it was noticed that, at strong interaction, the persistent current became paramagnetic for all samples, independent of the disorder realization.
The Wigner crystal ground state on such a lattice is of square form, as shown in Fig. 1. Thus, there are two lines with $`N_y=2`$ in which hopping sequences of $`L_x=6`$ hops can contribute to the persistent current. This explains immediately that the response is paramagnetic, since $`N_y`$ is even and $`\mathrm{sign}(P_S)=+1`$. The persistent current decreases in the strong interaction limit as $`t^6/U^5`$. A laborious evaluation of all the 18 sequences for each of the lines with $`N_y=2`$ yields the result
$$\stackrel{~}{I}\stackrel{~}{I}^{(6)}1.808\times 10^6\frac{t^6}{U^5}\text{for}\frac{t}{U}0.$$
(27)
As compared to the previous cases, the prefactor is much larger. This is caused by the bigger number of contributing sequences and, more importantly, the difference in interaction energy between the intermediate states and the unperturbed ground state $`ϵ_\alpha `$. The latter can be very small when the distance between the particles is large.
A numerical investigation by direct diagonalization of the corresponding Hamiltonian matrices SW\_unpub confirms the signs, the power laws and the numerical prefactors predicted by the above formulas, also for the last case of 4 particles on $`6\times 6`$ sites, where the persistent current is indeed found to follow the power law of (27) at strong interaction. However, it must be noticed that the sign of the persistent current is well established benenti at interaction values much lower than the ones where the agreement in amplitude with our formula starts to be good. Even though the data always follow the power laws at strong interaction, fitting an exponential interaction dependence, as done in Ref. benenti , is possible at moderately strong interaction, and might allow to extract useful informations in the regime where higher order terms of the perturbation theory are non-negligible.
### 4.2 Disorder effects at strong interaction
The role of the disorder and realization-dependent fluctuations of the persistent current vanish in the limit of strong interaction, when $`W/U0`$. In this section, we treat the lowest order correction in $`W/U`$ to the results presented above.
In addition, we consider the special case of a perfectly clean system $`W=0`$, in which the translational symmetry can considerably influence the interaction dependence of the persistent current.
#### 4.2.1 Disorder corrections to the persistent current
In order to take into account the disorder corrections in the perturbation theory for the longitudinal persistent current, it is not sufficient to consider in the denominator of the expansion terms the $`ϵ_\alpha `$ which completely neglect the disorder energy. Instead, the full energy differences $`E_0E_\alpha =U\stackrel{~}{ϵ}_\alpha `$ with
$$\stackrel{~}{ϵ}_\alpha =\frac{1}{U}\left(E_\alpha ^\mathrm{I}+E_\alpha ^\mathrm{D}E_0^\mathrm{I}E_0^\mathrm{D}\right)$$
(28)
account also for the differences in disorder energy between the intermediate states $`\alpha `$ and the ground state. The difference
$$\stackrel{~}{ϵ}_\alpha ϵ_\alpha =\frac{1}{U}\left(E_\alpha ^\mathrm{D}E_0^\mathrm{D}\right)=\frac{W}{U}\underset{i}{}v_i\left(n_i(\alpha )n_i(0)\right),$$
(29)
is of the order $`W/U`$ and vanishes in the limit $`W/U0`$. With the definition
$$d_\alpha =\underset{i}{}v_i\left(n_i(\alpha )n_i(0)\right)$$
(30)
we can write $`\stackrel{~}{ϵ}_\alpha =ϵ_\alpha +\frac{W}{U}d_\alpha `$ and therewith express the energy difference terms in the denominator of (13) as
$$\frac{1}{E_0E_\alpha }=\frac{1}{U\stackrel{~}{ϵ}_\alpha }=\frac{1}{Uϵ_\alpha \left(1+\frac{W}{Uϵ_\alpha }d_\alpha \right)}.$$
(31)
Taking the lowest order term in $`W/U`$, and averaging over the ensemble yields
$$\frac{1}{E_0E\alpha }\frac{1}{Uϵ_\alpha }\left(1\frac{W}{Uϵ_\alpha }<d_\alpha >\right),$$
(32)
where the brackets $`<\mathrm{}>`$ denote the ensemble average over all disorder realizations. At first glance one could expect that the correction linear in $`W/U`$ vanishes because $`<d_\alpha >=0`$, when the disorder average is taken over all values of $`d_\alpha `$. However, since the ground state is the Wigner crystal pinned at the lowest disorder configuration, $`d_\alpha `$ is more likely positive than negative. The first correction is thus linear in $`W`$, and since $`ϵ_\alpha `$ is always positive, the lowest order correction to the persistent current due to disorder is reducing the contributions of all sequences, the result being
$`I_\mathrm{W}^{(L_x)}`$ $``$ $`2{\displaystyle \frac{t^{L_x}}{U^{L_x1}}}{\displaystyle \underset{S_\mathrm{f}}{}}{\displaystyle \frac{\mathrm{sign}(P_S)}{ϵ_{\alpha _1}ϵ_{\alpha _2}\mathrm{}ϵ_{\alpha _{L_x1}}}}`$ (33)
$`\times `$ $`\left(1{\displaystyle \underset{\alpha }{}}{\displaystyle \frac{1}{ϵ_\alpha }}{\displaystyle \frac{W<d_\alpha >}{U}}+O[(W/U)^2]\right).`$
#### 4.2.2 The particular case $`W=0`$
At $`W=0`$ we have $`H_0=H_I`$, and the ground state at $`t=0`$ is degenerate because of the translation symmetry. The number $`n_\mathrm{D}`$ of equivalent degenerate Wigner crystal positions and corresponding basis states $`|\psi _0^{(\beta )}`$ depends on the system size and the number of particles. The hopping terms however lead to a coupling of these degenerate basis states, and split the degenerate levels, except for special situations where symmetries persist. The coupling terms themselves can be expressed in terms of a perturbative expansion in $`t`$. The order $`p`$ in which the degeneracy of the ground state is lifted may be different from the lowest order $`n`$ which yields a finite persistent current. We address in the following the three different cases which may occur.
##### $`p>n`$
In this case, the splitting can be ignored as compared to the persistent current at strong interaction. The perturbation theory can be applied as in the disordered case, using one of the degenerate ground state configurations (the result will be the same for arbitrary superpositions of the degenerate basis states).
##### $`p=n`$
When the persistent current is given by terms which are of the same order as the ones which lift the ground state degeneracy, the flux dependent correction to the ground state energy and the persistent current are given by the lowest eigenvalue of the effective coupling matrix $`M`$ between the $`n_\mathrm{D}`$ degenerate basis states. The matrix elements are
$`M_{\beta ,\beta ^{}}={\displaystyle \underset{\alpha _1,\alpha _2,\mathrm{},\alpha _{p1}}{}}`$ (34)
$`{\displaystyle \frac{\psi _0^{(\beta )}|H_\mathrm{K}|\psi _{\alpha _1}\psi _{\alpha _1}|H_\mathrm{K}|\psi _{\alpha _2}\mathrm{}\psi _{\alpha _{p1}}|H_\mathrm{K}|\psi _0^{(\beta ^{})}}{(E_0E_{\alpha _1})(E_0E_{\alpha _2})\mathrm{}(E_0E_{\alpha _{p1}})}}.`$
Since in this case all matrix elements are $`t^n/U^{n1}`$, the lowest eigenvalue and thus the persistent current will follow such a power law too.
##### $`p<n`$
When the levels are split at an order which is lower than the one which yields contributions to the persistent current, the higher order terms must be calculated using the lowest energy ground state found from the diagonalization of $`M`$. Such a ground state will be a superposition
$$|\psi _0=\underset{\beta =1}{\overset{n_\mathrm{D}}{}}f_\beta |\psi _0^{(\beta )}\text{with}\underset{\beta }{\overset{n_\mathrm{D}}{}}|f_\beta |^2=1$$
(35)
of the different Wigner crystal positions. Plugging this ground state into the general expression for the corrections to the ground state energy (13), one gets
$`E_0^{(n)}(W=0)={\displaystyle \underset{\beta ,\beta ^{}}{}}f_\beta ^{}f_\beta ^{}^{}{\displaystyle \underset{\alpha _1,\alpha _2,\mathrm{},\alpha _{n1}}{}}`$ (36)
$`{\displaystyle \frac{\psi _0^{(\beta )}|H_\mathrm{K}|\psi _{\alpha _1}\psi _{\alpha _1}|H_\mathrm{K}|\psi _{\alpha _2}\mathrm{}\psi _{\alpha _{n1}}|H_\mathrm{K}|\psi _0^{(\beta ^{})}}{(E_0E_{\alpha _1})(E_0E_{\alpha _2})\mathrm{}(E_0E_{\alpha _{n1}})}},`$
and realizes that contributions can arise from hopping processes starting at any of the ground state components $`|\psi _0^{(\beta )}`$, and ending at an arbitrary $`|\psi _0^{(\beta ^{})}`$ where $`\beta `$ and $`\beta ^{}`$ can be equal or different. The processes with $`\beta \beta ^{}`$ can in some special cases have a lower power in $`\frac{1}{U}`$ than the processes with $`\beta =\beta ^{}`$, relevant in the disordered case, and therefore dominate in the limit of strong interaction.
The case $`W=0`$ can be qualitatively different from the disordered case $`W0`$, even at very weak disorder. The reason for this is the fact that the unperturbed ground state is a superposition of the degenerate positions of the Wigner crystal. This degeneracy is lifted by the disorder, and the typical energy difference between two positions is of the order $`W\sqrt{N}`$ (this arises from the $`N`$ different random on-site energies, all being of order $`W`$). The coupling between these electro-statically equivalent states is provided by the hopping terms and can itself be estimated from a perturbation theory. The lowest order coupling needs a number of hops which is typically $`p=N`$ since the Wigner crystal must in most cases be translated as a whole and each of the particles must hop at least once. In some cases, a $`p<N`$ can be sufficient, as in the example of 3 particles on $`4\times 2`$ sites (see below).
Then, the coupling decreases at strong interaction as $`t^p/U^{p1}`$ and becomes, in the limit of strong interaction $`U`$, always much smaller than the splitting due to the disorder. This holds at any finite disorder value (provided $`p>1`$), and prevents a mixing of the different Wigner crystal configurations in the ground state. Only for $`W=0`$, the splitting by the disorder vanishes, and an infinitesimal coupling leads to the superposition of the different configurations.
#### 4.2.3 Example for disorder corrections
##### $`4\times 2`$ sites, 4 particles
The disorder dependence of the flux sensitivity is illustrated for the example of 4 particles in $`(L_x,L_y)=(4,2)`$. One expects an $`1/U^3`$ dependence for the leading term of the persistent current (24) and, according to (33), a $`W/U^4`$ dependence for the linear correction in disorder. We have performed numerical calculations at strong interaction SW\_unpub , which are in good agreement with these power laws predicted by the lowest order of the perturbation theory. The persistent current decreases when the strength of the disorder is increased, as predicted by the negative correction in formula (33).
##### Effects of degeneracies at $`W=0`$
The clean case however has a very different amplitude due to the degeneracies in the Wigner crystal. The persistent current can be calculated from (36), and after considering all the possible sequences between the different degenerate configurations, one obtains $`\stackrel{~}{I}^{(4)}(W=0)2.1238\stackrel{~}{I}^{(4)}`$, with $`\stackrel{~}{I}^{(4)}`$ taken from (24). A numerical check SW\_unpub confirms this result at strong interaction.
##### $`4\times 2`$ sites, 3 particles, $`W=0`$
In order to show that the absence of disorder does not only change the prefactor of the interaction dependence at large $`U`$, but can even influence the power of the $`U`$-dependence, we consider 3 particles in $`(L_x,L_y)=(4,2)`$. The particle which is alone in one line has two possible positions which are degenerate at $`t=W=0`$, and there exists a sequence of only two one-particle hops which links the two degenerate states with each other (see Fig. 3). Since the ground state at $`W=0`$ is a superposition of all possible degenerate configurations, this process leads to a contribution $`t^2/U`$ to the persistent current which dominates at strong interaction. In the disordered case, in contrast, the degeneracy of the ground state configurations is lifted and only sequences consisting of at least four hops can link the ground state to itself. Then, the interaction dependence of the lowest order contribution to the persistent current is $`t^4/(WU^2)`$. Interestingly, in this example, one of the intermediate states touched by a sequence of 4 hops can be a different, but electro-statically equivalent configuration of the Wigner crystal, and $`ϵ_\alpha =0`$, such that even at $`UW`$, the disorder energy cannot be neglected.
A numerical check SW\_unpub shows that by adding the disorder, the asymptotic dependence indeed follows this prediction and goes from $`U^1`$ to $`U^2`$.
### 4.3 Transverse current
#### 4.3.1 General considerations
Without disorder, the current in $`y`$-direction must vanish because of the symmetry of the problem (no flux is applied in $`y`$-direction). While this remains true in the ensemble average, the presence of disorder breaks this symmetry in individual samples. A finite transverse current can be observed in addition to the longitudinal current driven by the flux $`\varphi _x`$. The transverse current can be calculated from
$$I_t(\varphi _x)=\frac{E_0(\varphi _x,\varphi _y)}{\varphi _y}|_{\varphi _y=0},$$
(37)
such that a theoretical approach to determine $`I_t(\varphi _x)`$ consists in using the perturbative expansion for the ground state energy, as for the longitudinal current, but with a finite $`\varphi _y`$. Then, the derivative with respect to $`\varphi _y`$ at $`\varphi _y=0`$ can be evaluated for the leading contribution in the strong interaction regime. Therefore, only terms in the series (13) which depend on $`\varphi _y`$ can make a finite contribution to the transverse current.
In order to get such contributions, one needs hopping sequences in $`y`$-direction which contain hops crossing the boundary between sites $`(x,L_y)`$ and $`(x,1)`$ with $`h_\mathrm{u}h_\mathrm{d}0`$. $`h_\mathrm{u}`$ and $`h_\mathrm{d}`$ denote the number of such hops from $`y=L_y`$ to $`y=1`$ (’upwards’) and from $`y=1`$ to $`y=L_y`$ (’downwards’).
However, the flux dependence of the energy due to sequences containing $`L_y`$ upward or $`L_y`$ downward hoppings is $`\mathrm{cos}\varphi _y`$, as for the longitudinal processes (20), which are $`\mathrm{cos}\varphi _x`$. As a consequence, the resulting transverse current from these sequences is $`\mathrm{sin}\varphi _y`$, and vanishes when $`\varphi _y=0`$ is taken, just like the longitudinal current (21) at $`\varphi _x=0`$. Therefore, no transverse current can be obtained in the order $`L_y`$ of the perturbation theory.
Nevertheless, higher order sequences exist which do lead to finite contributions. These are due to sequences of hops which cross both, the boundary in $`x`$, and the boundary in $`y`$-direction. The lowest order of the expansion in which transverse currents appear is $`m=L_x+L_y`$, corresponding to sequences which consist of a combination of only forward and downward hops, only backward and upward hops, or vice versa, which complete one round in both spatial dimensions of the system. An example of such a sequence is shown in Fig. 4.
As before, each of the sequences containing only forward and upward hops $`S_{\mathrm{fu}}`$ can be associated with the reverse process consisting of backward and downward hops $`S_{\mathrm{bd}}`$. The contribution of the reverse process differs only in the sign of the flux-dependent phase factors. The correction to the ground state energy (13) of all $`S_{\mathrm{fu}}`$ and $`S_{\mathrm{bd}}`$ sequences together can be written as
$$E_{0,\mathrm{fu}+\mathrm{bd}}^{(m)}=2\frac{t^m}{U^{m1}}\mathrm{cos}(\varphi _x+\varphi _y)\underset{S_{\mathrm{fu}}}{}\frac{\mathrm{sign}(P_{S_{\mathrm{fu}}})}{\stackrel{~}{ϵ}_{\alpha _1}\mathrm{}\stackrel{~}{ϵ}_{\alpha _{m1}}}.$$
(38)
The same association can be made for sequences containing only forward and downward hops $`S_{\mathrm{fd}}`$, which are regrouped with the corresponding reverse processes consisting of backward and upward hops, yielding
$$E_{0,\mathrm{fd}+\mathrm{bu}}^{(m)}=2\frac{t^m}{U^{m1}}\mathrm{cos}(\varphi _x\varphi _y)\underset{S_{\mathrm{fd}}}{}\frac{\mathrm{sign}(P_{S_{\mathrm{fd}}})}{\stackrel{~}{ϵ}_{\alpha _1}\mathrm{}\stackrel{~}{ϵ}_{\alpha _{m1}}}.$$
(39)
#### 4.3.2 Lowest order result for the transverse current
By taking the derivative with respect to $`\varphi _y`$ at $`\varphi _y=0`$, we obtain for the transverse current in $`m`$-th order of the perturbation theory
$`I_t^{(m)}(\varphi _x)=2{\displaystyle \frac{t^m}{U^{m1}}}\mathrm{sin}\varphi _x`$ (40)
$`\times \left({\displaystyle \underset{S_{\mathrm{fu}}}{}}{\displaystyle \frac{\mathrm{sign}(P_{S_{\mathrm{fu}}})}{\stackrel{~}{ϵ}_{\alpha _1}\mathrm{}\stackrel{~}{ϵ}_{\alpha _{m1}}}}{\displaystyle \underset{S_{\mathrm{fd}}}{}}{\displaystyle \frac{\mathrm{sign}(P_{S_{\mathrm{fd}}})}{\stackrel{~}{ϵ}_{\alpha _1}\mathrm{}\stackrel{~}{ϵ}_{\alpha _{m1}}}}\right),`$
with $`m=L_x+L_y`$ being the lowest order which yields a non-zero contribution. Even when $`UW`$, the disorder energy cannot be neglected as in the lowest order term for the longitudinal current, since the two contributions to the current in (40) cancel each other exactly when the disorder vanishes. It can also be seen that the lowest order result for the transverse current is proportional to $`\mathrm{sin}\varphi _x`$ and vanishes with the longitudinal current at $`\varphi _x=0`$.
#### 4.3.3 Example
In the case of 4 particles in a lattice of $`6\times 6`$ sites, we have seen above (see equation (27)) that the longitudinal persistent current decreases at strong interaction following the power law $`t^6/U^5`$. The suppression of the transverse current is much more pronounced, since the leading contribution at strong interaction is of the order $`L_x+L_y=12`$, and our theory (40) predicts that it decays $`t^{12}/U^{11}`$. This means that the particle mobility is more and more restricted to the longitudinal direction when the interaction is increased. That the suppression of the transverse current is much stronger than the one of the longitudinal persistent current has been noticed in the numerical study of Ref. benenti . The dominance of the longitudinal current at strong interaction also explains the observation of “plastic flow” in a study of the local persistent currents berko2 . The orientation of the local currents has been proposed as a signature of the insulator-metal transition occurring in interacting 2d systems benenti\_mor .
### 4.4 Persistent current for electrons with spin
In this section we present some examples of the effects of the spin of the electrons on the longitudinal persistent current, showing that our approach is not restricted to spinless fermions, but can also be used for electrons carrying spin.
For finite systems containing $`N`$ particles, it is well known that the spin polarization of the ground state can depend on the magnetic flux due to level crossings between states of different spin symmetry as a function of the flux haeusler . However, without an external magnetic field acting on the electrons in the ring, there is a degeneracy related to the operator $`S_z`$ in the subspace of fixed total spin and one can write $`E(S^2,S_z)=E(S^2)`$.
It is in principle possible (though difficult experimentally) to create a magnetic flux through the ring while the magnetic field remains vanishing at the positions of the electrons in the ring. Such a flux does not lift the spin degeneracy mentioned above. If we now want to follow the dependence of the ground state energy on the flux, we can restrict the study to the subspace of minimum absolute value of $`S_z`$, choosing $`S_z=0`$ for an even number of particles with $`N_{}=N/2`$ spins up and $`N_{}=N/2`$ spins down, or $`S_z=1/2`$ for an odd number of particles with $`N_{}=\frac{N+1}{2}`$ spins up and $`N_{}=\frac{N1}{2}`$ spins down.
#### 4.4.1 One-dimensional systems
In 1d, a general rule exists for the sign of the persistent current in the case of spinless fermions leggett . This rule is valid at arbitrary disorder and interaction. Below, we address the question of the sign of the persistent current for electrons with spin in different parameter regimes.
##### Non-interacting electrons with spin
In the case of non-interacting electrons with spin, one can separately consider the flux dependence of the energy of the particles with up spin and the one of the particles with down spin, and add up their contributions to the persistent current.
If $`N`$ is even, the contribution of the $`N/2`$ up-spins to the persistent current will have the same sign $`(1)^{N/2}`$ as the contribution of the $`N/2`$ down-spins. This determines the sign of the persistent current. For odd $`N`$, the contribution of the $`(N+1)/2`$ up spins will have the sign $`(1)^{(N+1)/2}`$ while the down spins contribute a term with sign $`(1)^{(N1)/2}`$. In this case, to know the sign of the persistent current one must compare the amplitudes of the two contributions and the result will in general depend on the disorder configuration. However, in the zero-disorder case at low filling, the sign of the persistent current around flux $`\varphi =0`$ is known to be paramagnetic for odd $`N>1`$ loss-gold .
##### $`N`$ even, strong interaction
At strong interaction, the charges form a Wigner crystal which is pinned by the disorder. In this limit, the spin degree of freedom can be treated by an effective spin Hamiltonian. The latter turns out to be the anti-ferromagnetic Heisenberg Hamiltonian when $`\varphi =0`$, the positions $`i=i_1,i_2,\mathrm{},i_N`$ of the charges of the Wigner crystal being the spin lattice sites. Thus, the expectation values of the occupation number operators vanish everywhere except on these sites where they satisfy $`<\widehat{n}_{i_k,}+\widehat{n}_{i_k,}>=1`$. For an even number $`N`$ of spins in 1d, according to a theorem by Marshall, the ground state $`|\psi _0`$ of this spin Hamiltonian is a singlet of total spin $`S=0`$ auerbach , and can be expressed in the form
$$|\psi _0=\underset{\beta }{}f_\beta \underset{i=i_1}{\overset{i_N}{}}\left(c_{i,}^+\right)^{n_{i,}(\beta )}\underset{i=i_1}{\overset{i_N}{}}\left(c_{i,}^+\right)^{n_{i,}(\beta )}|0$$
(41)
with real $`f_\beta >0`$ for all spin configurations $`\beta `$ with fixed $`S_z=0`$. Note that in this expression, the ordering of the operators is done firstly according to spin, and secondly according to the position.
As in the case of spinless fermions without disorder, the ground state is a superposition of different basis states and the lowest order contributions to the ground state energy, which are flux dependent, are again of the order $`n=L_x`$. Similar to (36), they can arise from different kinds of sequences. First, only the up spins are moved around the ring, giving rise to a cyclic permutation of the spin up operators in (41), yielding the sign $`(1)^N_{}`$. The sequences involving only down spins give the same result. In addition, a sequence which moves all of the particles to the position of their neighbor can also contribute to the flux dependence since the resulting spin configuration is also contained in the ground state. The sign however is the same as the one of the previous sequences, since only one of the particles crosses the boundary and therefore only the order of the operators for one spin direction has to be restored by the corresponding cyclic permutation, and because the prefactor $`f_\beta ^{}f_\beta ^{}^{}`$, arising from the weights of the different components in the ground state, is always real and positive.
Since all contributions have the same sign, the sign of the persistent current in one-dimensional disordered chains at strong interaction is always given by this sign $`(1)^{N/2}`$, provided the particle number is even, independent of the particle density $`N/L`$. Because the ground state (41) holds only at $`\varphi =0`$, this sign rule is granted only in the vicinity of $`\varphi =0`$.
The resulting sign is the same as the one found in the non-interacting case loss-gold , and for the Hubbard model at half filling fye , but it differs from the result for the Hubbard model at low filling yu ; kotlyar , where the current around $`\varphi =0`$ is found to be diamagnetic for even numbers of particles at strong interaction. This difference may be due to the fact that the Hubbard interaction is local, making the Hubbard model at low filling in the $`U\mathrm{}`$ limit equivalent to spinless fermions. The more realistic long-range Coulomb interactions considered here do not show this artifact and lead to a qualitatively different behavior. The result is also different from the one obtained for Coulomb interacting electrons with spin in a continuous model with one barrier haeusler , which is always found to be diamagnetic at strong interaction. In this case, one may attribute the result to the fact that the interaction could lead to a rigid Wigner crystal which, in a continuous model, would be equivalent to one single heavy particle.
#### 4.4.2 Two-dimensional systems
For two dimensional systems, there is the same degeneracy in the position of the up and the down spins in the minimum $`S_z`$ subspace. In the strong interaction limit, each line of the Wigner crystal should follow the one dimensional law explained above, but the different spin configurations contributing to the ground state in 2d may have different numbers of spin up and spin down electrons in a given line. This yields contributions of different signs to the persistent current which have to be considered explicitly. It seems not to be possible to provide a simple sign rule. However, the power law decay of the persistent current at large interaction strength should have the same exponent as in the spinless case.
### 4.5 Size dependence and limitation of the theory
Here, we address the size dependence of the persistent current and the limitation of our theory which is given by the presence of defects in the Wigner crystal.
#### 4.5.1 Localization length for the Wigner crystal
From the exponential size-dependence of equation (22), one can extract the localization length for the Wigner crystal. A comparison to the exponential size scaling
$$I(L_x)\mathrm{exp}(L_x/\xi )$$
(42)
yields for large $`L_x`$ a localization length
$$\xi =\left(\mathrm{ln}\left(\frac{U}{t}\right)\right)^1,$$
(43)
showing the same decrease with increasing interaction strength as the localization length of the Mott insulator appearing in the Hubbard model at strong interaction stafford .
#### 4.5.2 Defects in the Wigner crystal
In the thermodynamic limit, the Wigner crystal is no longer a single domain and the above theory cannot be applied in the present form. At any finite, even very strong interaction, defects and domain walls can arise which allow to gain an amount of potential energy which grows with the size of the domains. The crystal might then prefer to be divided in defect-free domains each being pinned by the disorder. In the weak disorder limit, one can use similar arguments as in Imry in order to estimate the size of such domains. In the following, we briefly address different types of defects and calculate the associated critical ratio $`\left(\frac{U}{W}\right)_c`$ above which the theory applies.
##### Point defects
In the case of point-like defects, consisting of a single charge of the Wigner crystal being deplaced by one lattice constant $`a`$ with respect to its position in the perfect crystal, the gain in disorder potential energy can be estimated to be $`\mathrm{\Delta }E_{\mathrm{disorder}}^{\mathrm{point}}W`$. On the other hand, the cost in interaction energy is $`\mathrm{\Delta }E_{\mathrm{interaction}}^{\mathrm{point}}\frac{Ua^2}{b^3}`$ with $`b`$ being the mean spacing between particles. Since the charge density is given by $`\nu =1/b^2`$, this yields the density depending criterion
$$\frac{U}{W}>\left(\frac{U}{W}\right)_c^{\mathrm{point}}\frac{1}{a^2\nu ^{3/2}}$$
(44)
for the stability of the Wigner crystal against the creation of point defects.
##### Domain walls
The creation of a domain wall costs an interaction energy which is of the order $`\mathrm{\Delta }E_{\mathrm{interaction}}^{\mathrm{wall}}\frac{Ua^2}{b^3}L^{\mathrm{wall}}`$ with $`L^{\mathrm{wall}}`$ being the length of the wall. By deplacing a domain of linear size $`R`$ containing $`N_d=R^2\nu `$ particles, one can typically gain an amount $`W\sqrt{N_d}`$ of disorder potential energy. However, the first domain-like defect will appear in the most favorable position. Since the maximally possible gain is the extremely unprobable value $`WN_d`$, one can speculate that the gain in the optimal position is given by an intermediate power $`\mathrm{\Delta }E_{\mathrm{disorder}}^{\mathrm{domain}}WN_d^\gamma `$ with $`1/2\gamma <1`$. The length of the corresponding wall around the domain is $`L^{\mathrm{wall}}\sqrt{N_d}`$ such that this results in the stability of the crystal against the creation of domains of size $`R`$ for
$$\frac{U}{W}>\left(\frac{U}{W}\right)_c^{\mathrm{domain}}(R)\frac{R^{2\gamma 1}}{a^2\nu ^{2\gamma }}.$$
(45)
At a fixed W, if $`U>U_c^{\mathrm{domain}}(L)`$, a crystal of size $`L`$ will not be perturbed by domains. Otherwise, one can use $`U=U_c^{\mathrm{domain}}(R_c)`$ to extract the typical domain size $`R_c`$.
In systems larger than $`R_c`$, the electron crystal is divided in many domains of size $`R_c`$. If we neglect the coupling energy between these domains, a rough estimate for the response to a flux $`\varphi _x`$ in $`x`$-direction of a system of size $`(L_x,L_y)`$ is to take the product of the amplitudes of the longitudinally aligned $`L_x/R_c`$ domains and to sum over the responses of the $`L_y/R_c`$ “channels”.
Although we have not addressed all possible kinds of defects, we can assume that there is always a critical threshold $`\left(\frac{U}{W}\right)_c^{\mathrm{defect}}`$ above which our perturbative theory applies.
## 5 Summary
We have presented a systematic perturbative treatment of persistent currents in 2d lattice models for the case of strong Coulomb interaction, when the electronic charge density forms a Wigner crystal. The contribution with the weakest interaction dependence corresponds to sequences of one-particle hops along the shortest paths around the ring. These sequences dominate in the limit of strong interaction and determine the sign of the persistent current. For spinless fermions, this sign follows simple rules which depend only on the structure of the Wigner crystal.
Furthermore, we have shown that, except special cases, the leading term for the persistent current and therewith the sign of the persistent current at strong interaction do not depend on the realization of the disordered potential. Only the complete absence of disorder can qualitatively change the behavior. We considered the disorder corrections systematically and showed that they decay as $`W/U`$ at strong interaction.
In addition, we have shown that transverse currents appearing in individual disordered samples are suppressed much faster than the longitudinal current, thereby establishing an orientation of the local currents in longitudinal direction when the interaction is increased. This explains the numerical observation of the realization-independent sign of the persistent current and of the orientation of the local currents in longitudinal direction reported in Refs. benenti ; benenti\_mor ; berko1 ; berko2 .
###### Acknowledgements.
We thank G. Benenti, W. Häusler, G.-L. Ingold, D. Loss, J.-L. Pichard, P. Schwab, and X. Waintal for useful discussions. Particularly helpful discussions with R. Jalabert are gratefully acknowledged. Partial support from the TMR network “Phase coherent dynamics of hybrid nanostructures” of the EU and from the Procope program of the DAAD/A.P.A.P.E. is gratefully acknowledged.
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# Excitation of resonators by electron beams
## 1 Introduction
The excitation of resonators is described by Maxwell equations in vacuum \-
$$div\stackrel{}{E}=4\pi \rho (a)rot\stackrel{}{H}=\frac{4\pi }{c}\stackrel{}{J}+\frac{1}{c}\frac{\stackrel{}{E}}{t}(b),$$
$$rot\stackrel{}{E}=\frac{1}{c}\frac{\stackrel{}{H}}{t}(c),div\stackrel{}{H}=0(d).$$
(1)
These equations are a set of two vector and two scalar equations for vectors of electric $`\stackrel{}{E}(\stackrel{}{r},t)`$ and magnetic $`\stackrel{}{H}(\stackrel{}{r},t)`$ field strengths or eight equations for six independent components of the electric and magnetic fields. We suppose that the charge density $`\rho (\stackrel{}{r},t)`$ and current density $`\stackrel{}{J}(\stackrel{}{r},t)`$ are given values. It means that only four components of the electromagnetic field strengths are independent.
The solution of these equations includes transverse electromagnetic field strengths of free electromagnetic waves $`\stackrel{}{E}^{tr}`$, $`\stackrel{}{H}^{tr}`$ and accompanied longitudinal electric field strengths $`\stackrel{}{E}^l`$ of Coulomb fields of the beam crossing the resonator. Transverse electromagnetic field strengths excited by the beam in the resonator comply the condition $`div\stackrel{}{E}^{tr}=div\stackrel{}{H}^{tr}=div\stackrel{}{H}=0`$. Longitudinal electric field strength comply the condition $`rot\stackrel{}{E}^l=0`$, $`div\stackrel{}{E}^l=4\pi \rho `$ \- <sup>1</sup><sup>1</sup>1In general case transverse fields are not only free electromagnetic waves. Both a static magnetic field, a magnetic field accompanying a homogeneously moving particle and arbitrary time depended magnetic field are transverse one. A part of the Coulomb electrical field accompanying a relativistic particle is transverse one. The most simple example of the transverse electric field strength is the electric field strength of the homogeneously moving relativistic particle $`\stackrel{}{E}^{tr}=\stackrel{}{E}\stackrel{}{E}^l`$, where $`\stackrel{}{E}=e\stackrel{}{r}/\gamma ^2r^3`$, $`\stackrel{}{E}^l=e\stackrel{}{r}/r^3`$, $`\stackrel{}{r}`$ is the radius vector directed from the particle to the observation point, $`\gamma =\sqrt{1\beta ^2}`$ relativistic factor of the particle, $`R^{}=(xvt)^2+(1\beta ^2)(x^2+y^2)`$, $`\beta =v/c`$, $`v`$ the velocity of the particle , . After a particle beam cross a resonator then only transverse free electromagnetic waves stay at the resonator.. Free electromagnetic fields in resonators are solutions of homogeneous Maxwell equations ($`\stackrel{}{ȷ}=\rho =0`$) with corresponding boundary conditions. These solutions are a sum of eigenmodes of the resonator which include a discrete set of eigenfrequences $`\omega _\lambda `$ and corresponding to them functions $`\stackrel{}{E}_\lambda (\stackrel{}{r},t),\stackrel{}{H}_\lambda (\stackrel{}{r},t)`$ for the electric and magnetic field strengths (further we will omit the superscripts $`tr`$ and $`l`$ in the fields). The subscript $`\lambda `$ includes three numbers ($`m,n,q`$) corresponding to transverse and longitudinal directions of the resonator axis. In the case of open resonators the transverse electromagnetic $`TEM_{mnq}`$ modes are excited. When the number $`q`$ is very high then this number is omitted. Usually in the open resonators many longitudinal modes are excited even in the case of free-electron lasers emitting rather monochromatic radiation.
The solution of the problem of excitation of resonators is simplified by introduction of a transverse vector potential $`\stackrel{}{A}(\stackrel{}{r},t)=_\lambda \stackrel{}{A}_\lambda (\stackrel{}{r},t)`$ of free electromagnetic fields in Coulomb gauge $`div\stackrel{}{A}=0`$, where scalar potential $`\phi =0`$ when $`\rho =0`$ (here we omitted the superscripts $`tr`$ and $`l`$ in the vectors $`\stackrel{}{A}^{tr}`$). The corresponding wave equation for this vector can be solved by the method of separation of variables when we suppose $`\stackrel{}{A}_\lambda (\stackrel{}{r},t)=q_\lambda (t)\stackrel{}{A}_\lambda (\stackrel{}{r})`$, where $`q_\lambda (t)`$ is the amplitude of the vector potential and $`\stackrel{}{A}_\lambda (\stackrel{}{r})`$ is the eigenfunction of the resonator normalized by the condition $`|\stackrel{}{A}_\lambda (\stackrel{}{r})|^2𝑑V=1`$. In this case the total free electromagnetic field in the resonator is described by the expression $`\stackrel{}{A}(\stackrel{}{r},t)=_\lambda \stackrel{}{q}_\lambda (t)\stackrel{}{A}_\lambda (\stackrel{}{r})`$.
The electric and magnetic field strengths of the transverse free fields in resonators can be expressed through the vector potential in the form $`\stackrel{}{E}_\lambda (\stackrel{}{r},t)=d\stackrel{}{A}_\lambda (\stackrel{}{r},t)/dct=\dot{q}_\lambda (t)\stackrel{}{A}_\lambda (\stackrel{}{r})/c`$, $`\stackrel{}{H}_\lambda (\stackrel{}{r},t)=rot\stackrel{}{A}_\lambda (\stackrel{}{r},t)=q_\lambda (t)rot\stackrel{}{A}_\lambda (\stackrel{}{r})`$, where $`\dot{q}_\lambda (t)=dq_\lambda (t)/dt`$. When the charge and current densities are in the resonator then a scalar $`\phi _\sigma `$ and a longitudinal vector potential $`\stackrel{}{A}^l`$ ($`rot\stackrel{}{A}^l=0`$) determine Coulomb fields of the beam in the resonator. We are not interesting them in this paper.
When active and diffractive losses in the open resonator are absent then the vector potential of a free electromagnetic field in the resonator excited by the beam can be presented in the form
$$\stackrel{}{A}(\stackrel{}{r},t)=\underset{\lambda }{}q_{m\lambda }\stackrel{}{A}_\lambda (\stackrel{}{r})e^{i\omega _\lambda t},$$
(2)
where the coefficient $`q_{m\lambda }`$ is the amplitude of the excited eigenmode.
The electromagnetic fields excited by the electromagnetic beam are determined by the nonhomogeneous Maxwell equations or the corresponding equation for the vector potential
$$\mathrm{\Delta }\stackrel{}{A}(\stackrel{}{r},t)\frac{1}{c^2}\frac{^2\stackrel{}{A}(\stackrel{}{r},t)}{t^2}=\frac{4\pi }{c}\stackrel{}{J}(\stackrel{}{r},t).$$
(3)
The solution of the Eq(3) can be found in the form $`\stackrel{}{A}(\stackrel{}{r},t)=_\lambda q_\lambda (t)\stackrel{}{A}_\lambda ^{tr}(\stackrel{}{r})+_\sigma q_\sigma (t)\stackrel{}{A}_\sigma ^l(\stackrel{}{r})`$ (here we stayed superscripts $`tr`$ and $`l`$ and used the conditions $`\stackrel{}{A}_{\lambda ,\sigma }(\stackrel{}{r},t)=q_{\lambda ,\sigma }(t)\stackrel{}{A}_{\lambda ,\sigma }(\stackrel{}{r})`$). If we will substitute this expression into equation (3), integrate over the volume of the resonator, use the condition of normalization $`|\stackrel{}{A}_\lambda (\stackrel{}{r})|^2𝑑V=|\stackrel{}{A}_\sigma (\stackrel{}{r})|^2𝑑V=1`$, $`\stackrel{}{A}_\lambda (\stackrel{}{r})\stackrel{}{A}_\lambda ^{^{}}(\stackrel{}{r})𝑑V=\stackrel{}{A}_\lambda (\stackrel{}{r})\stackrel{}{A}_\sigma (\stackrel{}{r})𝑑V=\stackrel{}{A}_\sigma (\stackrel{}{r})\stackrel{}{A}_\sigma ^{^{}}(\stackrel{}{r})𝑑V=0`$ and take into account that the vector $`\stackrel{}{A}_\lambda `$ comply with the condition $`\mathrm{\Delta }\stackrel{}{A}_\lambda =(\omega _\lambda /c)^2\stackrel{}{A}_\lambda `$ then we will receive the equation for change of the amplitude of the eigenmode $`q_\lambda `$ for free fields in the resonator
$$\ddot{q}_\lambda +\omega _\lambda ^2q_\lambda =\frac{4\pi }{c}_V\stackrel{}{J}(\stackrel{}{r},t)\stackrel{}{A}_\lambda (\stackrel{}{r})𝑑V.$$
(4)
The expression (4) is the equation of the oscillator of unit mass excited by a force $`f(t)=(4\pi /c)_V\stackrel{}{J}(\stackrel{}{r},t)\stackrel{}{A}_\lambda (\stackrel{}{r})𝑑V`$. It describes the excitation of both enclosed and open resonators \- . The same expression for force determine the excitation of waveguides .
The eigenmodes of the rectangular resonators (cavities) were discovered by J.Jeans in 1905 when he studied the low of thermal emission. The equations (4) was used later for quantization of the electromagnetic field in quantum electrodynamics .
## 2 Emission of electromagnetic radiation by electron beams in open resonators
The equation (4) does not take into account the energy losses of the emitted radiation in the resonator. These losses can be introduced through the quality of the resonator $`Q_\lambda `$
$$\ddot{q}_\lambda +\frac{\omega _r}{Q_\lambda }\dot{q}_\lambda +\omega _\lambda ^2q_\lambda =\frac{4\pi }{c}_V\stackrel{}{J}(\stackrel{}{r},t)\stackrel{}{A}_\lambda (\stackrel{}{r})𝑑V,$$
(5)
where in the case of the open resonator $`\omega _r=2\pi /T`$, $`T=2L/c`$ is the period of oscillations of the light wavepacket between the resonator mirrors when it passes along the axis of the resonator (notice that in general case the frequencies $`\omega _\lambda =\omega _{mnq}`$ depend on $`m,n,q`$ and slightly differ from frequencies $`\omega _rq`$). Here we have introduced a version of a definition of a resonator quality connected with the frequency $`\omega _r`$. Another version of a quality is usually connected with the frequency $`\omega _\lambda `$. Our definition is more convenient for the case of free-electron lasers using open resonators.
Using (5) we can derive the expression for the energy balance in the resonator. For this purpose we can multiply this equation by $`\dot{q}_\lambda `$ and integrate over the volume of the resonator. Then we receive the equation
$$\frac{1}{2}\frac{d}{dt}[\dot{q}_\lambda ^2+\omega _\lambda ^2q_\lambda ^2]+(\frac{\omega _r}{Q_\lambda })^2\dot{q}_\lambda ^2=4\pi _V\stackrel{}{J}(\stackrel{}{r},t)\stackrel{}{E}_\lambda (\stackrel{}{r},t)𝑑V.$$
(6)
If we take into account that $`\stackrel{}{E}_\lambda (\stackrel{}{r},t)=\dot{q}_\lambda (t)\stackrel{}{A}_\lambda (\stackrel{}{r})/c`$, $`\stackrel{}{H}_\lambda (\stackrel{}{r},t)=rot\stackrel{}{A}_\lambda (\stackrel{}{r})`$, $`rot\stackrel{}{A}_\lambda =\omega _\lambda \stackrel{}{A}_\lambda /c`$, $`|\stackrel{}{A}_\lambda (\stackrel{}{r})|^2𝑑V=1`$ then the energy of the free electromagnetic field in the resonator can be presented in the form $`\epsilon _\lambda ^{em}=[(|\stackrel{}{E}_\lambda |^2+|\stackrel{}{H}_\lambda |^2)/8\pi ]`$ $`dV=[\dot{q}_\lambda ^2+\omega _\lambda ^2q_\lambda ^2]/8\pi c^2`$ and the equation (7) can be presented in the another form
$$\dot{\epsilon }_\lambda ^{em}+(\omega _r/Q_\lambda )\epsilon _\lambda ^{em}=_V\stackrel{}{J}(\stackrel{}{r},t)\stackrel{}{E}_\lambda (\stackrel{}{r},t)𝑑V.$$
(7)
The equation (5) is the pendulum equation with a friction. It determine the time evolution of the electromagnetic field stored at the resonator, when the time dependence of the beam current $`\stackrel{}{J}(\stackrel{}{r},t)`$ is given. The amplitude $`q_\lambda (\stackrel{}{r},t)`$ according to (5) is determined by the coefficient of expansion of the given current into series of eigenfunctions of the resonator. Notice that the value $`\stackrel{}{A}_\lambda [\stackrel{}{r}_e(t)]`$ depends on $`t`$ only through $`\stackrel{}{r}_e(t)`$ and the value $`\stackrel{}{E}_\lambda [t,\stackrel{}{r}_e(t)]=\dot{q}_\lambda (t)\stackrel{}{A}_\lambda [\stackrel{}{r}_e(t)]/c`$ depends on $`t`$ directly through $`q_\lambda (t)`$ and through $`\stackrel{}{r}_e(t)`$.
In the case of one particle of a charge ”e” the beam current density $`\stackrel{}{J}(\stackrel{}{r},t)=e\stackrel{}{v}(t)\delta [\stackrel{}{r}\stackrel{}{r}_e(t)]`$. In this case the force $`f(t)=e\stackrel{}{v}[\stackrel{}{r}_e(t)]\stackrel{}{A}_\lambda [\stackrel{}{r}_e(t)]`$ and the power transferred from the electron beam to the resonator wave mode $`\lambda `$ excited in the resonator $`P_\lambda (t)=_V\stackrel{}{J}(\stackrel{}{r},t)\stackrel{}{E}_\lambda (\stackrel{}{r},t)𝑑V=e_i\stackrel{}{v}_{ei}(t)\stackrel{}{E}_\lambda [(\stackrel{}{r}_{e,ı}(t),t)]`$. Using these expressions of force and power for all electrons ”i” of the beam we can present the equations (5), (7) in the form
$$\ddot{q}_\lambda +\omega _\lambda ^2q_\lambda =\frac{4\pi e}{c}\underset{i}{}\stackrel{}{v}_{ei}(t)\stackrel{}{A}_\lambda [\stackrel{}{r}_{ei}(t)],$$
(8)
$$\dot{\epsilon }_\lambda ^{em}+(\omega _r/Q_\lambda )\epsilon _\lambda ^{em}=e\underset{i}{}\stackrel{}{v}_{ei}(t)\stackrel{}{E}_\lambda [(\stackrel{}{r}_{ei}(t),t)].$$
(9)
It follows from (5), (7) and (8), (9) that transverse resonator modes are excited only in the case when the force $`f(t)0`$ and the power $`P_\lambda (t)0`$ that is when the particle trajectory passes through the regions where the corresponding resonator modes have large intensities and when the particle velocity has transverse and/or longitudinal components directed along the direction of the electric field strength. Open resonators on the level with enclosed ones have modes with longitudinal components of electric field strength (see Appendix). It means that open resonators can be excited even in the case when the particle trajectory have no transverse components and its velocity is directed along the axis of the resonator<sup>2</sup><sup>2</sup>2In this case the transition radiation is emitted by particles when they pass the walls of a resonator. The electromagnetic radiation will be emitted in the form of thin spherical layers at the first and second resonator mirrors . It will be reflected then repeatedly by resonator mirrors. The expansion of the electromagnetic fields of the spherical layers will be described by the series (2).. Using external fields of a single bending magnet can increase the power of the generated radiation. Both in the case of lack of a banding magnet and presence of one bending magnet the broadband radiation is emitted. The experiment confirms this observations . Using external fields of undulators and beams bunched at frequencies of the emitted radiation can lead to emission of rather monochromatic radiation.
In the simplest case when the beam current density $`\stackrel{}{J}(\stackrel{}{r},t)`$ is a periodic function of time then the force can be expanded in the series $`f(t)=_V\stackrel{}{J}(\stackrel{}{r},t)\stackrel{}{A}_\lambda (\stackrel{}{r})𝑑V=_{\nu =\mathrm{}}^{\mathrm{}}f_{\lambda \nu }\mathrm{exp}[i(\nu \omega _bt\phi _{\lambda \nu })]`$, where $`\omega _b=2\pi /T_b`$ and $`T_b`$ are a period and frequency of the current density oscillation accordingly, $`f_{\lambda \nu }=(1/T_b)_{T_b/2}^{T_b/2}f(t)\mathrm{exp}(i\nu \omega _bt)𝑑t`$, are the known coefficients, $`\phi _{\lambda \nu }`$ phase. The value $`f_{\lambda \nu }=f_{\lambda \nu }^{}`$, where $`f_{\lambda \nu }^{}`$ is the complex conjugate of $`f_{\lambda \nu }`$. The solution of the equation (5) for the case of the established oscillations ($`tQ_\lambda T_b`$) is
$$q_\lambda (t)=\underset{\nu =1}{\overset{\mathrm{}}{}}A_{\lambda \nu }\mathrm{exp}[i(\nu \omega _bt\theta _{\lambda \nu })],$$
(10)
where
$$A_{\lambda \nu }=\frac{f_{\lambda \nu }}{\sqrt{(\omega _\lambda ^2\nu ^2\omega _b^2)^2+(\nu \omega _r\omega _b/Q_\lambda )^2}},$$
$$\theta _{\lambda \nu }=\phi _{\lambda \nu }+arctg\frac{\nu \omega _r\omega _b}{Q_\lambda (\omega _\lambda ^2\nu ^2\omega _b^2)}.$$
It follows from the equation (10) that the maximum of the amplitude of the vector potential $`A_{\lambda \nu }=Q_\lambda f_{\lambda \nu }/\omega _r\omega _\lambda ^2`$ takes place at resonance $`\nu \omega _b=\omega _\lambda =\omega _{mnq}\omega _rq`$. Notice that all modes $`\lambda `$ are excited at the same frequency $`\omega _b`$ of the oscillator. In general case $`\omega _b\omega _\lambda `$.
The equation (10) is the first order linear equation of the energy change in the resonator excited by the electron beam. It follows from this equation that after switching off the beam current at some moment $`t_0`$ ($`\stackrel{}{J}(\stackrel{}{r},t)|_{t>t_0}=0`$) the energy in the resonator will be changed by the law $`\epsilon _\lambda ^{em}=\epsilon _{\lambda ,/0}^{em}\mathrm{exp}[(tt_0)/\tau ]`$, where $`\tau =Q_\lambda /\omega _r`$, $`\epsilon _{\lambda \mathrm{\hspace{0.17em}0}}^{em}=\epsilon _\lambda ^{em}|_{t=t_0}`$. On the contrary after switching on the beam current at some moment $`t_0`$ the energy in the resonator will be changed by the law $`\epsilon _\lambda ^{em}=\epsilon _{\lambda m}^{em}(1\mathrm{exp}[(tt_0))/\tau ]`$, where the energy of the electromagnetic field in the resonator $`\epsilon _{\lambda m}^{em}`$ is determined by the parameters of the resonator and the beam.
The considered example describes the emission of an oscillator or a system of oscillators which are in phase with the excited mode and have zero average velocity (trajectory has a form $`r_e=r_{e0}+\stackrel{}{ı}a_0\mathrm{cos}\omega _0t`$, where $`\stackrel{}{ı}`$ is the unite vector directed along the axis $`x`$). More complicated examples of trajectories of particles using for excitation of resonators by electron beams can be considered (the arc of circle, sine- or helical-like trajectories in bending magnets and undulators).
## 3 Vector TEM modes of open resonators
The theory of high quality open resonators does not differ from enclosed ones. But eigenmodes of open resonators have some unique features. The spectrum of the open resonators is rarefied, the operating mode spectrum has maximum selectivity. The dimensions of open resonators are much higher then the excited wavelengths and the dimensions of the enclosed resonators are of the order of excited wavelengths. The quality of open resonators at the same wavelengths is higher then enclosed ones.
There are some methods of calculation of TEM modes in open resonators. Usually scalar wave equations are investigated , . There is a small information in technical publications about distribution of vectors of the electric and magnetic field strengths in such resonators. In this section we search some distributions. In the Appendix the foundations of the excitation of resonators by electron beams are presented .
We will present the result for the Cartesian coordinates. In this case the solution of the scalar wave equation (24) (see Appendix) has a form
$$V_{mn}(x,y,z)=\frac{C}{\sqrt{w_x(z)w_y(z)}}H_m\left(\frac{\sqrt{2}x}{w_x(z)}\right)H_n\left(\frac{\sqrt{2}y}{w_y(z)}\right)$$
$$\mathrm{exp}\left\{\frac{ik}{2}\left(\frac{x^2}{q_x(z)}+\frac{y^2}{q_y(z)}\right)i(m+\frac{1}{2})arctg\frac{\lambda z}{\pi w_{0x}^2}i(n+\frac{1}{2})arctg\frac{\lambda z}{\pi w_{0y}^2}\right\}$$
(11)
and for the cylindrical coordinates
$$V(r,\varphi ,z)=C\left(\frac{r}{w(z)}\right)^m\left(\genfrac{}{}{0pt}{}{\mathrm{sin}m\varphi }{\mathrm{cos}m\varphi }\right)L_n^m\left(\frac{2r^2}{w^2(z)}\right)\mathrm{exp}\left\{\frac{ikr^2}{q(z)}i(m+2n+1)arctg\frac{\lambda z}{\pi w_0^2}\right\}w(z)^1,$$
(12)
where $`H_m`$, $`H_n`$ are the Hermittian polynomials, $`L_n^m`$ the Lagerian polynomials, $`\lambda =2\pi c/\omega `$ is the wavelength, C = constant,
$$\frac{1}{q(z)}=\frac{1}{R(z)}+\frac{i\lambda }{\pi w^2(z)},R(z)=z\left[1+\left(\frac{\pi w_0^2}{\lambda z}\right)^2\right],w^2(z)=w_0^2\left[1+\left(\frac{\lambda z}{\pi w_0^2}\right)^2\right].$$
In (11), (12) $`R(z)`$ is the radius of the wave front of Gaussian beam, $`w(z)`$ the radius of the beam, $`w_0(z)`$ the radius of the waist of the beam.
At $`m=n=0`$ we have the main mode of the Gaussian beam. If $`w_{0x}=w_{0y}=w_0`$ then the main modes for the Cartesian and cylindrical coordinates are the same
$$U(x,y,z)=\frac{C}{w(z)}\mathrm{exp}\left\{\frac{x^2+y^2}{w^2(z)}+\frac{ik}{2}\frac{x^2+y^2}{R(z)}iarctg\frac{\lambda z}{\pi w_0^2}\right\}\mathrm{exp}^{i(kz\omega t)}.$$
(13)
We have the solutions (11), (12) of the scalar wave equation (24) for the space limited beam. Now we can find vectors of the electric and magnetic field strengths using the expressions (23) and possible ways of construction of Hertz vectors. Let us suppose the next compositions with the electric Hertz vector assuming that magnetic Hertz vector is zero:
1) $`\mathrm{\Pi }_x^e=U(x,y,z)`$, $`\mathrm{\Pi }_y^e=\mathrm{\Pi }_z^e=0`$.
2) $`\mathrm{\Pi }_x^e=0`$, $`\mathrm{\Pi }_y^e=U(x,y,z)`$, $`\mathrm{\Pi }_z^e=0`$.
3) $`\mathrm{\Pi }_x^e=0`$, $`\mathrm{\Pi }_y^e=0`$, $`\mathrm{\Pi }_z^e=U(x,y,z)`$.
In the first case
$`div\stackrel{}{\mathrm{\Pi }}=\mathrm{\Pi }_x/x=V/x\mathrm{exp}[i(kz\omega t)]`$, $`(rot\stackrel{}{\mathrm{\Pi }})_x=0`$,
$`(rot\stackrel{}{\mathrm{\Pi }})_y=(V/z+ikV)\mathrm{exp}[i(kz\omega t)]`$, $`(rot\stackrel{}{\mathrm{\Pi }})_z=(V/y)\mathrm{exp}[i(kz\omega t)]`$
and
$`E_x^1=^2V/x^2+k^2V`$, $`E_y^1=^2V/xy`$,
$`E_z^1=^2V/xz+ikV/x`$, $`H_x^1=0`$, $`H_y^1=ikV/zk^2V`$, $`H_z^1=ikV/y`$.
The upper superscript shows the first composition of the Hertz vector. A common multiple $`\mathrm{exp}[i(kz\omega t)]`$ for all field components is omitted.
The values $`^2V/x_ix_kkV/x_ik^2V`$. That is why in this case $`E_x^1E_y^1,E_z^1`$, $`H_y^1H_z^1`$.
The second case does not differ from the first one. It is necessary to substitute variable $`x`$ by $`y`$ and vise versa.
In the third case
$`div\stackrel{}{\mathrm{\Pi }}=\mathrm{\Pi }_x/z=V/x\mathrm{exp}[i(kz\omega t)]`$, $`(rot\stackrel{}{\mathrm{\Pi }})_x=V/y\mathrm{exp}[i(kz\omega t)]`$,
$`(rot\stackrel{}{\mathrm{\Pi }})_y=V/x\mathrm{exp}[i(kz\omega t)]`$, $`(rot\stackrel{}{\mathrm{\Pi }})_z=0`$
and
$`E_x^3=^2V/xz+ikV/x`$, $`E_y^3=^2V/zy+V/y`$,
$`E_z^3=2ikV/z`$, $`H_x^3=ikV/y`$, $`H_y^3=ikV/x`$, $`H_z^3=0`$.
It follows that in the case of the main mode the electric and magnetic field strengths corresponding to the electric Hertz vector have components:
$$E_x^1=k^2U(x,y,z),E_y^10,E_z^1=2ikx\left[\frac{1}{w^2(z)}+\frac{ik}{R(z)}\right]U(x,y,z),$$
$$H_x^10,H_y^1=k^2U(x,y,z),H_z^1=2iky\left[\frac{1}{w^2(z)}+\frac{ik}{R(z)}\right]U(x,y,z)$$
$$E_x^2=0,E_y^2=k^2U(x,y,z),E_z^2=2iky[\frac{1}{w^2(z)}+\frac{ik}{R(z)}]U(x,y,z),$$
$$H_x^2=k^2U(x,y,z),H_y^20,H_z^2=2ikx\left[\frac{1}{w^2(z)}+\frac{ik}{R(z)}\right]U(x,y,z)$$
$$E_x^3=2ikx\left[\frac{1}{w^2(z)}+\frac{ik}{R(z)}\right]U(x,y,z),E_y^3=2iky\left[\frac{1}{w^2(z)}+\frac{ik}{R(z)}\right]U(x,y,z),$$
(14)
$$E_z^3=2ik\left\{\frac{4\lambda (x^2+y^2)z}{(\pi w_0)^2w^3(z)}+\frac{ik(x^2+y^2)}{2R^2(z)}\left[1\left(\frac{\pi w_0^2}{\lambda z}\right)^2\right]\frac{i\lambda }{\pi w_0^2\left[1+\left(\frac{\lambda z}{\pi w_0^2}\right)^2\right]}\right\}U(x,y,z),$$
$$H_x^3=2iky\left[\frac{1}{w^2(z)}+\frac{ik}{R(z)}\right]U(x,y,z),H_y^3=2iky\left[\frac{1}{w^2(z)}+\frac{ik}{R(z)}\right]U(x,y,z),H_z^3=0.$$
The electric and magnetic field strengths received from magnetic Hertz vector can be received from the fields (14) as well. For this purpose we can take the vector of the electric field strength received from magnetic Hertz vector equal to the negative value of the magnetic field strength received from the electric Hertz vector $`\stackrel{}{E}^{^{}}\stackrel{}{H}`$ and by analogy we can take $`\stackrel{}{H}^{^{}}\stackrel{}{E}`$.
The general solution for the electromagnetic field strength of the main mode of Gaussian beam $`TEM_{00}`$ can be presented in the form
$$\stackrel{}{E}=c_1\stackrel{}{E}^1+c_2\stackrel{}{E}^2+c_3\stackrel{}{E}^3c_4\stackrel{}{H}^1c_5\stackrel{}{H}^2c_6\stackrel{}{H}^3,$$
$$\stackrel{}{H}=c_1\stackrel{}{H}^1+c_2\stackrel{}{H}^2+c_3\stackrel{}{H}^3+c_4\stackrel{}{E}^1+c_5\stackrel{}{E}^2+c_6\stackrel{}{E}^3,$$
(15)
where $`c_i`$ are the arbitrary coefficients determined by the conditions of excitation of the mode by the electron beam. Waves determined by the only coefficient $`c_i`$ (when another ones are equal to zero) can be excited independently.
Higher modes in the open resonator will be described by the expressions (15) and by the expressions similar to (14) for the electromagnetic field strengths of the main mode. They will form orthogonal and full set of fundamental waves. The arbitrary wave may be expanded into these waves. Of cause, real electric and magnetic field strengths are determined by the real part of the expression (15).
In the open resonators the same Gaussian beams are excited. They propagate between mirrors both in $`z`$ and in $`z`$ directions. However the resonators will be excited on discrete set of eigenfrequences (wavelengths) .
We can see that according to (14) all considered waves $`\stackrel{}{E}^i,\stackrel{}{H}^i`$ are transverse. At the same time they have longitudinal components. This is the general property of the convergent and divergent waves , . Such waves have longitudinal components which permit the lines of the electric and magnetic field strengths to be closed.
The fields $`\stackrel{}{E}^1,\stackrel{}{H}^1`$ describe an electromagnetic wave with one direction of polarization and the fields $`\stackrel{}{E}^2,\stackrel{}{H}^2`$ with another one. They have high transverse components of the electric and magnetic field strengths and zero longitudinal components on the axis $`z`$.
Electromagnetic fields $`\stackrel{}{E}^3,\stackrel{}{H}^3`$ are a new kind of fields. They have zero transverse components of the electric and magnetic field strengths and high value longitudinal component of the electric field strength at the axis $`z`$ (similar to the wave $`E_{01}`$ at the axial region of the cylindrical waveguide). It means that in this case the lines of the electric and magnetic field strengths are closed in the directions both at the central part of the beam propagation that is near to the axis $`z`$ and far from the axis that is near to the region of theirs envelopes (caustics)<sup>3</sup><sup>3</sup>3Notice that usually the divergent waves with high directivity emitted by antennas are described and drawn by the lines of the electric and magnetic field strengths which are closed in the directions far from the axis of the beam propagation near to the region of theirs envelops..
Usually the scalar functions $`V(x,y,z)`$ or $`U(x,y,z)=V(x,y,z)\mathrm{exp}[i(kz\omega t)]`$ are used when the modes in open resonators are investigated , , . It was supposed that the waves are transverse ones and the values of the electromagnetic field strengths are distributed near the same way as the values of the scalar functions. At that some features like the existence of the wave $`\stackrel{}{E}^3,\stackrel{}{H}^3`$ were hidden. Such waves have longitudinal components of the electric field strength and hence can be excited through the transition radiation emitted on the inner sides of the resonator walls by an electron homogeneously moving along the axis $`z`$. Such excitation was observed in the experiments published in .
## 4 Conclusion
Open resonators permit an effective generation of broadband radiation at the main and/or other transverse modes under conditions when many longitudinal modes are excited. The longitudinal modes are limited in the longwavelength region by the diffraction losses and in the short wavelength region by the longitudinal electron beam dimensions (coherence conditions). Open resonators can be excited in the case when the external fields in the resonator are absent and the particle trajectory is directed along the axis of the resonator. Using external fields of a single bending magnet can increase the power of the generated radiation .
Appendix
Generation and propagation of electromagnetic waves in vacuum is described by Maxwell equations (1). We noticed above that these equations are a set of eight equations for six independent components of the electric and magnetic fields. Only four components of the electromagnetic field are independent. These equations added with initial and boundary conditions describe all processes in electrodynamics.
There is no general solution of the system of Maxwell equations with boundary conditions similar to the Lienard-Viechert solution for the fields produced by charged particles moving along some trajectories at a given low in free space. It means that private problems must be solved separately for every concrete case. At that when the boundary conditions exist, interactions of particles with surrounding media and intrabeam interactions of particles are essential then the beam density and beam current can not be given and the dynamical Lorentz equations must be added. Below we will consider the case when the beam density and the density of the beam current (particle trajectories) are given.
One of the possible simplifications of the solution of the Maxwell equations is to transform linear Maxwell equations to the equations of the second order relative to the field strengths or potentials.
First of all the Maxwell equations can be transformed to the equations separately for the electric and magnetic fields. For this purpose we can differentiate equation (1.b) with respect to $`t`$, use equation (1.c) and employ the vector identity $`rotrot\stackrel{}{F}=graddiv\stackrel{}{F}\mathrm{\Delta }\stackrel{}{F}`$, where $`\mathrm{\Delta }`$ is the Laplacian operator. Such a way we will receive the equation for the electric field strength and then by analogy we will receive the equation for the magnetic field strength. They are
$$\mathrm{}\stackrel{}{E}=\frac{4\pi }{c^2}\dot{\stackrel{}{J}}+4\pi grad\rho ,(a)\mathrm{}\stackrel{}{H}=\frac{4\pi }{c}rot\stackrel{}{J}.(b)$$
(16)
where $`\mathrm{}=\mathrm{\Delta }^2/c^2t^2`$ is the d’Alembertian operator, $`\dot{\stackrel{}{J}}=\stackrel{}{J}/t`$.
The equations (16) are the nonhomogeneous linear equations of the second order. We must add the equations (1.a), (1.d) to the system of the equations (16). It means that we have again a system of two vector and two scalar equations (in components they are eight equations) for six unknown components of the electric and magnetic field strengths $`E_i`$, $`H_i`$.
The divergence of the equation (16.a) leads to a more general continuity equation $`(/t)(\rho /`$ $`t`$ $`+div\stackrel{}{J})=0`$ which is valid when the continuity equation $`\rho /t+div\stackrel{}{J}=0`$ is valid.
The solution of the Maxwell equations will be the solution of these second order equations. The second order equations are another equations. Strictly speaking they are not equivalent to Maxwell equations. We must check theirs solutions by substituting these solutions into the linear Maxwell equations to reject unnecessary solutions. This is very difficult problem even for simple cases. A way out can be found by introducing of electromagnetic field potentials. The vector potential $`\stackrel{}{A}`$ and scalar potential $`\phi `$ are introduced by the equations $`\stackrel{}{H}=rot\stackrel{}{A}`$, $`\stackrel{}{E}=grad\phi (1/c)(\stackrel{}{A}/t)`$. In this case both from Maxwell equations and from the equations (16) it follows the equations for vector and scalar potentials
$$\mathrm{}\stackrel{}{A}=\frac{4\pi }{c}\stackrel{}{J}(a),\mathrm{}\phi =4\pi \rho (b)$$
(17)
and additional condition coupling the potentials (Lorentz gauge)
$$div\stackrel{}{A}=\frac{1}{c}\frac{\phi }{t}.$$
(18)
It is convenient to use the electric and magnetic Hertz vectors as well. They permit to simplify the solutions of the problem of propagation of waves in resonators and free space which is described by the homogeneous wave equations ($`\rho =0`$, $`\stackrel{}{J}=0`$). Both the electric and magnetic Hertz vectors $`\stackrel{}{\mathrm{\Pi }}^e`$, $`\stackrel{}{\mathrm{\Pi }}^m`$ are introduced by the same expressions
$$\stackrel{}{A}=\frac{1}{c}\frac{\stackrel{}{\mathrm{\Pi }}^{e/m}}{t};\phi =div\stackrel{}{\mathrm{\Pi }}^{e/m}.$$
(19)
Such way defined potentials $`\stackrel{}{A}`$ and $`\phi `$ will satisfy the equation (11) simultaneously.
Different superscripts $`e/m`$ in this case are used on the stage of introduction of the connection between electric and magnetic field strengths through Hertz vectors. The electric field strength can be expressed through the electric and magnetic Hertz vector by the equations
$$\stackrel{}{E}=graddiv\stackrel{}{\mathrm{\Pi }}^e\frac{1}{c^2}\frac{^2\stackrel{}{\mathrm{\Pi }}^e}{t^2};\stackrel{}{H}=\frac{1}{c}\frac{}{t}rot\stackrel{}{\mathrm{\Pi }}^e,$$
(20)
$$\stackrel{}{E}=\frac{1}{c}\frac{}{t}rot\stackrel{}{\mathrm{\Pi }}^m,\stackrel{}{H}=graddiv\stackrel{}{\mathrm{\Pi }}^m\frac{1}{c^2}\frac{^2\stackrel{}{\mathrm{\Pi }}^m}{t^2}.$$
(21)
These manipulations are valid because of both definitions (20) and (21) satisfy Maxwell equations (1) and equations (16). This is because of homogeneous wave equations for electromagnetic fields
$$\mathrm{}\stackrel{}{F}=0$$
(22)
are symmetric relative to fields $`\stackrel{}{E}`$, $`\stackrel{}{H}`$ ($`\stackrel{}{F}=\stackrel{}{E},\stackrel{}{H}`$). If $`\stackrel{}{E}`$ and $`\stackrel{}{H}`$ are some solutions of the homogeneous Maxwell equations (1b), (1c) then vectors $`\stackrel{}{E}^{^{}}=\stackrel{}{H}`$ and $`\stackrel{}{H}^{^{}}=\stackrel{}{E}`$ will satisfy these and another Maxwell equations as well.
In general case the problem may be reduced to solving of wave equation if potentials $`\stackrel{}{\mathrm{\Pi }}^e`$, $`\stackrel{}{\mathrm{\Pi }}^m`$ will be introduced simultaneously in the form
$$\stackrel{}{E}=graddiv\mathrm{\Pi }^e\frac{1}{c^2}\frac{^2\stackrel{}{\mathrm{\Pi }}^e}{t^2}\frac{1}{c}\frac{}{t}rot\stackrel{}{\mathrm{\Pi }}^m,\stackrel{}{H}=\frac{1}{c}\frac{}{t}rot\stackrel{}{\mathrm{\Pi }}^e+graddiv\stackrel{}{\mathrm{\Pi }}^m\frac{1}{c^2}\frac{^2\stackrel{}{\mathrm{\Pi }}^m}{t^2}.$$
(23)
We can be convinced that $`\stackrel{}{E}`$ and $`\stackrel{}{H}`$ described by (23) fulfil to Maxwell equations at $`\rho =\stackrel{}{J}=0`$ when vectors $`\stackrel{}{\mathrm{\Pi }}^e`$ and $`\stackrel{}{\mathrm{\Pi }}^m`$ fulfil the wave equation (22) with replaced $`\stackrel{}{F}`$ on $`\stackrel{}{\mathrm{\Pi }}^e`$ and $`\stackrel{}{\mathrm{\Pi }}^m`$.
Equation (22) is valid for each component of vectors $`\stackrel{}{\mathrm{\Pi }}^e`$ and $`\stackrel{}{\mathrm{\Pi }}^m`$. That is why it is possible to use scalar wave equation
$$\mathrm{}U=0$$
(24)
and identify its solution $`U`$ with one of components of vectors $`\stackrel{}{\mathrm{\Pi }}^e`$ or $`\stackrel{}{\mathrm{\Pi }}^m`$ and the rest components of these vectors equate to zero (say we can take $`\stackrel{}{\mathrm{\Pi }}^e=\stackrel{}{e}_x0+\stackrel{}{e}_y0+\stackrel{}{e}_zU,`$ $`\stackrel{}{\mathrm{\Pi }}^m=0`$). Substituting the constructed such a way vector with one component in (16) we will find the electromagnetic field strengths $`\stackrel{}{E}`$, $`\stackrel{}{H}`$ which satisfy the Maxwell and wave equations. Then we can identify the same solution with another component of the Hertz vector, equate the rest components to zero and calculate another electromagnetic field strengths $`\stackrel{}{E}`$, $`\stackrel{}{H}`$ which satisfy the Maxwell and wave equations as well. After we will go through all compositions with components then we will have a set of six different solutions for field strengths $`\stackrel{}{E}`$, $`\stackrel{}{H}`$. These solutions will be six electromagnetic waves with different structures. Sum of these solutions with some coefficients will be a solution of the Maxwell equations as well. This will be algorithm of electromagnetic field determination through Hertz vector.
Equation (24) has many different solutions. We must find such solutions which will correspond to the problem under consideration to a considerable extent. Below we will deal with monochromatic light beams of the limited diameter related with resonator modes. In general case such beams can be written in the form
$$U(x,y,z)=V(x,y,z)e^{i(kz\omega t)}$$
(25)
where $`V(x,y,z)`$ is a function of coordinate slowly varying in comparison with $`\mathrm{exp}i(kz\omega t)`$. A complex form of values will be used for computations and then we will proceed to a real part of the form.
Substituting (25) into (24) and taking into account the slow variation of $`V(x,y,z)`$ compared with $`\mathrm{exp}i(kz\omega t)`$ that is the condition $`|^2V/z^2|2k|V/z|`$ and the condition $`k=\omega /c`$ we will receive the equation
$$i\frac{V}{z}+\frac{1}{2k}(\frac{^2V}{x^2}+\frac{^2V}{y^2})=0$$
(26)
which describes a space limited beam.
In the general case the limited in the transverse direction wave propagating in free space or in a resonator have rather complicated structure. That is why it is desirable to find full, orthogonal set of fundamental waves with the well known feature of propagation. Then an arbitrary wave may be expanded into series of these waves. Different series of fundamental waves can be found for this problem and the arbitrary wave can be expanded into one or another series. The method of separation of variables is used to solve the wave equation. For example, in the Cartesian coordinates $`V(x,y,z)=X(x,y,z)Y(x,y,z)`$ and in the cylindrical coordinates $`V(x,y,z)=G(u)\mathrm{\Phi }(\phi )\mathrm{exp}[ikr^2/2q(z)]\mathrm{exp}[iS(z)]`$, where $`r`$ and $`\phi `$ are cylindrical coordinates on a plane transverse to $`z`$, $`u=r/w(z)`$. These solutions are considered in .
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# 1 Introduction
## 1 Introduction
According to idea clearly expressed in ref. our observed world could be a brane embedded in a higher-dimensional space. Moreover, in the scenario four-dimensional gravity with acceptable properties may be recovered. This fact initiated enormous activity in the study of different aspects of brane worlds.
In some versions of brane-world scenarios one can construct the black holes on the branes . The investigation of their properties is extremely interesting because this may help to understand better the fundamental problems of quantum gravity. For example, one can try to describe from brane point of view such phenomena as Hawking radiation, Black Hole entropy origin, etc. It is most likely that brane-world scenario should be realized within the context of AdS/CFT correspondence (say, in its simplest form as 5d gauged SG/4d CFT). If it is so one should start from the scalar-tensor gravity as a bulk theory. Then the role of scalars (dilaton if only single scalar presents) should be carefully addressed. In the present paper working in this direction we construct the family of dilatonic brane-world black holes (including regular cases ,like de Sitter space) and carefully investigate their properties, the problem of localization of 4d gravity in such spaces and brane corrections to Newton constant.
In the next section we start from d+1-dimensional dilatonic gravity with d-dimensional brane vacuum energy and formulate the corresponding equations of motion. A flat brane solution with very small cosmological constant is possible. A family of brane-world black holes solutions (curved branes) is presented. They correspond to (anti) de Sitter space, Nariai space, Kerr, and Schwarzschild-de Sitter black holes,etc. Some their properties are investigated. In section 3 we look to gravity perturbations around such backgrounds. Not only graviton but also dilaton perturbations are found. It is explicitly shown that in some cases 4d gravity may be localized in the same fashion as in ref.. Section four is devoted to the study of corrections to Newton constant. The Newton potential is calculated and it is shown that corrections to Newton law near branes are very small. In section 5 we analyse the regular solution where de Sitter Universe is realised on the brane. Brane vacuum energy is found. Simple analysis indicates that there may be problems with gravity localization. The role of quantum effects of brane matter is investigated in section 6. We suggest to add the conformal anomaly induced effective action of boundary, dilaton coupled matter to the complete action. In such a way, Randall-Sundrum compactification may be naturally fitted with AdS/CFT correspondence.It is shown that with account of such effective action there still exists 4d de Sitter wall (our Universe) living in 5d dilatonic spacetime which is asymptotically AdS. Quantum corrections change the brane vacuum energies, they become explicitly time-dependent. RG flow of Newton constant in IR and UV is briefly discussed. Some resume is given in last section.
## 2 Dilatonic black hole solutions in the brane world
In analogy with Randall-Sundrum model , we start with the following action of the gravity coupled with dilaton $`\varphi `$:
$`S`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G_{d+1}}}[{\displaystyle _M}d^{d+1}x\sqrt{g}(R{\displaystyle \frac{1}{2}}_\mu \varphi ^\mu \varphi V(\varphi ))`$ (1)
$`{\displaystyle \underset{i=\mathrm{hid},\mathrm{vis}}{}}{\displaystyle _{B_i}}d^dx\sqrt{\gamma }U_i(\varphi )].`$
Here $`M`$ is the bulk manifold which usually corresponds to AdS and $`B_{\mathrm{hid}}`$ and $`B_{\mathrm{vis}}`$ are branes corresponding to hidden and visible sectors respectively. $`\gamma `$ is the metric on the brane induced by the metric $`g`$ in the bulk. Here $`U_i(\varphi )`$ corresponds to the vacuum energies on the branes in . One assumes $`U(\varphi )`$ is dilaton dependent and its form is explicitly given later from the consistency of the equations of motion. Some important examples of the dilaton potential are presented in , where $`V(\varphi )`$ is given in terms of the superpotential $`W(\varphi )`$ :
$$V=\left(\frac{W}{\varphi }\right)^2\frac{D1}{2(D2)}W^2,$$
(2)
and $`W`$ has the following form:
$$W=\sqrt{\frac{𝒩}{2}}g\left(\frac{1}{a_1}\mathrm{e}^{\frac{a_1\varphi }{2}}\pm \frac{1}{a_2}\mathrm{e}^{\frac{a_2\varphi }{2}}\right).$$
(3)
Here $`D=d+1`$, $`𝒩`$ is the number of the supercharges and the parameter $`g`$, $`a_1`$ and $`a_2`$ depend on the model features but $`a_1>a_2>0`$ in general. As in , we only consider the case of $``$ sign of $`\pm `$ in (3). Note that potentials of above type appear as a result of sphere reduction in M-theory or string theory .
We now assume the metric has the following form:
$$ds^2=dz^2+\mathrm{e}^{2A(z)}\eta _{ij}dx^idx^j,$$
(4)
and $`\varphi `$ only depends on $`z`$. We also suppose the hidden and visible branes sit on $`z=z_{\mathrm{hid}}`$ and $`z=z_{\mathrm{vis}}`$, respectively. Then the equations of motion are given by
$`\varphi ^{\prime \prime }+(D1)A^{}\varphi ^{}={\displaystyle \frac{V}{\varphi }}+{\displaystyle \underset{i=\mathrm{hid},\mathrm{vis}}{}}{\displaystyle \frac{U_i(\varphi )}{\varphi }}\delta (zz_i),`$ (5)
$`(D1)A^{\prime \prime }+(D1)(A^{})^2+{\displaystyle \frac{1}{2}}(\varphi ^{})^2`$
$`={\displaystyle \frac{V}{D2}}V{\displaystyle \frac{D1}{2(D2)}}{\displaystyle \underset{i=\mathrm{hid},\mathrm{vis}}{}}U_i(\varphi )\delta (zz_i),`$ (6)
$`A^{\prime \prime }+(D1)(A^{})^2={\displaystyle \frac{1}{D2}}V{\displaystyle \frac{1}{2(D2)}}{\displaystyle \underset{i=\mathrm{hid},\mathrm{vis}}{}}U_i(\varphi )\delta (zz_i).`$ (7)
Here $`{}_{}{}^{}\frac{d}{dz}`$. For purely bulk sector ($`z_{\mathrm{hid}}<z<z_{\mathrm{vis}}`$, as $`z_{\mathrm{hid}}<z_{\mathrm{vis}}`$), the explicit solutions are given in . Eqs. (5-7) have the following first integrals (in the bulk sector):
$`\varphi ^{}=\sqrt{2}{\displaystyle \frac{W}{\varphi }},A^{}={\displaystyle \frac{1}{\sqrt{2}(D2)}}W.`$ (8)
Near the branes, Eqs. (5-7) have the following form :
$$\varphi ^{\prime \prime }\frac{U_i(\varphi )}{\varphi }\delta (zz_i),A^{\prime \prime }\frac{U_i(\varphi )}{2(D2)}\delta (zz_i),$$
(9)
or
$$2\varphi ^{}\frac{U_i(\varphi )}{\varphi },2A^{}\frac{U_i(\varphi )}{2(D2)},$$
(10)
at $`z=z_i`$. Comparing (10) with (8), we find
$$U_{\mathrm{hid}}(\varphi )=2\sqrt{2}W(\varphi ),U_{\mathrm{vis}}(\varphi )=2\sqrt{2}W(\varphi ).$$
(11)
For simplicity, let $`z_{\mathrm{hid}}=0`$, where $`\varphi =0`$ in the solution in and we only consider this solution in the following. Then at $`z=z_{\mathrm{vis}}`$, $`\varphi `$ is negative. Since the superpotential is given by the exponential of $`\varphi `$ with positive sign, the vacuum energy $`U_{\mathrm{vis}}`$, which can be identified with the cosmological constant, can be small even if $`|\varphi |`$ is not so large. The result might explain why the cosmological constant is small. We should note that the visible brane corresponds to the 4d universe where we live. Then the visible brane part in the action (1) mainly describes the dynamics of the universe and $`U_{\mathrm{vis}}`$ corresponds to the vacuum energy or cosmological constant in our universe.
As an extension, one can consider the case that the brane is curved. Instead of (4), we take the following metric:
$$ds^2=dz^2+\mathrm{e}^{2A(z)}\stackrel{~}{g}_{ij}dx^idx^j,$$
(12)
Here $`\stackrel{~}{g}_{ij}`$ is the metric of the Einstein manifold, which is defined by
$$\stackrel{~}{R}_{ij}=k\stackrel{~}{g}_{ij},$$
(13)
where $`\stackrel{~}{R}_{ij}`$ is the Ricci tensor given by $`\stackrel{~}{g}_{ij}`$ and $`k`$ is a constant. Then Eqs.(5) and (2) do no change but one obtains the following equation instead of (7):
$$A^{\prime \prime }+(D1)(A^{})^2=k\mathrm{e}^{2A}\frac{1}{D2}V\frac{1}{2(D2)}\underset{i=\mathrm{hid},\mathrm{vis}}{}U_i(\varphi )\delta (zz_i).$$
(14)
Especialy when $`k=0`$, we get the previous solution for $`\varphi `$, $`A`$ and $`U_i`$. We should note, however, that $`k=0`$ does not always mean the brane is flat. As well-known, the Einstein equations are given by,
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R+\frac{1}{2}\mathrm{\Lambda }g_{\mu \nu }=T_{\mu \nu }^{\mathrm{matter}}.$$
(15)
Here $`T_{\mu \nu }^{\mathrm{matter}}`$ is the energy-momentum tensor of the matter fields. If we consider the vacuum solution where $`T_{\mu \nu }^{\mathrm{matter}}=0`$, Eq.(15) can be rewritten when $`d=4`$ as (see also )
$$R_{\mu \nu }=\frac{\mathrm{\Lambda }}{2}g_{\mu \nu }.$$
(16)
If we put $`\mathrm{\Lambda }=2k`$, Eq.(16) is nothing but the equation for the Einstein manifold. The Einstein manifolds are not always homogeneous manifolds like flat Minkowski, (anti-)de Sitter space
$$ds_4^2=V(r)dt^2+V^1(r)dr^2+r^2d\mathrm{\Omega }^2,V(r)=1\frac{\mathrm{\Lambda }}{6}r^2,$$
(17)
or Nariai space
$$ds_4^2=\frac{1}{\mathrm{\Lambda }}\left(\mathrm{sin}^2\chi d\psi ^2d\chi ^2d\mathrm{\Omega }^2\right).$$
(18)
but they can be some black hole solutions like Schwarzschild-(anti-)de Sitter black hole
$$ds_4^2=V(r)dt^2+V^1(r)dr^2+r^2d\mathrm{\Omega }^2,V(r)=1\frac{\stackrel{~}{G}_4M}{r}\frac{\mathrm{\Lambda }}{6}r^2.$$
(19)
As a special case, one can also consider $`k=0`$ solution like Schwarzschild black hole,
$$ds_4^2\stackrel{~}{g}_{ij}dx^idx^j=\left(1\frac{\stackrel{~}{G}_4M}{r}\right)dt^2+\frac{dr^2}{\left(1\frac{\stackrel{~}{G}_4M}{r}\right)}+r^2d\mathrm{\Omega }^2,$$
(20)
or Kerr one
$`ds_4^2`$ $`=`$ $`\mathrm{\Delta }\stackrel{~}{A}dt^2{\displaystyle \frac{\mathrm{\Sigma }^2}{\mathrm{\Delta }}}dr^2\mathrm{\Sigma }^2d\theta ^2{\displaystyle \frac{\mathrm{sin}^2\theta }{\stackrel{~}{A}}}\left(d\phi \mathrm{\Omega }dt\right)`$ (21)
$`\stackrel{~}{A}{\displaystyle \frac{\mathrm{\Sigma }^2\left(\mathrm{\Delta }a^2\mathrm{sin}^2\theta \right)}{\mathrm{\Sigma }^4\mathrm{\Delta }4a^2\stackrel{~}{G}_4^2M^2r^2\mathrm{sin}^2\theta }},\mathrm{\Omega }{\displaystyle \frac{2a\stackrel{~}{G}_4Mr\stackrel{~}{A}}{\mathrm{\Sigma }^2}}`$
$`\mathrm{\Delta }r^22\stackrel{~}{G}_4Mr+a^2,\mathrm{\Sigma }^2r^2+a^2\mathrm{cos}^2\theta .`$
In (19), (20) and (21), $`M`$ is the mass of the black hole on the brane and the effective gravitational constant $`\stackrel{~}{G}_4`$ on the 3-brane (here $`d=4`$) is given by
$$\frac{1}{\stackrel{~}{G}_4}=\frac{1}{G_5}_{z_{\mathrm{hid}}}^{z_{\mathrm{vis}}}𝑑z\mathrm{e}^{(d2)A}.$$
(22)
In (21), the parameter $`a`$ is related with the angular momentum $`J`$ of the black hole on the brane by
$$J=Ma.$$
(23)
In these solutions, the curvature singularity at $`r=0`$ has a form of line penetrating the bulk 5d universe and the horizon makes a tube surrounding the singularity. The singularity and the horizon connect the hidden and visible branes. Thus we presented a big family of dilatonic solutions in the brane world.
It would be interesting to see how the black hole looks like in the bulk (see corresponding discussion in ). We now consider the case $`|z_{\mathrm{vis}}|`$ is large but $`z_{\mathrm{vis}}<0`$ or $`z_{\mathrm{vis}}>0`$. When $`z_{\mathrm{vis}}`$ is negative and its absolute value $`|z_{\mathrm{vis}}|`$ is large, the asymptotic behaviour of $`A`$ is given by
$$A\gamma z,\gamma \frac{g}{a_1d\sqrt{𝒩}}.$$
(24)
Then since the metric has the form in (12), the typical proper size $`l_{\mathrm{BH}}^{}`$ of the black hole (or the typical proper distance to the horizon) transverse to the brane would be given by
$$l_{\mathrm{BH}}^{}=\frac{1}{\gamma }\mathrm{ln}\left(\gamma r_0\right)=\frac{1}{\gamma }\mathrm{ln}\left(\gamma \stackrel{~}{G}_4M\right).$$
(25)
Then the transverse size $`l_{\mathrm{BH}}^{}`$ grows like $`\mathrm{ln}M`$ when the black hole mass $`M`$ increases although the (horizon) size $`r_0`$ along the brane is linear to $`M`$. This might tell that matter can pass around the black hole through the fifth dimension similar to the effect suggested in .
On the other hand, we can consider the case that $`z_{\mathrm{vis}}`$ is positive and large. The asymptotic behaviour of $`A`$ when $`z`$ is large, is given by
$$\mathrm{e}^{2A}\alpha ^2z^{2\beta },\alpha \left\{c\sqrt{\frac{𝒩}{8}}g\left(\frac{1}{4}a_2\sqrt{𝒩}g\right)^{\frac{a_1}{a_2}}\right\}^{\frac{1}{d}},\beta \frac{a_1}{a_2d}.$$
(26)
Then if one defines a new coordinate $`w`$ by
$$y=\frac{z^{1\beta }}{\alpha (1\beta )},$$
(27)
the metric in (12) can be rewritten in the following form:
$$ds^2\alpha ^2\left\{\alpha (1\beta )y\right\}^{\frac{2\beta }{1\beta }}\left(dw^2+\stackrel{~}{g}_{ij}dx^idx^j\right).$$
(28)
Then we can expect the horizon in the bulk is typically given by $`yr_0`$ and the typical proper size $`l_{\mathrm{BH}}^+`$ of the black hole transverse to the brane would be given by
$$l_{\mathrm{BH}}^+=\left\{\alpha (1\beta )r_0\right\}^{\frac{1}{1\beta }}=\left\{\alpha (1\beta )\stackrel{~}{G}_4M\right\}^{\frac{1}{1\beta }}.$$
(29)
Since $`a_1=2\sqrt{\frac{5}{3}}`$ and $`a_2=\frac{4}{\sqrt{15}}`$ for $`d=D1=4`$ and $`𝒩=1`$ model in , we have $`\frac{1}{1\beta }=\frac{5}{3}`$ and the transverse size $`l_{\mathrm{BH}}^+`$ grows like $`M^{\frac{5}{3}}`$ when the black hole mass $`M`$ increases. This would tell that matter cannot pass through a black hole, which is different from the case when $`z_{\mathrm{vis}}\mathrm{}`$. Of course, this is a preliminary qualitative discussion of brane-world black holes properties.
## 3 Gravity perturbations
Our next problem will be the description of the gravity on the brane. For that purpose one can consider the perturbation from the above obtained solution by assuming the metric in the following form:
$$ds^2=dz^2+\mathrm{e}^{2A(z)}\stackrel{~}{g}_{ij}dx^idx^j+h_{ij}dx^idx^j,|h_{ij}|1.$$
(30)
Here we choose the gauge where $`g_{zz}=1`$, $`g_{zi}=g_{iz}=0`$, $`D^ih_{ij}=0`$ and $`h_i^i=0`$. We should note that the metric $`\stackrel{~}{g}_{ij}`$ on the brane is not necessary to be flat but can be any Einstein manifold (13). We put $`k=0`$ for simplicity. Then $`A(z)`$ is given by solving (8) . From the Einstein equation, we obtain the following linearized equation:
$$0=\mathrm{e}^{2A}\stackrel{~}{\mathrm{}}^{(0)}h_{ij}+_z^2h_{ij}+(d4)_zA_zh_{ij}4\left\{(d1)_zA^2+_z^2A\right\}h_{ij}.$$
(31)
Here $`\stackrel{~}{\mathrm{}}^{(0)}`$ is the d’Alembertian on the brane given by $`\stackrel{~}{g}_{ij}`$. Changing the coordinate $`z`$ to $`\zeta `$ by
$$d\zeta =\mathrm{e}^Adz,$$
(32)
one rewrites (31) in the following form:
$$0=\stackrel{~}{\mathrm{}}^{(0)}h_{ij}+_\zeta ^2h_{ij}+(d5)_\zeta A_\zeta h_{ij}4\left\{(d2)_\zeta A^2+_\zeta ^2A\right\}h_{ij}.$$
(33)
In order to solve Eq.(33), we assume the following form for $`h_{ij}`$:
$$h_{ij}=\psi (\zeta )\widehat{h}_{ij}(x)$$
(34)
and assume $`\widehat{h}_{ij}(x)`$ satisfies
$$\stackrel{~}{\mathrm{}}^{(0)}\widehat{h}_{ij}(x)=m^2\widehat{h}_{ij}(x).$$
(35)
Here $`m`$ corresponds to the mass. Then one has
$$0=m^2\psi +_\zeta ^2\psi +(d5)_\zeta A_\zeta \psi 4\left\{(d2)_\zeta A^2+_\zeta ^2A\right\}\psi .$$
(36)
As we have a dilaton field $`\varphi `$, we should consider the perturbation of $`\varphi `$:
$$\varphi =\varphi ^{(0)}+\phi ,|\phi ||\varphi ^{(0)}|.$$
(37)
Here $`\varphi ^{(0)}`$ is given by solving (8). Linearizing the equation of motion given by the variation of $`\varphi `$,
$$\mathrm{}\varphi =\frac{V(\varphi )}{\varphi }+\underset{i}{}\frac{U_i(\varphi )}{\varphi }\delta (zz_i),$$
(38)
one gets
$$\mathrm{e}^{2A}\stackrel{~}{\mathrm{}}^{(0)}\phi +_z^2\phi +d_zA_z\phi =\left\{\frac{^2V(\varphi )}{\varphi ^2}+\underset{i}{}\frac{^2U_i(\varphi )}{\varphi ^2}\delta (xz_i)\right\}\phi .$$
(39)
Using the coordinate $`\zeta `$ defined by (32) and assuming $`\varphi `$ in the following form:
$$\phi (\zeta ,x)=\theta (\zeta )\widehat{\phi }(x),\stackrel{~}{\mathrm{}}^{(0)}\widehat{\phi }(x)=m^2\widehat{\phi }(x),$$
(40)
we can rewrite Eq.(39) as:
$$m^2\theta +_\zeta ^2\theta +(d1)_\zeta A_\zeta \theta =\mathrm{e}^{2A}\left\{\frac{^2V(\varphi )}{\varphi ^2}+\underset{i}{}\frac{^2U_i(\varphi )}{\varphi ^2}\delta (zz_i)\right\}\theta .$$
(41)
In the solutions found in , there appear two branches. In the first branch, $`z`$ runs from $`\mathrm{}`$ to $`0`$, $`\varphi `$ from $`0`$ to $`+\mathrm{}`$ and there appears a curvature singularity at $`z=0`$. In the second one $`z`$ runs from $`\mathrm{}`$ to $`+\mathrm{}`$, $`\varphi `$ from $`0`$ to $`\mathrm{}`$ and there does not appear any curvature singularity. In both of the branches, the spacetime approaches to AdS.
We now consider the perturbation when $`z_{\mathrm{vis}}`$ is negative and its absolute value $`|z_{\mathrm{vis}}|`$ is large and the region $`zz_{\mathrm{vis}}`$. Then we find the following asymptotic behaviour of $`A`$ and $`\varphi `$:
$$A\gamma z,\varphi \mathrm{e}^{d\gamma z}0,\gamma \frac{g}{a_1d\sqrt{𝒩}}.$$
(42)
Then $`\zeta `$ in (32) is given by
$$\zeta =\frac{\mathrm{e}^{\gamma z}}{\gamma }$$
(43)
and Eq.(36) has the following form
$$0=m^2\psi +_\zeta ^2\psi +\frac{5d}{\zeta }_\zeta \psi \frac{4(d1)}{\zeta ^2}\psi .$$
(44)
When $`m>0`$, the corresponding solution is called Kaluza-Klein (KK) mode and the solution of $`m=0`$ corresponds to the graviton on the brane. When $`m^2=0`$, the solution of (44) is given by the power of $`\zeta `$:
$$\psi =\zeta ^a.$$
(45)
The exponent $`a`$ can be found by solving the following algebraic equation:
$$0=a^2+(4d)a4(d1).$$
(46)
Especially when $`d=4`$, one gets
$$a=\pm 2\sqrt{3}.$$
(47)
The existence of the normalizable solution tells that the gravity is localized near the brane. This situation does not change even if the brane is flat or a (4d) black hole spacetime.
When $`m^2>0`$, the solution of (44) is given by Bessel functions $`J_\nu `$ and $`N_\nu `$ ($`=Y_\nu `$):
$`\psi (\zeta )=\zeta ^b\left(c_1J_\nu (m\zeta )+c_2N_\nu (m\zeta )\right)`$
$`12b=5d,b^2\nu ^2=4(d1).`$ (48)
Here $`c_1`$ and $`c_2`$ are constants of integration, which should be determined by the boundary condition at $`z=z_{\mathrm{vis}}`$:
$`_z\psi (z=z_{\mathrm{vis}})`$ $`=`$ $`\mathrm{e}^{A(\zeta =\zeta _{\mathrm{vis}})}_\zeta \psi (\zeta =\zeta _{\mathrm{vis}})`$ (49)
$`=`$ $`{\displaystyle \frac{2\sqrt{2}}{d3}}W\left(\varphi (z=z_{\mathrm{vis}})\right)\psi (z=z_{\mathrm{vis}}).`$
Here $`\zeta _{\mathrm{vis}}`$ is the value of $`\zeta `$ corresponding to $`z=z_{\mathrm{vis}}`$. The boundary condition comes from the $`\delta `$-function behaviour of $`_z^2A`$ at $`z=z_{\mathrm{vis}}`$. Especially for $`d=4`$, we find
$$b=0,\nu =2\sqrt{3}.$$
(50)
If there is a solution with $`m^2<0`$, the system becomes unstable. When $`m^2<0`$, the solution of (44) is given by modified Bessel functions $`I_\nu `$ and $`K_\nu `$:
$`\psi (\zeta )=\zeta ^b\left(c_1I_\nu (\mu \zeta )+c_2K_\nu (\mu \zeta )\right)`$
$`12b=5d,b^2\nu ^2=4(d1),\mu ^2=m^2.`$ (51)
Since $`I_\mu `$ increases exponentially for large $`\zeta `$, $`c_1`$ must vanish. On the other hand, $`K_\nu `$ behaves as $`K_\mu (m\zeta )\mathrm{e}^{\mu \zeta }`$ for large $`\zeta `$. When $`\zeta `$ is large, $`z\mathrm{}`$, $`\varphi 0`$ and $`W\sqrt{\frac{𝒩}{2}}g\left(\frac{1}{a_1}\frac{1}{a_2}\right)<0`$. Therefore there is an unstable solution which satisfies the boundary condition (49) if one chooses
$$\mu \frac{2\sqrt{2}}{d3}\mathrm{e}^{A(z=z_{\mathrm{vis}})}W\left(\varphi (z=z_{\mathrm{vis}})\right)\frac{2\sqrt{2}}{d3}\mathrm{e}^{\gamma z_{\mathrm{vis}}}\sqrt{\frac{𝒩}{2}}g\left(\frac{1}{a_1}\frac{1}{a_2}\right).$$
(52)
As $`\mu `$ vanishes in the limit of $`z\mathrm{}`$, the brane will be driven to $`z\mathrm{}`$.
We now consider the perturbation by the dilaton $`\varphi `$. When $`\zeta `$ is large, Eq.(41) can be rewritten as
$$m^2\theta +_\zeta ^2\theta \frac{d1}{\zeta }_\zeta \theta \frac{V_0^{\prime \prime }}{\gamma ^2\zeta ^2}\theta =\mathrm{e}^{2A}\underset{i}{}\frac{^2U_i(\varphi )}{\varphi ^2}\delta (zz_i)\theta .$$
(53)
Here
$$V_0^{\prime \prime }\frac{^2V(\varphi )}{\varphi ^2}|_{\varphi =0}=\frac{1}{g^2}\left(1\frac{a_2}{a_1}\right)>0.$$
(54)
Then one finds $`\theta `$ is given by the power of $`\zeta `$ when $`m^2=0`$ or Bessel functions when $`m^2>0`$. From the $`\delta `$-function, we find that $`\theta `$ should satisfy the following boundary condition $`\zeta =\zeta _{\mathrm{vis}}`$ :
$$_\zeta \theta =\frac{\sqrt{2}}{\gamma \zeta _{\mathrm{vis}}}\frac{^2W}{\varphi ^2}|_{\varphi =0}\theta .$$
(55)
Since
$$\frac{^2W}{\varphi ^2}|_{\varphi =0}=\sqrt{\frac{𝒩}{2}}\frac{g(a_1a_2)}{4}>0,$$
(56)
the coefficient $`\frac{\sqrt{2}}{\gamma \zeta _{\mathrm{vis}}}\frac{^2W}{\varphi ^2}|_{\varphi =0}`$ is negative. Then the boundary condition is consistent with the modified Bessel function $`K_\nu `$ and the condition could be satisfied by properly choosing $`m^2`$. This tells that there exists an unstable mode corresponding to $`m^2<0`$.
In the second branch in the solution in , we can consider the case that $`z_{\mathrm{vis}}`$ is positive and large. We also consider the region $`zz_{\mathrm{vis}}`$. Since the asymptotic behaviour of $`A`$ when $`z`$ is large is given by (26), one gets
$$\zeta \frac{z^{1\beta }}{1\beta }$$
(57)
and Eq.(33) has the following form
$$0=\stackrel{~}{\mathrm{}}^{(0)}h_{ij}+_\zeta ^2h_{ij}+\frac{(d5)\stackrel{~}{\beta }}{\zeta }_\zeta h_{ij}4\left\{(d2)\stackrel{~}{\beta }^2\stackrel{~}{\beta }\right\}h_{ij}.$$
(58)
Here
$$\stackrel{~}{\beta }\frac{\beta }{1\beta }.$$
(59)
and we obtain
$$0=m^2\psi +_\zeta ^2\psi +\frac{(d5)\stackrel{~}{\beta }}{\zeta }_\zeta \psi 4\left\{(d2)\stackrel{~}{\beta }^2\stackrel{~}{\beta }\right\}\psi .$$
(60)
When $`m=0`$, the solution is given by
$$\psi =(\zeta )^a.$$
(61)
Here $`a`$ is the solution of the following algebraic equation:
$$0=a^2+\left\{(d5)\stackrel{~}{\beta }1\right\}a4\left\{(d2)\stackrel{~}{\beta }^2\stackrel{~}{\beta }\right\}.$$
(62)
For $`d=D1=4`$ and $`𝒩=1`$ model in , where $`a_1=2\sqrt{\frac{5}{3}}`$ and $`a_2=\frac{4}{\sqrt{15}}`$, we have
$$a=\frac{4\pm 2\sqrt{39}}{3}=5.49\mathrm{},2.83\mathrm{}$$
(63)
Since we are considering the case that $`z_{\mathrm{vis}}`$ is positive and $`zz_{\mathrm{vis}}`$, the first positive $`a`$ would correspond to a normalizable solution. The existence of the normalizable solution tells that the gravity near the brane is localized.
When $`m>0`$, the solution of (60) is given by Bessel functions $`J_\nu `$ and $`N_\nu `$ ($`=Y_\nu `$):
$`\psi (\zeta )=z^b\left(c_1J_\nu (m\zeta )+c_2N_\nu (m\zeta )\right)`$
$`12b=(d5)\stackrel{~}{\beta },b^2\nu ^2=4\left\{(d2)\stackrel{~}{\beta }^2\stackrel{~}{\beta }\right\}.`$ (64)
Here $`c_1`$ and $`c_2`$ are constants of integration, which should be again determined by the boundary condition (49) at $`z=z_{\mathrm{vis}}`$. The boundary condition comes from the $`\delta `$-function behaviour of $`_z^2A`$ at $`z=z_{\mathrm{vis}}`$. When $`\zeta `$ ($`z`$) is large, the Bessel functions $`J_\nu (m\zeta )`$ and $`N_\nu (m\zeta )`$ behave as
$`J_\nu (m\zeta )`$ $``$ $`\sqrt{{\displaystyle \frac{2}{\pi m\zeta }}}\mathrm{cos}\left(m\zeta {\displaystyle \frac{(2\nu +1)\pi }{4}}\right),`$
$`N_\nu (m\zeta )`$ $``$ $`\sqrt{{\displaystyle \frac{2}{\pi m\zeta }}}\mathrm{sin}\left(m\zeta {\displaystyle \frac{(2\nu +1)\pi }{4}}\right).`$ (65)
For $`d=D1=4`$ and $`𝒩=1`$ model in , we have
$$b=\frac{4}{3},\nu ^2=\frac{439}{9}=(4.16\mathrm{})^2.$$
(66)
The existence of the unstable mode corresponding to modified Bessel function $`K_\mu (\mu \zeta )`$ $`(\mu ^2=m^2`$ depends on the asymptotic behaviour of $`W`$ as in $`z_{\mathrm{vis}}\mathrm{}`$ case. From (3) ($``$ sign is chosen), we find $`W<0`$ when $`\varphi \mathrm{}`$ and $`\mathrm{e}^{A\left(\varphi (z)\right)}W\left(\varphi (z)\right)`$ vansihes when $`z(\zeta )+\mathrm{}`$. Therefore an unstable solution given by $`K_\mu (\mu \zeta )`$ exists and the brane moves to $`+\mathrm{}`$, where the vacuum energy of the brane vanishes since $`\mu 0`$ when $`z_{\mathrm{vis}}\mathrm{}`$.
We also consider the perturbation by the dilaton $`\varphi `$. When $`\zeta `$ is large, Eq.(41) can be rewritten as
$$m^2\theta +_\zeta ^2\theta +\frac{\stackrel{~}{\beta }(d1)}{\zeta }_\zeta \theta +\frac{\stackrel{~}{V}_0}{\gamma ^2\zeta ^2}\theta =\mathrm{e}^{2A}\underset{i}{}\frac{^2U_i(\varphi )}{\varphi ^2}\delta (zz_i)\theta .$$
(67)
Here
$$\stackrel{~}{V}_0\frac{1}{2}g^2\left(1\frac{a_1}{da_2}\right)^2\left(\frac{a_2}{a_1}\right)\left(\frac{1}{4}a_2\sqrt{𝒩}g\right)^2.$$
(68)
Then we find $`\theta `$ is given by the power of $`\zeta `$ when $`m^2=0`$ or Bessel functions when $`m^2>0`$. The existence of the unstable mode corresponding the modified Bessel function depends on the sign of $`\frac{^2W}{\varphi ^2}`$ at $`z=z_{\mathrm{vis}}`$. Since
$$\frac{^2W}{\varphi ^2}|_{z=z_{\mathrm{vis}}}\sqrt{\frac{𝒩}{2}}g\frac{a_2}{4}\mathrm{e}^{\frac{a_2\varphi }{2}}<0,$$
(69)
for large positive $`z_{\mathrm{vis}}`$, the sign is negative, which is different from the case $`z_{\mathrm{vis}}\mathrm{}`$. Therefore, remarkably there does not exist an unstable mode corresponding to $`m^2<0`$ when we choose $`z_{\mathrm{vis}}+\mathrm{}`$, where the cosmological constant becomes very small. To conclude, the above analysis shows the possibility to get the localized gravity near the brane corresponding to our black hole solution. Note that excellent introduction to universal aspects of gravity localization in brane world can be found in ref..
## 4 Newton law correction
It is interesting to discuss now the correction to the Newton law coming from KK mode.
First we consider $`z_{\mathrm{vis}}\mathrm{}`$ case in (3). The constants of integration should be determined by the boundary condition at $`z=z_{\mathrm{vis}}`$ in (49). As the overall factor $`\sqrt{\frac{2}{\pi m\zeta }}`$ should be absorbed into the measure for the normalization, we replace it with unity and we impose a constraint as follows:
$$c_1^2+c_2^2=1.$$
(70)
Then when $`d=4`$, and $`m\zeta `$ is large, $`\psi (\zeta )`$ behaves as
$$\psi (\zeta )c_1\text{cos}\left(m\zeta \frac{(4\sqrt{3}+1)\pi }{4}\right)+c_2\text{sin}\left(m\zeta \frac{(4\sqrt{3}+1)\pi }{4}\right)$$
(71)
so then $`_\zeta \psi (\zeta )`$ is given by
$$_\zeta \psi (\zeta )mc_1\text{sin}\left(m\zeta \frac{(4\sqrt{3}+1)\pi }{4}\right)+mc_2\text{cos}\left(m\zeta \frac{(4\sqrt{3}+1)\pi }{4}\right).$$
(72)
We should note that there might be some ambiguities in the limiting procedure. Since we consider $`\zeta _{\mathrm{vis}}+\mathrm{}`$, we first put $`m\zeta `$ to be large and after that we choose the mass $`m`$ in the KK mode to be small, which is relevant to the long range force. Then by using the boundary condition (49) and the constraint (70), we get $`c_1,c_2`$ in following forms:
$`c_1`$ $`=`$ $`\pm \kappa \sqrt{{\displaystyle \frac{1}{\kappa ^2+1}}}`$
$`c_2`$ $`=`$ $`\pm \sqrt{{\displaystyle \frac{1}{\kappa ^2+1}}}`$ (73)
Here $`\kappa `$ is defined by
$`\kappa `$ $``$ $`{\displaystyle \frac{\text{cos}\zeta ^{}+Z(\zeta =\zeta _{\mathrm{vis}})\text{sin}\zeta ^{}}{\text{sin}\zeta ^{}Z(\zeta =\zeta _{\mathrm{vis}})\text{cos}\zeta ^{}}}`$
$`Z(\zeta =\zeta _{\mathrm{vis}})`$ $``$ $`{\displaystyle \frac{2\sqrt{2}e^{A(\zeta =\zeta _{\mathrm{vis}})}W(\varphi (z=z_{\mathrm{vis}}))}{m}}`$
$`\zeta ^{}`$ $``$ $`m\zeta _{\mathrm{vis}}{\displaystyle \frac{(4\sqrt{3}+1)\pi }{4}}`$ (74)
One can find the value of $`\psi `$ (71) at $`\zeta =\zeta _{\mathrm{vis}}`$ using (4), (4):
$`\psi (\zeta _{\mathrm{vis}})`$ $``$ $`\sqrt{{\displaystyle \frac{1}{\kappa ^2+1}}}\left(\kappa \text{cos}\zeta ^{}+\text{sin}\zeta ^{}\right)`$ (75)
$`=`$ $`\sqrt{{\displaystyle \frac{1}{\kappa ^2+1}}}\left({\displaystyle \frac{1}{\text{sin}\zeta ^{}+Z(\zeta =\zeta _{\mathrm{vis}})\text{cos}\zeta ^{}}}\right)`$
$`=`$ $`{\displaystyle \frac{\text{sin}\zeta ^{}+Z(\zeta =\zeta _{\mathrm{vis}})\text{cos}\zeta ^{}}{\sqrt{1+Z(\zeta =\zeta _{\mathrm{vis}})^2}}}\left({\displaystyle \frac{1}{\text{sin}\zeta ^{}+Z(\zeta =\zeta _{\mathrm{vis}})\text{cos}\zeta ^{}}}\right)`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{1+Z(\zeta =\zeta _{\mathrm{vis}})^2}}}`$
If we take $`m`$ is small (or $`|Z|`$ is large), then $`\psi (\zeta _{\mathrm{vis}})`$ is given by
$`\psi {\displaystyle \frac{1}{Z(\zeta =\zeta _{\mathrm{vis}})}}={\displaystyle \frac{m}{2\sqrt{2}\mathrm{e}^{A(\zeta =\zeta _{\mathrm{vis}})}W(\varphi (z=z_{\mathrm{vis}}))}}`$ (76)
Then the correction to Newton’s Law is
$`V(r)`$ $``$ $`\stackrel{~}{G}_4{\displaystyle \frac{m_1m_2}{r}}+{\displaystyle _0^{\mathrm{}}}𝑑mG_5{\displaystyle \frac{m_1m_2e^{mr}}{r}}\psi (\zeta =\zeta _{\mathrm{vis}})^2`$ (77)
$`=`$ $`\stackrel{~}{G}_4{\displaystyle \frac{m_1m_2}{r}}+{\displaystyle _0^{\mathrm{}}}𝑑mG_5{\displaystyle \frac{m_1m_2e^{mr}}{r}}{\displaystyle \frac{m^2}{8\mathrm{e}^{2A(\zeta =\zeta _{\mathrm{vis}})}W(\varphi (z=z_{\mathrm{vis}}))^2}}`$
$`=`$ $`\stackrel{~}{G}_4{\displaystyle \frac{m_1m_2}{r}}\left(1+{\displaystyle \frac{G_5}{\stackrel{~}{G}_4}}{\displaystyle \frac{1}{r^34\mathrm{e}^{2A(\zeta =\zeta _{\mathrm{vis}})}W(\varphi (z=z_{\mathrm{vis}}))^2}}\right)`$
We should note that that the correction is not given by $`\frac{1}{r^3}`$ as in but $`\frac{1}{r^4}`$. This is mainly due to the limiting procedure where we first have put $`m\zeta `$ to be large and after that we have chosen the mass $`m`$ in the KK mode to be small.
Next we consider the case $`z_{\mathrm{vis}}\mathrm{}`$ given in (3). When $`\zeta `$ ($`z`$) is large, the Bessel functions $`J_\nu (m\zeta )`$ and $`N_\nu (m\zeta )`$ behave as
$`J_\nu (m\zeta )`$ $``$ $`\sqrt{{\displaystyle \frac{2}{\pi m\zeta }}}\mathrm{cos}\left(m\zeta {\displaystyle \frac{(2\nu +1)\pi }{4}}\right),`$
$`N_\nu (m\zeta )`$ $``$ $`\sqrt{{\displaystyle \frac{2}{\pi m\zeta }}}\mathrm{sin}\left(m\zeta {\displaystyle \frac{(2\nu +1)\pi }{4}}\right).`$ (78)
For $`d=D1=4`$ and $`𝒩=1`$ model in , we have the parameters $`b`$, $`\nu `$ given in (66). Then $`c_1`$ and $`c_2`$ in (3) should be determined by the boundary condition at $`z=z_{\mathrm{vis}}`$ in (49) and constraint (70) again. Since
$$_\zeta \psi (\zeta )mc_1\text{sin}\left(m\zeta \frac{(2\nu +1)\pi }{4}\right)mc_2\text{cos}\left(m\zeta \frac{(2\nu +1)\pi }{4}\right),$$
(79)
we can get $`c_1,c_2`$ in the same way as (4).
$`c_1`$ $`=`$ $`\pm \kappa \sqrt{{\displaystyle \frac{1}{\kappa ^2+1}}}`$
$`c_2`$ $`=`$ $`\pm \sqrt{{\displaystyle \frac{1}{\kappa ^2+1}}}`$ (80)
Here $`\kappa `$ is defined by
$`\kappa `$ $``$ $`{\displaystyle \frac{\text{cos}\zeta ^{}+Z(\zeta =\zeta _{\mathrm{vis}})\text{sin}\zeta ^{}}{\text{sin}\zeta ^{}Z(\zeta =\zeta _{\mathrm{vis}})\text{cos}\zeta ^{}}}`$
$`Z(\zeta =\zeta _{\mathrm{vis}})`$ $``$ $`{\displaystyle \frac{2\sqrt{2}e^{A(\zeta =\zeta _{\mathrm{vis}})}W(\varphi (z=z_{\mathrm{vis}}))}{m}}`$
$`\zeta ^{}`$ $``$ $`m\zeta {\displaystyle \frac{(2\nu +1)\pi }{4}}`$ (81)
And $`\psi `$ (3) is written by using (4), (4).
$`\psi (\zeta )`$ $``$ $`\sqrt{{\displaystyle \frac{1}{\kappa ^2+1}}}\left(\kappa \text{cos}\zeta ^{}+\text{sin}\zeta ^{}\right)`$ (82)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{1+Z(\zeta =\zeta _{\mathrm{vis}})^2}}}`$
Then the correction to Newton’s law in the limit that $`m`$ is small (or $`|Z|`$ is large) is
$`V(r)`$ $``$ $`\stackrel{~}{G}_4{\displaystyle \frac{m_1m_2}{r}}+{\displaystyle _0^{\mathrm{}}}𝑑mG_5{\displaystyle \frac{m_1m_2e^{mr}}{r}}\psi ^2`$ (83)
$`=`$ $`\stackrel{~}{G}_N{\displaystyle \frac{m_1m_2}{r}}+{\displaystyle _0^{\mathrm{}}}𝑑mG_5{\displaystyle \frac{m_1m_2e^{mr}}{r}}{\displaystyle \frac{m^2}{8e^{2A(\zeta =\zeta _{\mathrm{vis}})}W(\varphi (z=z_{\mathrm{vis}}))^2}}`$
$`=`$ $`\stackrel{~}{G}_4{\displaystyle \frac{m_1m_2}{r}}\left(1+{\displaystyle \frac{G_5}{\stackrel{~}{G}_4}}{\displaystyle \frac{1}{4e^{2A(\zeta =\zeta _{\mathrm{vis}})}W(\varphi (z=z_{\mathrm{vis}}))^2}}{\displaystyle \frac{1}{r^3}}\right).`$
This is qualitatively the same type of correction as in ref.. Thus, the observer living on the brane Universe does not see drastic changes in the Newton law.
## 5 Dilatonic de Sitter brane Universe
When $`V`$ is constant, instead of (2), the solution is given in , as follows
$`ds^2`$ $`=`$ $`f(y)dy^2+y{\displaystyle \underset{i,j=0}{\overset{d1}{}}}\stackrel{~}{g}_{ij}(x^k)dx^idx^j`$
$`f`$ $`=`$ $`{\displaystyle \frac{d(d1)}{4y^2\lambda ^2\left(1+\frac{c^2}{2\lambda ^2y^d}+\frac{kd}{\lambda ^2y}\right)}}`$
$`\varphi `$ $`=`$ $`c{\displaystyle 𝑑y\sqrt{\frac{d(d1)}{4y^{d+2}\lambda ^2\left(1+\frac{c^2}{2\lambda ^2y^d}+\frac{kd}{\lambda ^2y}\right)}}}.`$ (84)
Here we define $`V=\lambda ^2`$ and $`g_{ij}`$ is the metric of the Einstein manifold, which is defined by $`r_{ij}=k\stackrel{~}{g}_{ij}`$, where $`r_{ij}`$ is the Ricci tensor constructed with $`\stackrel{~}{g}_{ij}`$ and $`k`$ is a constant. Especially when $`k`$ is given by $`k=d\widehat{c}^2>0`$, the metric on the brane can correspond to the cosmological solution
$$\underset{i,j=0}{\overset{d1}{}}\stackrel{~}{g}_{ij}(x^k)dx^idx^j=\frac{1}{\widehat{c}^2t^2}\left(dt^2+\underset{i=1}{\overset{d1}{}}\left(dx^i\right)^2\right).$$
(85)
This solution describes the wall expanding and travelling in 5d universe. It may be presented as regular solution in terms of radial coordinate similarly to black hole.
If one defines a new coordinate $`z`$ by
$$z=𝑑y\sqrt{\frac{d(d1)}{4y^2\lambda ^2\left(1+\frac{c^2}{2\lambda ^2y^d}+\frac{kd}{\lambda ^2y}\right)}}$$
(86)
and solves $`y`$ with respect to $`z`$, we obtain the warp factor $`\mathrm{e}^{2A}=y(z)`$. We should note that there is a curvature singularity at $`y=0`$ . Therefore we cannot put only one brane but two branes and consider the region sandwiched by them to avoid the singularity.
Using (10), one finds the vacuum energy on the brane as follows
$`U_{\mathrm{hid}}`$ $`=`$ $`4\lambda \sqrt{{\displaystyle \frac{d1}{d}}\left(1+{\displaystyle \frac{c^2}{2\lambda ^2y_{\mathrm{hid}}^d}}+{\displaystyle \frac{kd}{\lambda ^2y_{\mathrm{hid}}}}\right)}`$
$`U_{\mathrm{vis}}`$ $`=`$ $`4\lambda \sqrt{{\displaystyle \frac{d1}{d}}\left(1+{\displaystyle \frac{c^2}{2\lambda ^2y_{\mathrm{vis}}^d}}+{\displaystyle \frac{kd}{\lambda ^2y_{\mathrm{vis}}}}\right)}.`$ (87)
Here we assume $`z_{\mathrm{vis}}>z_{\mathrm{hid}}`$ and consider the region $`z_{\mathrm{vis}}zz_{\mathrm{hid}}`$. $`y_{\mathrm{hid},\mathrm{vis}}`$ is the value of $`y`$ corresponding to $`z=z_{\mathrm{hid},\mathrm{vis}}`$. From (10), one gets
$$\frac{U_{\mathrm{hid}}}{\varphi }=\frac{2c}{y_{\mathrm{hid}}^{\frac{d}{2}}},\frac{U_{\mathrm{vis}}}{\varphi }=\frac{2c}{y_{\mathrm{vis}}^{\frac{d}{2}}},$$
(88)
which tells that the vacuum energies $`U_{\mathrm{vis}}`$ on the brane depend on the dilaton field.
One can consider the perturbation around the solution (5), (85). As it is difficult to consider the general case, we first investigate the region $`y<y_{\mathrm{vis}}`$ and $`y_{\mathrm{vis}}0`$, that is, the brane is near the singularity. We should note, however, it is sufficient to consider the asymptotic region to check the existence of the normalized zero mode and continuous KK modes. Using (86), when $`y0`$, one gets
$$\zeta \frac{1}{c}\sqrt{\frac{d}{d1}}y^{\frac{d1}{2}}.$$
(89)
This tells that we now consider the region where $`\zeta `$ is small. Since $`A=\frac{1}{2}\mathrm{ln}y`$, Eq.(36) has the following form :
$$0=m^2\psi +_\zeta ^2\psi +\frac{d5}{d1}\frac{1}{\zeta }_\zeta \psi +\frac{4}{(d1)^2}\frac{1}{\zeta ^2}\psi .$$
(90)
When $`m^2=0`$, the solution is given by
$$\psi _0=\zeta ^{\frac{2}{d1}}\left(\stackrel{~}{c}_1+\stackrel{~}{c}_2\mathrm{ln}\zeta \right).$$
(91)
On the other hand, when $`m^2>0`$, we find
$$\psi _m=\zeta ^{\frac{2}{d1}}\left(c_1J_0(m\zeta )+c_2N_0(m\zeta )\right).$$
(92)
The mode corresponding to $`m=0`$ does not seem to be normalizable as in flat dilatonic brane . Then in order to localize the gravity, we need to put a brane corresponding to the hidden sector.
## 6 Quantum effective action for dilatonic brane and RG flow of Newton constant
Our starting point is again the following action of 5d dilatonic gravity (gauged supergravity):
$$S=\frac{1}{16\pi }\left[_Md^{d+1}x\sqrt{g}\left(R\frac{1}{2}_\mu \varphi ^\mu \varphi V(\varphi )\right)\right].$$
(93)
The dilatonic potential is not specified for the moment. As it was discussed in section 2, the very common choice for $`V`$ is the exponential of dilaton: It corresponds to the effective action for the breathing-mode scalar and gravity which follows from KK sphere reduction from M-theory or strings . For the exponential potentials there are singular domain wall solutions in above theory . It is interesting that there are usually problems with localization of 4d gravity when using only the action (93). As a result one should consider the inclusion of four dimensional action (for walls they correspond to wall source terms). Then
$$S_{\mathrm{source}}=\underset{i=\mathrm{hid},\mathrm{vis}}{}d^4x\sqrt{\gamma }\left\{L_{\mathrm{QFT}}\mathrm{e}^{\alpha _i\varphi }+U_i(\varphi )\right\}.$$
(94)
Here $`U_i(\varphi )`$ are vacuum energies on branes, normally $`U_i(\varphi )`$ are dictated by the form of dilatonic potential. For example, if $`V\mathrm{e}^{\kappa \varphi }`$ then $`U_i`$ has also the exponential form. $`L_{\mathrm{QFT}}`$ is an arbitrary Lagrangian corresponding to massless QFT (say, QED, QCD, SM, GUT) which is classically conformally invariant in the background $`\stackrel{~}{g}_{\mu \nu }`$. There is dilaton coupling on the brane which is typical for Brans-Dicke gravity. Usually it is assumed to be invisible in 4d world.
We are going to search for the solutions of the sort
$$ds^2=dz^2+\mathrm{e}^{2A(z)}\stackrel{~}{g}_{ij}dx^idx^j$$
(95)
where 4 dimensional $`\stackrel{~}{g}_{ij}=a^2(\eta )\eta _{ij}`$. It is assumed that branes sit on $`z=z_{\mathrm{hid}}`$ and $`z=z_{\mathrm{vis}}`$.
We integrate over quantum fields in theory (94). Supposing that there is only gravitational background, the interaction of dilaton coupled quantum fields leads to effective action induced by conformal anomaly :
$$\mathrm{\Gamma }_{\mathrm{source}}=\underset{i=\mathrm{hid},\mathrm{vis}}{}V_3𝑑\eta \left\{2b_1\sigma _1\sigma _1^{\prime \prime \prime \prime }2(b_1+b)\left(\sigma _1^{\prime \prime }\sigma _{1}^{}{}_{}{}^{2}\right)^2\right\}$$
(96)
where $`\sigma =\mathrm{ln}a(\eta )`$, $`\sigma _1=A+\sigma +\frac{\alpha _1\varphi }{3}`$. For the sake of simplicity, we adopt the large $`N`$-expansion (that justifies the neglecting of proper quantum gravity contribution to (96)). If spinors give the leading contribution then $`b=\frac{3N_{\frac{1}{2}}}{60(4\pi )^2}`$, $`b_1=\frac{11N_{\frac{1}{2}}}{360(4\pi )^2}`$. One can take the contribution above as corresponding to maximally SUSY Yang-Mills theory (which only changes the coefficients of above effective action). That corresponds to implementing above compactification to AdS/CFT scheme. Note that the suggestion to take into account the boundary matter quantum effects (via conformal anomaly induced effective action for SUSY Yang- Mills theory) in brane-world scenario has appeared in ref.. It was shown the possibility of creation of de Sitter or Anti-de Sitter 4d Universe in 5d AdS space. In ref. the same idea on application of conformal anomaly has been expressed and effective brane tension due to such boundary quantum contribution for 4d de Sitter world has been found.
Thus our complete action will be given by sum of three terms:
$$S_{\mathrm{complete}}=S+\mathrm{\Gamma }_{\mathrm{source}}+\underset{i=\mathrm{hid},\mathrm{vis}}{}V_3𝑑\eta \mathrm{e}^{4A}a^4(\eta )U_i(\varphi ).$$
(97)
One can now consider the solution of the equations of motion given from the action (97). In the bulk 5d universe, the action is identical with the previous one (1), then the solutions in the bulk are also given by the previous ones. Especially when $`V(\varphi )`$ is a constant, we obtain the solution in (5). Near the brane, however, one obtains the following equations for $`d=D1=4`$ instead of (9) :
$`\varphi ^{\prime \prime }`$ $``$ $`\left[{\displaystyle \frac{U_i(\varphi )}{\varphi }}+{\displaystyle \frac{\alpha _1}{3}}\mathrm{e}^{4A}\left\{4b_1\sigma _1^{\prime \prime \prime \prime }4(b+b_1)(\sigma _1^{\prime \prime \prime \prime }6\sigma _{1}^{}{}_{}{}^{2}\sigma _1^{\prime \prime })\right\}\right]\delta (zz_i),`$
$`A^{\prime \prime }`$ $``$ $`{\displaystyle \frac{1}{6}}\left\{U_i(\varphi )4b_1\sigma _1^{\prime \prime \prime \prime }4(b+b_1)(\sigma _1^{\prime \prime \prime \prime }6\sigma _{1}^{}{}_{}{}^{2}\sigma _1^{\prime \prime })\right\}\delta (zz_i).`$ (98)
Then by substituting the solution in (5), (86), we find $`U_i`$ and $`\frac{U_i}{\varphi }`$ become time dependent:
$`U_{\mathrm{hid}}`$ $`=`$ $`4\lambda \sqrt{{\displaystyle \frac{3}{4}}\left(1+{\displaystyle \frac{c^2}{2\lambda ^2y_{\mathrm{hid}}^4}}+{\displaystyle \frac{4k}{\lambda ^2y_{\mathrm{hid}}}}\right)}{\displaystyle \frac{6b_1y_{\mathrm{hid}}^2}{t^4}}`$
$`U_{\mathrm{vis}}`$ $`=`$ $`4\lambda \sqrt{{\displaystyle \frac{3}{4}}\left(1+{\displaystyle \frac{c^2}{2\lambda ^2y_{\mathrm{vis}}^4}}+{\displaystyle \frac{4k}{\lambda ^2y_{\mathrm{vis}}}}\right)}{\displaystyle \frac{6b_1y_{\mathrm{vis}}^2}{t^4}}`$
$`{\displaystyle \frac{U_{\mathrm{hid}}}{\varphi }}`$ $`=`$ $`{\displaystyle \frac{2c}{y_{\mathrm{hid}}^{\frac{d}{2}}}}{\displaystyle \frac{8\alpha _1b_1y_{\mathrm{hid}}^2}{t^4}}`$
$`{\displaystyle \frac{U_{\mathrm{vis}}}{\varphi }}`$ $`=`$ $`{\displaystyle \frac{2c}{y_{\mathrm{vis}}^{\frac{d}{2}}}}{\displaystyle \frac{8\alpha _1b_1y_{\mathrm{vis}}^2}{t^4}}.`$ (99)
Hence, the price one paids for keeping the same de Sitter brane-world solution is in change of brane vacuum energies. Quantum corrections explicitly give contribution to vacuum energies which become time-dependent (or dependent from the radius of de Sitter space in radial coordinates). That indicates that effective brane tension will be changed.
Let us discuss now RG flow of 4d Newton constant (for a recent review of holographic RG, see ). Using (22) for the solution in (5) for $`d=4`$, we find
$$\frac{1}{\stackrel{~}{G}_4}=\frac{1}{G_5}_{y_{\mathrm{hid}}}^{y_{\mathrm{vis}}}𝑑y\frac{\sqrt{3}}{\lambda \sqrt{1+\frac{c^2}{2\lambda ^2y^4}+\frac{kd}{\lambda ^2y}}}.$$
(100)
If we define $`U`$ by $`y_{\mathrm{vis}}=U^2`$, we can identify $`U`$ with the energy scale on the visible brane from AdS/CFT correspondence . Therefore Eq.(100) expresses the scale dependence of the gravitational coupling. When $`U`$ ($`y_{\mathrm{vis}}`$) is small, we find
$$\frac{1}{\stackrel{~}{G}_4}\frac{1}{G_5}\frac{\left(y_{\mathrm{vis}}^3y_{\mathrm{hid}}^3\right)}{c\sqrt{3}}=\frac{1}{G_5}\frac{\left(U^6y_{\mathrm{hid}}^3\right)}{c\sqrt{3}}.$$
(101)
Therefore the gravitational coupling $`\stackrel{~}{G}_4`$ becomes large in the IR region, which would be due to the curvature singularity at $`y=0`$. On the other hand, when $`U`$ ($`y_{\mathrm{vis}}`$) is large, we find
$$\frac{1}{\stackrel{~}{G}_4}\frac{1}{G_5}\frac{\sqrt{3}}{\lambda }\left(y_{\mathrm{vis}}+Y(y_{\mathrm{hid}})\right)=\frac{1}{G_5}\frac{\sqrt{3}}{\lambda }\left(U^2+Y(y_{\mathrm{hid}})\right).$$
(102)
Here $`Y`$ depends on $`y_{\mathrm{hid}}`$ but does not on $`y_{\mathrm{vis}}`$. Eq.(102) seems to tell that the gravitational coupling $`\stackrel{~}{G}_4`$ becomes small in the UV region.
We can also consider the scale dependence of the Newton constant for the black hole type solution in Section 2. If the scale $`U=\mathrm{e}^{A(z=z_{\mathrm{vis}})}`$ is small (IR region), $`z_{\mathrm{vis}}`$ becomes negative and large. Then from (24) and (22), we find
$`U`$ $``$ $`\mathrm{e}^{\gamma z_{\mathrm{vis}}},`$
$`{\displaystyle \frac{1}{\stackrel{~}{G}_4}}`$ $`=`$ $`{\displaystyle \frac{1}{G_5}}{\displaystyle ^{z_{\mathrm{vis}}}}𝑑z\mathrm{e}^{2A}{\displaystyle \frac{1}{G_5}}{\displaystyle ^{z_{\mathrm{vis}}}}𝑑z\mathrm{e}^{2\gamma z}={\displaystyle \frac{1}{G_5}}{\displaystyle \frac{\mathrm{e}^{2\gamma z_{\mathrm{vis}}}}{2\gamma }}={\displaystyle \frac{1}{G_5}}{\displaystyle \frac{U^2}{2\gamma }}.`$ (103)
Here we assumed $`z_{\mathrm{vis}}>z_{\mathrm{hid}}\mathrm{}`$ and put the constant of the integration to vanish. Therefore $`\stackrel{~}{G}_4`$ becomes large in the IR region. On the other hand, when the scale $`U=\mathrm{e}^{A(z=z_{\mathrm{vis}})}`$ is large (UV region), $`z_{\mathrm{vis}}`$ becomes positive and large. By using (26), we obtain
$`U`$ $``$ $`\alpha z_{\mathrm{vis}}^\beta ,`$
$`{\displaystyle \frac{1}{\stackrel{~}{G}_4}}`$ $``$ $`{\displaystyle \frac{1}{G_5}}{\displaystyle ^{z_{\mathrm{vis}}}}𝑑z\alpha ^2z^{2\beta }={\displaystyle \frac{1}{G_5}}\left({\displaystyle \frac{\alpha ^2z_{\mathrm{vis}}^{2\beta +1}}{2\beta +1}}+\stackrel{~}{Y}(y_{\mathrm{hid}})\right)`$ (104)
$`=`$ $`{\displaystyle \frac{1}{G_5}}\left({\displaystyle \frac{U^{2+\frac{1}{\beta }}}{(2\beta +1)\alpha ^{\frac{1}{\beta }}}}+\stackrel{~}{Y}(y_{\mathrm{hid}})\right).`$
Here $`\stackrel{~}{Y}`$ depends on $`y_{\mathrm{hid}}`$ but does not on $`y_{\mathrm{vis}}`$. For $`d=D1=4`$ and $`𝒩=1`$ model in , we have $`2+\frac{1}{\beta }=\frac{18}{5}`$. Therefore the gravitational coupling $`\stackrel{~}{G}_4`$ becomes small in the UV region again. We should note that the parameters specifying the black hole on the brane do not enter in the above expressions (6) and (6), which is significantly different from the case that there is a black hole in the bulk .
## 7 Discussion
In summary, we presented the family of brane-world solutions of 5d dilatonic gravity. This family includes flat brane with effectively small cosmological constant, (anti) de Sitter and Nariai spaces, and brane-world dilatonic black holes. The study of gravity perturbations around such black holes shows that 4d gravity may be trapped. Corrections to Newton law near branes are calculated. The proposal to take into account brane matter quantum effects is made.( Actually, such proposal presented earlier in refs. helps to formulate the problem in terms of AdS/CFT.) The corresponding anomaly induced effective action is used to estimate the role of quantum effects in realization of de Sitter branes in asymptotically AdS dilatonic space. It is demonstrated that quantum corrections change the brane vacuum energies. RG flow of Newton constant in IR and UV is discussed.
There are many related problems which are left for future study. In particular, it looks that the picture under discussion should be realized within AdS/CFT correspondence. As it has been partially demonstrated it is possible to do (at least in case of de Sitter brane). However, the details of such quantum corrected brane-world Universe should be investigated deeply. The research of the role of brane matter quantum effects to black holes is also extremely interesting topic. Another open problem is the correct interpretation of dilatonic brane-world black holes as holographic RG flows. The presented example of RG flow of 4d Newton constant for dilatonic black hole represents the modest step in this direction.
## 8 Acknoweledgements.
We thank M.Ryan for the interest in this work. This research has been supported in part by CONACyT grant 2845E.
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# Modal field theory and quasi-sparse eigenvector diagonalization
## 1 Introduction
Of the many approaches to non-perturbative quantum field theory, we can identify two general computational strategies. The first is the method of Monte Carlo. The main advantages of this approach is that it can treat many higher dimensional field theories, requires relatively little storage, and can be performed with massively parallel computers. The other strategy is the method of explicit diagonalization. The strong points of this method are that it is immune to the fermion sign problem, can handle complex-valued actions, and yields direct information about the spectrum and eigenstate wavefunctions.
Modal field theory is a simple Hamiltonian framework which can accomodate either computational strategy. The first step is to approximate field theory as a finite-dimensional quantum mechanical system. The approximation is generated by decomposing field configurations into free wave modes and has been explored using both spherical partial waves and periodic box modes . From there we can analyze the properties of the reduced system using Monte Carlo, diagonalization, or some other computational method. In this short review we present two different approaches using the modal field formalism which address the example of $`\varphi ^4`$ theory in $`1+1`$ dimensions. We begin with a summary of the Monte Carlo calculation of the critical behavior of $`\varphi ^4`$ theory in Euclidean space. We then discuss the computational challenges of a diagonalization-based approach and conclude with a calculation of the lowest energy states of the theory using a diagonalization technique known as the quasi-sparse eigenvector method.
## 2 Monte Carlo
In this section we review the basic features of modal field theory in a periodic box and discuss spontaneous symmetry breaking in Euclidean two-dimensional $`\varphi ^4`$ theory using the method of diffusion Monte Carlo. The field configuration $`\varphi `$ is subject to periodic boundary conditions $`\varphi (t,xL)=\varphi (t,x+L).`$ Expanding in terms of periodic box modes, we have
$$\varphi (t,x)=\sqrt{\frac{1}{2L}}\underset{n=0,\pm 1,\mathrm{}}{}\varphi _n(t)e^{in\pi x/L}.$$
(1)
The sum over momentum modes is regulated by choosing some large positive number $`N_{\mathrm{max}}`$ and throwing out all high-momentum modes $`\varphi _n`$ such that $`\left|n\right|>N_{\mathrm{max}}`$. In this theory renormalization can be implemented by normal ordering the $`\varphi ^4`$ interaction term. After a straightforward calculation (details shown in ), we find that the counterterm Hamiltonian has the form
$$\frac{6\lambda b}{4!2L}\underset{n=N_{\mathrm{max}},\mathrm{}N_{\mathrm{max}}}{}\varphi _n\varphi _n,$$
(2)
where
$$b=\underset{n=N_{\mathrm{max}},\mathrm{}N_{\mathrm{max}}}{}\frac{1}{2\omega _n},\omega _n=\sqrt{\frac{n^2\pi ^2}{L^2}+\mu ^2}.$$
(3)
We represent the canonical conjugate pair $`\varphi _n`$ and $`\frac{d\varphi _n}{dt}`$ using the Schrödinger operators $`q_n`$ and $`i\frac{}{q_n}`$. Then the functional integral for $`\varphi ^4`$ theory is equivalent to that for a quantum mechanical system with Hamiltonian
$`H`$ $`={\displaystyle \underset{n=N_{\mathrm{max}},\mathrm{}N_{\mathrm{max}}}{}}\left[\frac{1}{2}\frac{}{q_n}\frac{}{q_n}+\frac{1}{2}(\omega _n^2\frac{\lambda b}{4L})q_nq_n\right]`$ (4)
$`+\frac{\lambda }{4!2L}{\displaystyle \underset{{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{n_i=N_{\mathrm{max}},\mathrm{}N_{\mathrm{max}}}{n_1+n_2+n_3+n_4=0}}}{}}q_{n_1}q_{n_2}q_{n_3}q_{n_4}.`$
The existence of a second-order phase transition in two-dimensional $`\varphi ^4`$ theory has been discussed in the literature . The phase transition is due to $`\varphi `$ developing a non-zero expectation value and the spontaneous breaking of $`\varphi \varphi `$ reflection symmetry. It is believed that this theory belongs to the same universality class as the two-dimensional Ising model and therefore shares the same critical exponents. From the Ising model we expect to find
$$\nu =1,\beta =\frac{1}{8}.$$
(5)
$`\nu `$ is the exponent associated with the inverse correlation length or, equivalently, the mass of the one-particle state, $`m`$. It is defined by the behavior of $`m`$ near the critical coupling $`\lambda _c,`$
$$m(\lambda _c\lambda )^\nu .$$
(6)
We determine the behavior of the mass as we approach the critical point from the symmetric phase of the theory. All computations are done using the method of diffusion Monte Carlo (DMC). The idea of DMC is to model the dynamics of the imaginary-time Schrödinger equation using the diffusion and decay/production of simulated particles. The kinetic energy term in the Hamiltonian determines the diffusion rate of the simulated particles and the potential energy term determines the local decay/production rate. A self-contained introduction to DMC can be found in .
$`\beta `$ is the critical exponent describing the behavior of the vacuum expectation value. In the symmetric phase the vacuum state is unique and invariant under the reflection transformation $`\varphi \varphi `$ (or equivalently $`q_nq_n,`$ for each $`n`$). In the broken-symmetry phase the vacuum is degenerate as $`L\mathrm{}`$, and $`q_0`$, the zero-momentum mode, develops a vacuum expectation value. In the $`L\mathrm{}`$ limit tunnelling between vacuum states is forbidden. One ground state, $`|0^+,`$ is non-zero only for values $`q_0>0`$ and the other, $`|0^{}`$, is non-zero only for $`q_0<0`$. We will choose $`|0^+`$ and $`|0^{}`$ to be unit normalized. We can determine $`\beta `$ from the behavior of the vacuum expectation value as we approach the critical coupling,
$$0^+\left|\varphi \right|0^+(\lambda \lambda _c)^\beta .$$
(7)
In our calculations we used $`(L,N_{\mathrm{max}})=(2.5\pi ,`$ $`8),(2.5\pi ,`$ $`10),(5\pi ,`$ $`16),`$ and $`(5\pi ,`$ $`20).`$ For convenience we measure all quantities in units where $`\mu =1`$. For each set of parameters $`L`$ and $`N_{\mathrm{max}},`$ the curves for $`m`$ and $`0^+\left|\varphi \right|0^+`$ near the critical coupling were fitted using the parameterized forms
$$m=a\left(\frac{\lambda _c^m}{4!}\frac{\lambda }{4!}\right)^\nu $$
(8)
and
$$0^+\left|\varphi \right|0^+=b\left(\frac{\lambda }{4!}\frac{\lambda _c^\varphi }{4!}\right)^\beta .$$
(9)
Combining the results for the various values of $`L`$ and $`N_{\mathrm{max}}`$ and taking into account finite-size effects, we can extrapolate to the $`L`$, $`N_{\mathrm{max}}\mathrm{}`$ limit. The results we find are
$`\frac{\lambda _c^m}{4!}`$ $`=2.5\pm 0.2\pm 0.1,\nu =1.3\pm 0.2\pm 0.1,a=0.43\pm 0.05\pm 0.02,`$ (10)
$`\frac{\lambda _c^\varphi }{4!}`$ $`=2.5\pm 0.1\pm 0.1,\beta =0.13\pm 0.02\pm 0.01,b=0.71\pm 0.04\pm 0.03.`$
The first error bounds include inaccuracies due to Monte Carlo statistics, higher energy states (for the mass curves), and extrapolation. The second error bounds represent estimates of the systematic errors due to our choice of initial state, time step parameter, and bin sizes in the DMC simulations. Our results for the critical exponents are consistent with the Ising model predictions (5). The results for the critical coupling $`\frac{\lambda _c^m}{4!}`$ and $`\frac{\lambda _c^\varphi }{4!}`$ are in agreement with a recent lattice computation
$$\frac{\lambda _c}{4!}=2.56_{.01}^{+.02}$$
(11)
as well as the discrete light-cone quantization result
$$\frac{\lambda _c}{4!}\frac{(3+\sqrt{3})\pi }{6}2.48.$$
(12)
## 3 Diagonalization
In principle diagonalization techniques can provide significant detailed information about a quantum system. The basic idea is to diagonalize the Hamiltonian matrix defined over a truncated finite-dimensional Fock space. In practise however few field theories are computationally feasible and typically only theories in $`1+1`$ dimensions. Sparse matrix techniques such as the Lanczos or Arnoldi schemes allow us to push the dimension of Fock space to about 10<sup>5</sup> or 10<sup>6</sup> states. This may be sufficient to do accurate calculations for $`\varphi _{1+1}^4`$ theory near the critical point $`\frac{\lambda }{4!}2.5`$ for larger values of $`L`$ and $`N_{\mathrm{max}}.`$ It is, however, near the upper limit of what is possible using current computer technology and existing algorithms. Unfortunately field theories in $`2+1`$ and $`3+1`$ dimensions will require much larger Fock spaces, probably at least 10<sup>12</sup> and 10<sup>18</sup> states respectively. In order to tackle these larger Fock spaces it is necessary to venture beyond standard diagonalization approaches.
In view of this difficulty one could ask how the Monte Carlo method avoids this problem. The answer is by means of importance sampling. While the space of all configurations is impossibly large, it is often adequate to sample only the most important configurations. Unfortunately Monte Carlo techniques are not effective for a number of important problems. These include systems with dynamical fermions and computations involving complex-valued actions. In general difficulties arise when the functional integral measure is not positive definite and produce contributions with oscillating signs or phases.
In a new importance sampling method was proposed for use within a diagonalization scheme. The aim was to exploit the extreme sparsity of the Hamiltonian matrix, a result of the restricted form that characterizes local renormalizable interactions. If $`N`$ is the number of non-zero entries per row or column in the Hamiltonian matrix, it is shown in that a typical eigenvector is dominated by its largest $`\sqrt{N}`$ elements. There are some exceptions to this rule related to degeneracies which are discussed in . The algorithm for diagonalizing $`H`$ is as follows. We start by choosing a set of orthonormal basis vectors such that the Hamiltonian matrix $`H_{ij}`$ is sparse and the eigenvectors are quasi-sparse. The term quasi-sparse indicates that the norm of a vector is dominated by the contribution from a small fraction of its components. The remaining steps are as follows:
1. Select a subset of basis vectors $`\{e_{i_1},\mathrm{},e_{i_n}\}`$ and call the corresponding subspace $`S`$.
2. Diagonalize $`H`$ restricted to $`S`$ and find one eigenvector $`v.`$
3. Sort the basis components $`e_{i_j}|v`$ according to absolute size and throw away the least important basis vectors.
4. Replace the discarded basis vectors by new basis vectors. These are selected at random from a pool of candidate basis vectors which are connected to the old basis vectors through non-vanishing matrix elements of $`H`$.
5. Redefine $`S`$ as the subspace spanned by the updated set of basis vectors and repeat steps 2 through 5.
If the subset of basis vectors is sufficiently large, the exact eigenvectors will be stable fixed points of the update process. We will refer to this diagonalization technique as the quasi-sparse eigenvector (QSE) method.
We now test the QSE method on $`\varphi ^4`$ theory in $`1+1`$ dimensions. The theory can be treated rather well using conventional sparse matrix techniques and a Fock space of about 10<sup>6</sup> states. However the test we propose will be more stringent. The size of Fock space and the number of non-zero transitions per state for $`\varphi ^4`$ theory in $`3+1`$ dimensions are the cubes of the corresponding numbers in $`1+1`$ dimensions. Therefore a more useful test is to solve the $`1+1`$ dimensional system in a manner that can be directly generalized and expanded to the $`3+1`$ dimensional problem. Setting an upper limit of $`10^6`$ states for the $`3+1`$ dimensional case, we will attempt to solve the $`1+1`$ dimensional system using no more than 10$`0`$ (cube root of 10<sup>6</sup>) basis vectors at a time.
We again use the Hamiltonian in (4) and set $`L=5\pi `$ and $`N_{\mathrm{max}}=20`$. This corresponds with a momentum cutoff scale of $`\mathrm{\Lambda }=4.`$ We also place auxiliary constraints on the states in our momentum Fock space. We keep only those states with $`13`$ particles and kinetic energy $`2\mathrm{\Lambda }`$. A precise definition of the kinetic energy cutoff is provided in , and in terms of the parameter $`K_{\mathrm{max}}`$ introduced there we use $`K_{\mathrm{max}}=41`$. In our calculations we also restrict our attention to the zero momentum sector.
With these constraints our Fock space contains about $`2\times 10^6`$ states and the Hamiltonian matrix has about $`10^3`$ transitions per state. However, we will restrict our QSE calculation to include only 100 basis vectors at a time. Results for the energy eigenvalues are shown in Fig. 1. From our Monte Carlo calculation we know that the theory has a phase transition at $`\frac{\lambda }{4!}2.5`$ corresponding with spontaneous breaking of the $`\varphi \varphi `$ reflection symmetry. In the broken phase there are two degenerate ground states and we refer to these as the even and odd vacuum states. In Fig. 1 we see signs of a second order phase transition near $`\frac{\lambda }{4!}2.5`$. Since we are working in a finite volume the spectrum is discrete, and we can track the energy eigenvalues as functions of the coupling. Crossing the phase boundary, we see that the vacuum in the symmetric phase becomes the even vacuum in the broken phase while the one-particle state in the symmetric phase becomes the odd vacuum. The energy difference between the states is also in agreement with the Monte Carlo calculation of the same quantities. The state marking the two-particle threshold in the symmetric phase becomes the one-particle state above the odd vacuum, while the state at the three-particle threshold becomes the one-particle state above the even vacuum. These one-particle states should be degenerate in the infinite volume limit. One rather unusual feature is the behavior of the first two-particle state above threshold in the symmetric phase. In the symmetric phase this state lies close to the two-particle threshold. But as we cross the phase boundary the state which was the two-particle threshold is changed into a one-particle state. Thus our two-particle state is pushed up even further to become a two-particle state above the even vacuum and we see a pronounced level crossing.
A conventional sparse matrix algorithm would require about $`10`$ GB RAM and 10<sup>5</sup> Gflops to determine the energy and wavefunctions of the lowest lying states for this system. The QSE method was able to perform the same task with about $`30`$ KB RAM and 10<sup>2</sup> Gflops. The reductions in floating point operations and memory are about three and six orders of magnitude respectively and show the advantages of QSE diagonalization over other numerical techniques. Other large non-perturbative systems should now also be amenable to numerical solution using this method.
## Acknowledgments
I thank my collaborators on the works cited here and the organizers and participants of the Workshop on Light-Cone QCD and Nonperturbative Hadron Physics in Adelaide. Financial support provided by the National Science Foundation.
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# 1. Some properties of inherent oscillations and waves
## 1. Some properties of inherent oscillations and waves
The traditional models of molecules and crystals in which the atoms are replaced by physical points with masses of atoms do not correspond to adiabatic Born-Oppenheimer’s theory and do not give opportunities to study inherent oscillations. The opportunity of inherent oscillations and waves existence in the condensed materials can be proved by using of adiabatic model of solids . In such model each atom is submitted as a nucleus and electron shell connected with each other by elastic force and capable to carry out oscillatory movements ($`\alpha `$ type of oscillations). Each atom in adiabatic model of a crystal can be submitted as inherent $`\alpha `$ oscillator consisting from a nucleus and electron shell. Elastic force of inherent oscillator is almost proportional to displacement of a nucleus from the centre of an electron shell. The value of appropriate factor depend on electrical potential field which is created by electron shell near her centre. The frequency of inherent oscillations depends on conditions in which there is an atom. In particular the collective properties of electrons in a crystal can be expressed so that their displacements (displacements of electron shells) are coherent (synchronous), within the some length ($`\lambda `$) and the coherent area acquire superconductive properties. The size $`\lambda `$ can exceed tens of space lattice and to achieve significant size. In such conditions the nucleus of each atom has an opportunity to oscillate relatively common mass of coherent area. The amplitude of a nucleus oscillations is insignificant $`(10^2A^0)`$.
The energy spectrum of ($`\alpha `$type) inherent oscillations in harmonic approach can be presented by the expression
$$E_\upsilon =E_\alpha (1/2+\upsilon )$$
(1)
where oscillatory quantum number $`\upsilon `$ = 0, 1, 2…, the elementary oscillatory quantum is defined by mass of a nucleus $`M_z`$ with number Z: $`E_\alpha =\mathrm{}(\zeta /M_z)^{1/2}`$, $`\zeta `$\- factor connecting force with displacement of a nucleus from an equilibrim position at the centre of an electronic envelop, $`\mathrm{}`$=h/$`(2\pi )`$, h - Plank’s constant. According to accounts (conformable to experimental data) at increasing Z magnitude $`E_\alpha `$ decreases from 0.402 eV at Z = 2 up to 0.22 eV at Z = 8 and then increase up to 0.406 eV at Z = 80. The inherent oscillations of a $`\beta `$\- type caused by joint displacements of nucleus and K - electron subshell relatively other part of an electron shell and also inherent oscillations of a $`\gamma `$\- type caused by joint displacement of a nucleus, K - and L - electron subshells relatively other part of an electron shell of atom are possible also. The elementary quantums of inherent oscillations $`\beta `$\- and $`\gamma `$\- types are less than the appropriate quantums of $`\alpha `$\- oscillations: $`E_\alpha >E_\beta >E_\gamma `$. Thus energy of some inherent oscillations with small $`\upsilon `$ in a crystal can not exceed energies of the solved electron transitions between various electronic states i and j: $`E_\upsilon <|E_iE_j|`$. This inequality is a condition of applicability an adiabatic principle in correspondence with . Thus in adiabatic approach in the condensed materials the inherent oscillations $`\alpha `$-, $`\beta `$\- and $`\gamma `$\- types can exist and propagate. They represents the oscillatory displacement of nucleuses relatively electronic environments. The speeds of inherent oscillations waves can achieve speed of a sound.
The theoretical opportunity of existence of inherent oscillations and waves followed from adiabatic model of a crystal, based on adiabatic Born-Oppenheimer’s principle, is confirmed by experimental data about temperature dependence of crystals resistivity, about phonon drag at high temperatures and about optical electron-vibrational spectra formed with participation of inherent oscillations. The energy levels of inherent oscillations are shown as deep energy levels and the electron transitions between these levels or on these levels or from these levels are electron-vibrational. Electron-vibrational transitions are accompanied by absorption or emission several (on the average S) crystal phonons, where S - constant of electron-phonon interaction.
The following ways of excitation of inherent oscillations can be practically important:
\- thermal, at the expense of phonons energy;
\- optical, at the expense of absorbed photons energy;
\- shock , for the account of kinetic energy of (hot) electrons or holes;
\- recombination , at the expense of energy recombination ;
\- combined, as a combination of two or several specified ways.
In semiconductors it is convenient to raise inherent oscillations at the expense of recombination energy of electrons (holes). For this purpose is necessary rather strong interaction of electrons with lattice oscillations. The constant S equal to an average number of phonons participating in the act of electron-phonon interaction and serves a measure of electron interactions with phonons. The size S in ideal silicon monocrystal calculated on the data on meaning of deformation potential constants does not exceed 0.03 for optical phonons and 0.15 for acoustic phonons . Hence excitation of inherent oscillations at the expense of interaction with electrons in ideal crystals is rare event that are in agreement with adiabatic theory according to which electrons and the oscillations of a lattice in an ideal crystal do not interact with each other. The effective production of inherent oscillations at the expense of recombination energy is possible in crystals containing such defects of a lattice for which the strong interaction of electrons with phonons is characteristic. These defects have received the name of the electron-vibrational centres (EVC). EVC are the Jan-Teller’s centres. They represents local defects of a crystal which equilibrum position or frequencies of oscillations depends on their electronic condition. For EVC the large meanings S ($`S>1`$) are characteristic. According to the theory - the quantity of S is close to an average phonons participating in electron-vibrational transition to energy levels of EVC. By the theoretical estimations S can exceed 50. The inherent oscillations of atoms of the basic substance can exist and to spread in crystals and crystal structures but effectively to raise such fluctuations and waves at the expense of electrons or holes energy is possible at presence EVC. In this connection the inherent oscillations and waves caused by atoms of the basic substance and also inherent oscillations and waves caused by impurity atoms are possible in crystals.
## 2. Impurity oxygen atoms in silicon
Impurity oxygen atoms in silicon monocrystals are electrical inactive and their presence can be established on characteristic optical absorption on length of a wave about 9 micrometer -.The intensity of the specified band of absorption reflects the contents of oxygen atoms on an optical path in a crystal. The increase of concentration of oxygen in silicon is accompanied by occurrence of quartz disseminations and by displacements of a maximum of the characteristic optical band of absorption to 10…11micrometers . This absorption is characteristic for Si oxides and it is explained traditionally in the literature by optical excitation of chains Si-O-Si oscillations. In the monography (p. 179) was shown convincingly that the reason of the specified optical band which determine a kind of a spectrum in the field of 9…11 micrometers consist not simply in specified chains of atoms but is contained in structure of a lattice (probably in structural defects containing atoms of oxygen). The similar bands of optical absorption in the specified spectral area are characteristic for oxides of InAs, InSb, InP and oxides of other semiconductors that also allows to connect these bands with defects of a lattice containing atoms of oxygen. Moreover there are weighty bases to believe that these defects containing atoms of oxygen are the Jan-Teller’s centres and they are EVC.
The presence of EVC inherent oscillators determine an optical properties of crystals. The oscillatory model describing optical spectra of crystals and supposing interaction of oscillators with a wide set of crystal phonons is stated in . In accordance to the optical reflectivity
$$R_{\mathrm{}}=\frac{(n1)^2+k^2}{(n+1)^2+k^2},$$
(2)
where n - optical factor of refraction and k - parameter of absorption are defined from the following expressions:
$$n^2k^2=\epsilon _{opt}+\omega _p^2\frac{\mathrm{\Omega }^2\omega ^2}{(\mathrm{\Omega }^2\omega ^2)^2+\omega ^2\nu ^2};$$
(3)
$$2nk=\omega _p^2\frac{\omega \nu }{(\mathrm{\Omega }^2\omega ^2)^2+\omega ^2\nu ^2},$$
(4)
where: $`\mathrm{\Omega }`$ \- frequency of oscillator, $`\nu `$ \- frequency of phonon, $`\omega `$ \- optical frequency, $`\omega _p=\sqrt{Ne^2/(M\epsilon _0)}`$, N - concentration of oscillators, e - charge of oscillator, M - mass of oscillator, $`\epsilon _0`$ \- electrical constant, $`\epsilon _{opt}`$ \- optical dielectric permeability. The optical reflection spectra of silicon oxides and reflection spectra of specified semiconductors oxides can be presented as sum of separate oscillators reflection spectra. Such decomposition of a spectrum corresponds to the quantum theory . These theories can be applied to the analysis of IR reflection by inherent oscillators of EVC.
Every separate oscillator have ”zero” frequency of oxygen inherent $`\alpha `$ \- oscillators ($`\mathrm{}\mathrm{\Omega }`$ = 0.11 eV) which cooperate mainly with various lattice frequencies between $`\mathrm{\Omega }`$ and $`\omega _p`$. The greatest contribution in oxydes reflection spectra is brought by frequencies appropriate to ”zero” fluctuations of inherent oxygen oscillators and phonons for which ($`\mathrm{\Omega }/\omega _p`$) = 0.25, $`\epsilon _{opt}`$=1.2, G/$`\mathrm{\Omega }`$=0.011, G - factor of attenuation. The rather small meaning $`\epsilon _{opt}`$ allows to carry him only to the local centre but not to all crystal. The experimental spectrum of polarized light reflection and calculated on the Eq. (2-4) spectrum of reflection are given on fig. 1. The experimental spectrum was measured at a corner of fall of the linearly polarized radiation of 45 degrees for oriented along an axis C quartz monocrystal. The comparison of spectra submitted in a fig. 1 specifies presence of oscillators in quartz with oscillatory energy which is equal to energy of ”zero” oscillations of inherent oxygen oscillators, calculated on the Eq. (1) at $`\upsilon `$ = 0 for atom of oxygen. Such transitions with frequency of ”zero” oscillations are forbidden for quantum oscillator but EVC are capable to show duality of properties . The conditions in a minimum of oscillator potential and transitions with frequency of ”zero” oscillations for EVC are accordingly allowable. The excess of experimental absorption above calculated absorption at energies of quantums more than 0.11 eV is determined by the contribution with participation of frequencies $`\omega _p`$ conterminous with frequencies of inherent $`\alpha `$\- oxygen oscillators calculated on the Eq. (1) at other meanings of $`\upsilon `$. The same conclusion can be made as a result of similar approach of reflection spectra for melted quartz, silicate glasses and oxides of a number of semiconductors. According to this result the optical band which is characteristic for quartz and others oxides of semiconductors is possible to identify with excitation of ”zero” inherent oscillations of oxygen atom ($`\upsilon `$ = 0).
Impurity atoms of oxygen in silicon irradiated by radiation formes associations with vacancies of a crystal lattice known as A-centres. As a result of study of electron spin resonance spectra and his dependence from orientation data and about energy levels of defects, items of information on a situation of oxygen atoms in silicon lattice in the model of A-centre was constructed. A-centre represents anisotropic defect whose electrical dipole moment is directed lengthways . A-centre have deep electron-vibrational energy levels in the forbidden energy zone of silicon which are distant from boundaries ($`E_c`$ and $`E_v`$) of the forbidden energy zone on sizes multiple 0.11 eV and are close to calculated on the Eq. (1) for oxygen atom. It allows to identify them with inherent $`\alpha `$-oscillations of oxygen atoms. Acceptor $`E_c`$ \- (0.16…0.22) eV and donor $`E_v`$ \+ (0.27…0.33) eV levels are most active at low concentration of A - centres. These levels are close to the calculated energy levels of inherent oscillations of oxygen atom. Specified changes of the energy levels are explained by change of complete oscillatory energy of A - centre caused by the interaction of the centres with each other at changing of their concentration from $`510^{17}cm^3`$ up to $`10^{13}cm^3`$ and depends on average distance (R) between them as $`R^3`$. It coordinates with aeolotropic structure of A - centre and specifies one-dimensional character of interaction between them. The given result corresponds to oscillatory model of A-centre as one-dimensional oscillator that justifies application Eq. (1) for the description of his electron-vibrational energy levels.
The application of the linearly polarized radiation and monocrystals with rather low concentration of EVC allows to create necessary conditions for measurement of electron-vibrational spectra with participation of one phonons type. Such spectra may be analysed on the basis of the Pekar-Huang- Rhys theory \- . The typical spectra of photoconductivity and optical transmittance connected with photoionization of A - centre by polarized IR radiation with an electrical vector E $``$ are given on a fig. 2. One can see on fig. 2 that spectra are modulated by phonons. The minima of optical transmittance (maxima of optical absorption) correspond to minima of photoconductivity. The established conformity of extremums in spectra of photoconductivity and absorption is explained by characteristic negative photoconductivity when electron-vibrational transitions (from valence band on A - centre) occur and proves the fact of auto localization of electrons and holes on EVC simultaneously. If EVC interacts only with one type ocillations of a lattice (similar to spectra on a fig. 2) then in according to the spectral distribution of electron-vibrational transitions depends from S, function Bose ($`f`$), contains product of the modified Bessel’s function of the order p on $`_n\delta `$(n-p) , where n = 1, 2,… and $`\delta `$ \- delta function, p - number of phonons participating in electron-vibrational transition. Because of presence $`\delta `$ \- functions p accepts only integer meaning. p$`>`$0 corresponds to absorption of optical quantum and p$`<`$0 corresponds to radiation of optical quantum with emission of p phonons. Accordingly spectrum contains two wings adequate electron-vibrational transitions with absorption and emission of phonons. One wing lays below than energy of transition without phonons (p$`=`$0) another wing lays higher than energy of transition without phonons. The wing of a spectrum appropriate to absorption phonons disappears with downturn of temperature . At participation of one type phonons the spectrum represents a series of discrete lines which energy differ on size multiple phonons. Expression for spectral distribution includes only one parameter S. Size of S may be selected to approach conveniently a contour of an experimental spectrum and thus to determine S. The number of phonons which are emitted by the centre in a maximum of an optical band is equal S. The energy level on which the electron-vibrational transition is carried out corresponds to transition without phonons (p = 0). In the field of high temperatures when $`4S[f(f+1)]^{1/2}>p>1`$ the discrete lines extend the periodic structure in spectra disappears and the contour of a spectral band can be approached by function of Gauss:
$$G(\mathrm{}\omega )=exp\left\{\frac{1}{4S[f(f+1)]}(\frac{\omega \omega _{max}}{\omega _0})^2\right\}$$
(5)
where $`\omega `$\- optical frequency, $`\omega _{max}`$ \- frequency in a maximum of spectral distribution, $`\omega _0`$ \- optical frequency of electronic transition at = 0. In the region of low temperatures when f$``$0 and $`>>`$ Sf(f+1) contour of a spectral band follow the dependence
$$G(\mathrm{}\omega )=S^p/p!.$$
(6)
Approach to spectra on basis of expressions (5, 6) has allowed us to determine S = 5 for A - centre, types and energy of phonons connected with the centres, energy levels of A - centre corresponding to transitions with p = 0. In particular, data about TA phonons received from the analysis of experimental spectra connected with A- centre are given in Table 1 together with the appropriate literary data.
Table 1: Energy of TA phonons in silicon
| Direction | Energy of phonons (meV) determined by methods: | | | |
| --- | --- | --- | --- | --- |
| of phonon | calculating | dissipation | indirect | photocon- |
| wave | | of nutrons | IR | ductivity |
| vector | | | absorbtion | on A-centre |
| 100 | 23.0 | 21.0 | 22.0 | 22.0 |
| 110 | 18.0 | 17.9 | 18.0 | 18.3 |
| 111 | 16.5 | 16.7 | 17.0 | 16.9 |
On an insertion of a fig. 2 are given the experimental data about energy splitting of connected with A- centre LO and TA phonons arising as a result of interaction A- centres among themselves depending on their concentration (N). The size S also changes from 5 up to 2 at increase N up to $`510^{17}cm^3`$. Electrons and holes carrying out transitions to levels EVC inevitably cooperate with p=S phonons. In result at energy levels of EVC, which are inherent oscillatory energy levels of these centres and are shown as deep energy levels, appears located electrons, holes and phonons. Hence EVC represents complex formation consisting from impurity atom, his inherent oscillations , electrons, holes and phonons. These particles in structure EVC form the interconnected auto localized system submitting other statistical and dynamic laws versus free phonons and electrons of conductivity. In this connection EVC can cause unusual physical properties to crystals and crystal structures which are poorly connected to a structure of energy bands of a crystal and properties of conductivity electrons. The similar conclusion was made earlier as a result of research of phonon drag in semiconductors where the phonon drag is caused only by those phonons which are strongly connected with electrons and consequently have other properties .
The modulation of spectra by phonons stops in samples on the basis of various semiconductors when the concentration of free charge carriers at room temperature (concentration of doping impurity) exceeds $`n_{max}`$ = $`210^{17}cm^3`$. It gives the basis to consider that concentration of electrons and holes which are auto located on EVC and capable to be superconducting does not exceed $`n_{max}`$.
A- centre arise at technological processings of silicon and silicon structures, in particular, in structures metal - oxide of silicon - silicon (MOS-structures). The presence of A- centres in structures is shown on characteristic negative photoconductivity in silicon under oxide, modulation of spectra by phonons and on kinetics of photo-emf which is described by not less than triad of time constant (t). These constants are mutually connected by proportions $`t_1:t_2:t_3=1!/S^1:2!/S^2:3!/S^3`$ at S = 2 also corresponds to probability of electron-vibrational transitions which are described by Eq. (6) at participation p = 1, 2 and 3 phonons. A - centres also influence on volt - capacitance characteristics of MOS-struktures. On the fig. 3 are submitted typical volt - capacitance (C-V) characteristics of aluminium contact by the area $`4.910^4cm^2`$ to a plate of silicon containing A- centres in the concentration $`510^{15}cm^3`$. The characteristics are measured at room temperatureT ($`T<T_c`$) on various frequencies. The submitted characteristics are depended from frequency. The dependence of the characteristics from frequency is defined by available A - centres in samples. On various frequencies the capacitance is nonmonotonic function of back bias that does not correspond to traditional theory of capacitance for ideal contact metal - semiconductor with Schottky barrier. Untraditional physical model of contact metal - semiconductor is necessary for the analysis of such C-V dependences. It is possible to approximate the experimental C-V curves submitted in a fig. 3 by known dependence of contact capacitance (C) from the enclosed voltage (V): $`C=[S(V)\epsilon \epsilon _0]/L`$, where $`\epsilon `$ \- relative dielectric permeability of the semiconductor and $`e_0`$ \- electrical constant, L - thickness of depletion region, accounting that the effective area of contact (S) depends on a voltage: S = S(V). Such dependences corresponds of contact model which is taking into account the presence of fine regions (embedments) with high onductivity in the semiconductor under a field electrode which are coherent areas. In process of penetration of an electrical field into the semiconductor (at increase of back-biasing voltage) the border of depletion region achieves some of these embedments and area of equipotential surface grows. Thus differential capacitance of contact is accordingly increased. On experience the smooth increase of capacitance is observed at achievement of some bias voltage that it is possible to explain by volumetric distribution of the embedments which serially achieves equipotential surface of depletion region border. The unmonotonous dependence C(V) with several extremums in such model can be explained by presence of several layers of embedments. Identical high-resistance and low-resistance layers, as it is well known, exists in such and similar structures strongly influences conductivity by elastic waves arising in depletion region of contact. Besides it is known that the periodicity of layers also depends on a voltage on contact and from voltage frequency. It brings contribution in frequency dependence of contact capacitance. The appropriate model of contact which allow to explain frequency dependence of capacitance contains coherence areas of the small size in the semiconductor under a field electrode. These areas are distributed disorderly and are grouped in layers parallel to a field electrode. It is possible to explain the experimental C-V curves just by a discrete structure of layers formated by fine coherence areas. Monolithic low-resistance layers are unsuitable because in this case the differential capacitance can not change in relation to capacitance of depletion region of the semiconductor. The presence of embedments with conductive (superconductive) pieces in the semiconductor under field contact quite corresponds to an opportunity of formation of coherence areas with the characteristic size $`\lambda `$. Experimentally observable frequency dependence of differential capacitance in semiconductor contacts confirms presence of coherence areas in samples with electron-vibrational centres.
## 3. Estimation of parameters
Superconductive property of a material are doubtlessly connected to occurrence of coherence areas which possess of superconductive properties due to coherence of electronic displacements. The law of occurrence of coherence areas is defined by concentration those EVC in which the inherent oscillations are exited and also by concentration of phonons. The interaction between EVC determines occurrence of coherence areas and is executed by means of an exchanging of the centres with each other by phonons. It is possible to assume that for an estimation of the minimal concentration ($`N_{min}`$) of EVC at which the formation of coherence areas is possible when such interaction is effective on distances between EVC about length of phonon wave ($`\mathrm{\Lambda }`$). The minimal concentration of exited EVC can be estimated as $`\mathrm{\Lambda }^3`$. Accepting the maximal speed of a sound in silicon $`V=9.1310^5`$ cm/c and phonon frequency $`F=1.2510^{10}c^1`$ we shall receive $`N_{min}=2.5610^{12}cm^3`$.
Normal-superconducting transition temperature ($`T_c`$) can be estimated proceeding from the following reasons. If to consider the intrinsic semiconductor material with effective mass for density of conditions for electrons ($`m_{nd}`$) and holes ($`m_{pd}`$) when at certain temperature ($`T`$) speed of thermal generation of electrons and holes (G) in individual volume and at time of life ($`\tau _i`$) according to the statistical theory define inherent concentration $`n_i`$ = $`\sqrt{N_cN_v}exp(\frac{E_g}{2kT})`$ = G$`\tau _i`$. $`N_c=2(\frac{2\pi m_{nd}kT}{h^2})^{3/2}`$ and $`N_v=2(\frac{2\pi m_{nd}kT}{h^2})^{3/2}`$ \- effective number of conditions in conduction energy band and in valence energy band, $`E_g`$ \- width of energy forbidden gape of the semiconductor. In a stationary condition for the same time $`\tau _i`$ carriers of charges recombinate with the same rate G. If recombination occurs on energy levels of EVC then each recombination act gives rise to EVC inherent oscillations with energy determined by Eq. (1) and else S phonons is rised. Let’s consider that it is enough for synchronization of electronic displacement in coherence areas of one phonon per one EVC. Therefore for formation coherence area is necessary not less $`n_=(N_{min}/S)f(E)`$ recombination acts in time $`\tau _i`$, where $`f(E)`$ \- the function Bose for given oscillation energy $`E(\mu )`$. It is possible to write down a condition for existence of the minimal concentration of exited EVC as $`n_i=n`$. By substituting in this equality the expressions for $`n_i`$ and $`n_k`$we shall receive a relation which is possible to consider as the equation for definition of temperature $`T=T_c`$:
$$\sqrt{N_cN_v}exp(\frac{E_g}{2kT})=\frac{N_{min}}{S(exp\frac{E(\mu )}{kT}1)},$$
(7)
where E($`\mu `$) - discrete oscillatory energy of EVC. It is necessary to tell more about possible meanings of E($`\mu `$). E($`\mu `$) can not exceed Eg/2 as in this case she will be transferred with a high probability to electrons and is spent on electronic transitions through the forbidden zone of the semiconductor. It is known EVC show quantum and classical properties : E($`\mu `$) can accept meanings conterminous with energies of charmonical quantum vibrator according to Eq. (1) and also meanings multiple $`E_\alpha `$. Therefore for $`\alpha `$-type inherent EVC oscillator E($`\mu `$) = $`(E_\alpha /2)\mu <Eg/2`$, $`\mu `$ = 1,2,3,…
Using known tabulared meanings of parameters and also S = 5, $`\upsilon `$=0 and $`E_\alpha `$ = 0.22 for A- centre in silicon by numerical way is received from (7) the meaning of the $`T_c`$=305K at $`\mu `$=1, $`T_c`$=235K at $`\mu `$=2 and $`T_c`$=162K at $`\mu `$=3 for silicon. The similar estimation for Ge with the oxygen centres gives meaning of the $`T_c`$=182K, $`T_c`$=130K and $`T_c`$=33K. According to this results the superconductivity in silicon with $`\mu =2`$ , $`E(\mu )=E_\alpha `$, for A- centres necessary to expect at temperatures above 235K and in germanium - at temperatures is higher 130K.
The estimation of critical density of current ($`J_k`$) can be executed proceeding from the data about possible concentration of superconducting charges (n) and about maximal their drift speed (v). The meaning of n can be in limits from n = $`n_k`$ up to $`n_{max}`$ = 2$`10^{17}cm^3`$ and the speed of drift is probably limited by dissipation on borders of coherence areas and can not exceed speed of a sound V. Believing $`J_k`$ = enV we come to the following estimation of critical density of a current in silicon at room temperature: $`0,55/cm^2<J_k<73A/cm^2`$.
For an estimation of critical magnetic intensity it is possible to involve the following experimental data received at constant density of a current in silicon samples in a magnetic field. Such data are given on fig. 4 and represent typical dependence of derivative of voltage V (V - on sample) with respect to magnetic induction (B), that is dV/dB from B when current density $`J`$ = const, measured in a normal condition of Si sample (at $`T<T_c`$) containing EVC with concentration $`510^{15}cm^{}3`$. The dependence submitted on fig. 4 contains periodic maxima and can be explained by influence of a magnetic field on charged oscillators as electron connected with EVC by means of crystal phonon. The movement equation of such oscillator can be written down as follows:
$$\frac{d^2X}{dt^2}+(\frac{e}{m^{}})\frac{dX}{dt}B\omega ^2X=0,$$
(8)
where X - generalized coordinate , e - charge and m\* - effective mass of electron, $`\omega `$ \- cyclic frequency of fluctuations, B - magnetic induction. Component containing $`\frac{dX}{dt}`$ takes into account action of Lorentz's force when $`\frac{dX}{dt}`$ and B is mutually orthogonal. The given equation (8) supposes oscillating solution only under condition: $`\left\{4\omega ^2[\frac{eB}{m^{}}]^2\right\}>0`$.
The oscillatory movements of oscillator are possible when $`B<\frac{2m^{}\omega }{e}`$ that is when B is not too great. Otherwise oscillations are suppressed by a magnetic field and are impossible.
Maxima of dependences submitted on fig. 4 are possible to identify with frequencies those phonons which provide a coherence condition. If to accept for silicon $`m^{}=0.98m_0`$, where $`m_0`$ \- mass at rest of electron, then the appropriate frequencies multiple $`2.110^{10}`$ $`c^1`$, lay in limits from $`2.110^{10}c^1`$ up to $`1.2610^{11}c^1`$ and correspond to acoustic fluctuations of a crystal. It is seen from these data that the frequency $`F=1.2510^{10}c^1`$ accepted for an estimation of the minimal concentration of A- centres ($`N_{min}`$) is less but close to $`2.110^{10}`$ $`c^1`$. Oscillations with the large frequencies are absent when B $`>`$ 1.3 tesla hence coherence area can not arise also superconductivity is impossible. Therefore it is possible to expect that in magnetic fields with B $`<`$ 1.3 tesla the superconductivity caused EVC is possible but at B $`>`$ 1.3 tesla she may be impossible.
## 4. Experimental results and discussion
The experimental temperature dependences of resistance for two Si samples are submitted on fig. 5. Within the separate temperature areas the resistance is described by some activation energies which equal to energies of phonons cooperating with inherent oscillations of oxygen atoms. One can see on fig. 5 when the temperature is raised up to $`T=T_c`$ the resistivity sharply decreases up to zero (superconductivity arises at $`T=T_c`$). The superconductivity exist up to rather high temperatures (probably up to temperature of a crystal melting and is higher). One can see on fig. 5 else that curve 1 contains two steps of resistance sharp reduction (at $`T`$ = 250K and at $`T`$ = $`T_c`$ = 309K) but curve 2 contains one such step (at $`T=T_c=389.6K`$). The similar characteristics of some samples can contain 3 or 4 such steps. The presence of such steps resistance reduction at approach of temperature to $`T_c`$ specifies discrete (quantum) character of superconductivity arising. The occurrence of such steps can be connected with processes of arising, reorganization and cooperation of coherence areas with each other. In some samples alongside with the specified steps of resistance reduction the steps of sharp resistance increase are observed also and in such case superconductivity does not arise that it is possible to explain by deformation or destruction of coherency areas. In any case the sharp changing of samples resistance at discrete temperatures can definitely be connected with changing of a role of various quantums of EVC inherent oscillations because at absence of EVC the sharp jumps of resistance and superconductivity is not observed.
Becouse the given superconductivity arises due to oscillations of a crystal with very high frequencies appropriate to hypersound it is expedient to name the given superconductivity, caused by the electron-vibration centres, as hyperconductivity to distinguish her mechanism from other mechanisms of superconductivity.
In silicon samples we observed hyperconductivity in magnetic fields with B up to 0.5 tesla that does not contradict an estimation of allowable intensity of a magnetic field. Hyperconductivity existed at density of a current up to 10 $`A/cm^2`$ that not contradicts estimations of critical density for a current density.
The occurrence hyperconductivity is connected to occurrence of thermopower (TEP) features. On fig.6 are given temperature dependences of TEP for the sample designated by 2 on fig. 5. Such alternating-sign of temperature dependence TEP is characteristic for normal-hyperconducting transition. It seems the behaviour TEP reflects complex dynamic processes of formation and interaction of coherence areas among themselves.
Occurrence of superconductivity is accompanied by huge increase of thermal conductivity. In accordance with experimental data the meaning of thermal conductivity in some samples increase more than $`10^5`$ times. Probably the hyperconductivity is accompanied by heat superconductivity.
Hyperconductivity arises in various crystals containing EVC. On fig. 7 the available experimental data about meanings of the $`T_c`$ in various hyperconductor samples are submitted. One can see on fig. 7 that normal-hyperconducting transition temperature $`T_c`$ shows the tendency to reduction at increasing of average size of the basic substance nuclear number ($`Z_{mid}`$). Straight linees a and b on fig. 7 restrict the area in which the experimental points are placed. One can see on fig. 7 that the experimental meanings $`T_c`$ for every material can change over a wide range between linees a and b. On the other hand situation of straight linees a and b differs on insignificant size of energy $`kT17`$ meV appropriate to acoustic phonons. It speaks about importance not only inherent oscillations of EVC but acoustic phonons simultaneously for formation of hyperconductivity in different materials. The minimal experimental meanings of $`T_c`$ for Si and Ge are in agreement with given estimation of $`T_c`$ when $`\mu `$=2 for $`\alpha `$-inherent oxygen oscillator. According to Eq. (7) the increase S corresponds to reduction $`T_c`$ and promotes hyperconductivity existence in wide area of temperatures. The meanings $`S1`$ are necessary for existence hyperconductivity from near to room temperatures up to higher temperatures that it is possible to provide by introducing EVC in crystals. The experimental data about meanings $`T_c`$ for semiconductors with small $`E_g`$ can be coordinated with Eq. (7) only by taking into account $`\beta `$-type or $`\gamma `$-type of inherent oscillations because solutions of Eq. (7) with participation of $`\alpha `$ \- inherent oscillations are absent.
## 5. Conclusion
The superconductivity caused by electron-vibrational centres - hyperconductivity can exist in the field of room temperatures and at higher temperatures in various materials in wide electrical current density and magnetic intensity intervals. That allows ones to hope on opportunity of hyperconductivity application in science and engineering.
## Figure captions
Fig. 1. Experimental spectrum IR reflection of polarized light from a surface of quartz (1). The calculated spectrum of reflection by "zero" inherent fluctuations of oxygen (2).
Fig. 2. Spectrum of photoconductivity (1) and spectrum of IR transmittance (2) aused by photoionization of A - centres in silicon at T = 80K. The arrows designate the appropriate each other minima of spectra and the number participating phonons is specified by figures. On an insertion the experimental data about splitting LO and TA phonons are given depending on concentration of A - centres.
Fig. 3. haracteristics volt - capacitance of contact Si - Al containing A - centres are measured at room temperature $`T<T_c`$ on various frequencies: 0.2 MHz (1); 0.5 MHz (2); 1.0 MHz (3); 5.0 MHz (4); 10.0 MHz (5); 20.0 MHz (6).
Fig. 4. Derivative of voltage on Si sample with respect to magnetic intensity (B) at room temperature $`T<T_c`$.
Fig. 5. Temperature dependences of resistance for two of Si samples containing various concentration of A - centres: $`N_1510^{14}cm^3`$ and $`N_2610^{15}cm^3`$.
Fig. 6. Temperature dependence of thermopower for the same Si sample whose characteristic is designated on fig. 5 by number 2.
Fig. 7. The experimental meanings of normal-hyprconducting transition temperatures ($`T_c`$) according to average meanings of nuclear number ($`Z_{mid}`$) in various samples containing electron-vibrational centres. Straight linees a and b limit area in which the experimental points lay.
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