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warning/0003/nucl-th0003058.html	ar5iv	text	# Preserving the gauge invariance of meson production currents in the presence of explicit final-state interactions ## I Introduction Gauge invariance is one of the central issues when attempting to describe how photons interact with hadronic systems. Concentrating on the simplest case—pion photoproduction with real or virtual photons, gauge invariance can easily be shown to follow if the $`\pi N`$ and $`\gamma N`$ problems are treated completely and consistently on an equal footing . In practice, however, one often needs to revert to some approximate treatment of one or more of the contributing reaction mechanisms and this usually leads to a violation of gauge invariance. To restore it, the neglected reaction mechanisms must be approximated by auxiliary currents constructed such that the gauge-invariance-violating contributions to the four-divergence of the total production amplitude are cancelled. Such a procedure cannot be unique, of course, since one may always add arbitrary transverse currents without affecting the four-divergence. At the tree-level, where one does not resolve the internal mechanisms entering the interaction current (cf. Fig. 1), various recipes exist to preserve gauge invariance. The simplest case concerns the choice of bare vertices with pseudovector coupling for the $`\pi NN`$ vertex, where the corresponding Kroll–Ruderman contact current follows from the minimal substitution procedure. The case of extended nucleons, whose internal structure is described in terms of (phenomenological) form factors, is treated in Refs. . In the present work, we want to go beyond the tree level and investigate how one can preserve gauge invariance if the internal structure of the interaction current is taken into account explicitly. The reaction mechanisms that enter the interaction current are summarized within the dashed box in Fig. 1. Specifically, we are interested in preserving gauge invariance in the explicit presence of hadronic final-state interactions. This is achieved by introducing auxiliary currents which cancel the gauge-invariance-violating contributions. In particular, we show how one may exploit the constraints following from the generalized Ward–Takahashi identities to construct these currents. The present discussion is restricted to nucleons and pions only to facilitate the presentation. We do not include here possible resonances or other transition mechanisms since their couplings to the electromagnetic field are transverse and have no bearing on the question of gauge invariance. None of these restrictions are essential, however, and one may easily adapt the present formalism to accommodate more complex situations. ## II Gauge Invariance The pion photoproduction current of the nucleon is shown in Fig. 1 . According to the diagrams in the first line of this figure, the total current $`M^\mu `$ may be broken up into four main contributions: The three Born terms due to the $`s`$-, $`u`$-, and $`t`$-channel currents stemming from the photon coupling to the three external legs of the $`\pi NN`$ vertex, and the interaction current $`M_{\mathrm{int}}^\mu `$ where the photon attaches itself to an internal leg of the $`\pi NN`$ vertex, i.e., $$M^\mu =M_s^\mu +M_u^\mu +M_t^\mu +M_{\mathrm{int}}^\mu .$$ (1) While the first three contributions are relatively straightforward, the last one—as it is shown in the last two lines of Fig. 1—explicitly involves the full complexity of the internal reaction dynamics of the underlying $`\pi N`$ scattering problem summarized in Fig. 2. For gauge invariance of the total current $`M^\mu `$ to hold true, its four-divergence must satisfy a generalized Ward–Takahashi identity $`k_\mu M^\mu `$ $`=`$ $`[F_s\tau ]S_{p+k}Q_iS_p^1+S_p^{}^1Q_fS_{p^{}k}[F_u\tau ]`$ (3) $`+\mathrm{\Delta }_q^1Q_\pi \mathrm{\Delta }_{qk}[F_t\tau ],`$ where $`p`$ and $`k`$ are the four-momenta of the incoming nucleon and photon, respectively, and $`p^{}`$ and $`q`$ are the four-momenta of the outgoing nucleon and pion, respectively, related by momentum conservation $`p^{}+q=p+k`$. $`S`$ and $`\mathrm{\Delta }`$ are the propagators of the nucleons and pions, respectively, with their subscripts denoting the available four-momentum for the corresponding hadron; $`Q_i`$, $`Q_f`$, and $`Q_\pi `$ are the initial and final nucleon and the pion charge operators, respectively. $`[F_x\tau ]`$ denotes the $`\pi NN`$ vertex (including coupling and isospin operators), with the subscript $`x`$ labeling the kinematic situation appropriate for the $`s`$-, $`u`$-, or $`t`$-channel diagrams appearing in Fig. 1. The vertex isospin $`\tau `$ does not commute with the charge operators and $`\tau Q_iQ_f\tau Q_\pi \tau =0`$ describes charge conservation at the vertex in a symbolic manner. Equation (3) is easily obtained from Eq. (1) upon using the Ward–Takahashi identities for the nucleon and pion currents, $`k_\mu J_\mathrm{n}^\mu `$ $`=`$ $`S_{p_\mathrm{n}+k}^1Q_\mathrm{n}Q_\mathrm{n}S_{p_\mathrm{n}}^1,`$ (5) $`k_\mu J_\pi ^\mu `$ $`=`$ $`\mathrm{\Delta }_{q_\pi +k}^1Q_\pi Q_\pi \mathrm{\Delta }_{q_\pi }^1,`$ (6) where $`p_\mathrm{n}`$ and $`q_\pi `$ are the respective initial hadron momenta of the electromagnetic vertices, and using the fact that for gauge invariance to be true the interaction current must obey $$k_\mu M_{\mathrm{int}}^\mu =[F_s\tau ]Q_i+Q_f[F_u\tau ]+Q_\pi [F_t\tau ].$$ (7) ### A Preserving gauge invariance As alluded to above, the preceding relations can easily be shown to be true if the $`\pi N`$ and $`\gamma N`$ problems are treated consistently on an equal footing . In practical applications, however, approximations are inevitable which usually violate the gauge invariance. To see how one may preserve gauge invariance in such a situation, let us explicitly write the four-divergence of the interaction current using the relevant parts of Fig. 1 as guidance. One finds $`k_\mu M_{\mathrm{int}}^\mu `$ $`=`$ $`k_\mu m_{\mathrm{bare}}^\mu +k_\mu U^\mu G_0[F\tau ]`$ (10) $`+XG_0\{k_\mu J_\mathrm{n}^\mu S_{p_\mathrm{n}k}[F_u\tau ]+k_\mu J_\pi ^\mu \mathrm{\Delta }_{q_\pi k}[F_t\tau ]`$ $`+k_\mu m_{\mathrm{bare}}^\mu +k_\mu U^\mu G_0[F\tau ]\},`$ where $`G_0=S_{p_\mathrm{n}}\mathrm{\Delta }_{q_\pi }`$ denotes the product of the intermediate nucleon and pion propagators, with respective four-momenta $`p_\mathrm{n}`$ and $`q_\pi `$ denoting the integration variables, and $`X`$ is the nonpolar $`\pi N`$ amplitude (see Fig. 2) which mediates the hadronic final-state interaction; $`U^\mu `$ subsumes all exchange currents and $`m_{\mathrm{bare}}^\mu `$ is the bare contact current. For gauge invariance to hold true, the bare current must satisfy the condition $`k_\mu m_{\mathrm{bare}}^\mu =[f_s\tau ]Q_i+Q_f[f_u\tau ]+Q_\pi [f_t\tau ],`$ (11) where $`[f_x\tau ]`$ denotes the bare $`\pi NN`$ vertex the same way the notation $`[F_x\tau ]`$ was used above for the dressed vertex. This is the analog of Eq. (7) for the bare current; it is usually satisfied as a matter of course. One of the simplest nontrivial examples is the case of pure pseudovector coupling without form factors, where $`m_{\mathrm{bare}}^\mu `$ is the Kroll–Ruderman contact current ; see Eq. (27). Combining now Eq. (10) with the necessary condition (7), and making use of Eq. (11) in the Born terms, produces $`0`$ $`=`$ $`([F_s\tau ][f_s\tau ])Q_iQ_f([F_u\tau ]f_u\tau ])`$ (15) $`Q_\pi \left([F_t\tau ][f_t\tau ]\right)+k_\mu U^\mu G_0[F\tau ]`$ $`+XG_0\{k_\mu J_\mathrm{n}^\mu S_{p_\mathrm{n}k}[F_u\tau ]+k_\mu J_\pi ^\mu \mathrm{\Delta }_{q_\pi k}[F_t\tau ]`$ $`+k_\mu m_{\mathrm{bare}}^\mu +k_\mu U^\mu G_0[F\tau ]\}`$ as a necessary off-shell requirement for gauge invariance to be satisfied. In other words, as long as the basic Ward–Takahashi identities (2) for the hadron currents are true, any approximation of the full reaction mechanisms constructed in such a manner that the condition (15) is satisfied will also preserve gauge invariance as a matter of course. In view of the arbitrariness of transverse contributions, there are of course infinitely many ways this can be achieved. The prescription we give in the following applies to the simplifying assumption that one completely omits the explicit treatment of exchange currents $`U^\mu `$. However, even if they are taken into account in some partial manner, it is a straightforward exercise to adapt the following formulation to accommodate such situations. Omitting explicit exchange currents, we maintain gauge invariance by constructing auxiliary currents $`j_0^\mu `$ and $`j_1^\mu `$ which provide the same effect as the exchange currents $`U^\mu `$ as far as the preservation of gauge invariance is concerned. To this end, we make the replacements $`U^\mu G_0[F\tau ]`$ $``$ $`j_0^\mu +\mathrm{\Delta }j_0^\mu ,`$ (17) $`XG_0U^\mu G_0[F\tau ]`$ $``$ $`j_1^\mu +\mathrm{\Delta }j_1^\mu ,`$ (18) and demand that $`j_0^\mu `$ and $`j_1^\mu `$ satisfy $`k_\mu j_0^\mu `$ $`=`$ $`\left([F_s\tau ][f_s\tau ]\right)Q_i+Q_f\left([F_u\tau ][f_u\tau ]\right)`$ (20) $`+Q_\pi \left([F_t\tau ][f_t\tau ]\right)`$ and $`k_\mu j_1^\mu `$ $`=`$ $`XG_0\{k_\mu J_\mathrm{n}^\mu S_{p_\mathrm{n}k}[F_u\tau ]+k_\mu J_\pi ^\mu \mathrm{\Delta }_{q_\pi k}[F_t\tau ]`$ (22) $`+k_\mu m_{\mathrm{bare}}^\mu \}.`$ Clearly, if these conditions are met, then $$k_\mu \left(\mathrm{\Delta }j_0^\mu +\mathrm{\Delta }j_1^\mu \right)=0,$$ (23) i.e., $`\mathrm{\Delta }j_0^\mu +\mathrm{\Delta }j_1^\mu `$ is purely transverse. Without explicit treatment of exchange currents, these contributions are inaccessible and will be dropped. The resulting production current, $`M^\mu `$ $`=`$ $`M_s^\mu +M_u^\mu +M_t^\mu +m_{\mathrm{bare}}^\mu +j_0^\mu `$ (25) $`+XG_0\left\{M_u^\mu +M_t^\mu +m_{\mathrm{bare}}^\mu \right\}+j_1^\mu ,`$ will then satisfy the generalized Ward–Takahashi identity (3) and therefore it will be gauge invariant. To be more specific as to how to implement the conditions (20) and (22), allowing for a mixture of pseudoscalar and pseudovector couplings, let us write the dressed $`\pi NN`$ vertex as $$F=g_{\mathrm{ps}}\gamma _5G_{\mathrm{ps}}+g_{\mathrm{pv}}\frac{\gamma _5q/_\pi }{2m}G_{\mathrm{pv}},$$ (26) where the indices ps and pv stand for pseudoscalar and pseudovector contributions, respectively; $`G_{\mathrm{ps}}`$ and $`G_{\mathrm{pv}}`$ denote the corresponding normalized form factors (with their strength parameters $`g_{\mathrm{ps}}`$ and $`g_{\mathrm{pv}}`$ adding up to the physical coupling constant, $`g_{\pi NN}=g_{\mathrm{ps}}+g_{\mathrm{pv}}`$), $`q_\pi `$ is the four-momentum of the pion, and $`m`$ the nucleon mass. The bare vertex $`f`$ is given by the same equation with all $`G`$’s removed, and the corresponding bare current, $$m_{\mathrm{bare}}^\mu =g_{\mathrm{pv}}\frac{\gamma _5\gamma ^\mu }{2m}Q_\pi \tau ,$$ (27) is just the usual Kroll–Ruderman contact term . Equation (20) now reads explicitly $`k_\mu j_0^\mu `$ $`=`$ $`\gamma _5\left(G_sg\right)\tau Q_i`$ (32) $`+\gamma _5\left(G_ug\right)Q_f\tau `$ $`+\gamma _5\left(G_tg\right)Q_\pi \tau `$ $`g_{\mathrm{pv}}\gamma _5{\displaystyle \frac{k/}{2m}}\left(G_{\mathrm{pv},s}\tau Q_iG_{\mathrm{pv},u}Q_f\tau \right)`$ $`k_\mu m_{\mathrm{bare}}^\mu ,`$ where $$G_x=g_{\mathrm{ps}}G_{\mathrm{ps},x}+g_{\mathrm{pv}}\frac{p/p/^{}}{2m}G_{\mathrm{pv},x}$$ (33) (with $`x=s`$, $`u`$, $`t`$) denotes the kinematic situations in which the vertex functions $`G_{\mathrm{ps}}`$ and $`G_{\mathrm{pv}}`$ appear; $`g`$ is given by the same equation with all $`G`$’s removed. Note that all the terms containing $`g`$ in Eq. (32) add up to zero. We may therefore replace $`g`$ by an arbitrary function $`\widehat{F}`$ with impunity. Furthermore, using the Mandelstam variables $`s=(p+k)^2`$, $`u=(p^{}k)^2`$, and $`t=(qk)^2`$, we may then rewrite the resulting equation as $`k_\mu j_0^\mu `$ $`=`$ $`k_\mu \{{\displaystyle \frac{(2p+k)^\mu }{sp^2}}\gamma _5(G_s\widehat{F})\tau Q_i`$ (38) $`{\displaystyle \frac{(2p^{}k)^\mu }{up^2}}\gamma _5\left(G_u\widehat{F}\right)Q_f\tau `$ $`{\displaystyle \frac{(2qk)^\mu }{tq^2}}\gamma _5\left(G_t\widehat{F}\right)Q_\pi \tau `$ $`g_{\mathrm{pv}}{\displaystyle \frac{\gamma _5\gamma ^\mu }{2m}}\left(G_{\mathrm{pv},s}\tau Q_iG_{\mathrm{pv},u}Q_f\tau \right)`$ $`m_{\mathrm{bare}}^\mu \},`$ which allows us to put $`j_0^\mu `$ $`=`$ $`{\displaystyle \frac{(2p+k)^\mu }{sp^2}}\gamma _5\left(G_s\widehat{F}\right)\tau Q_i`$ (42) $`{\displaystyle \frac{(2p^{}k)^\mu }{up^2}}\gamma _5\left(G_u\widehat{F}\right)Q_f\tau `$ $`{\displaystyle \frac{(2qk)^\mu }{tq^2}}\gamma _5\left(G_t\widehat{F}\right)Q_\pi \tau `$ $`g_{\mathrm{pv}}{\displaystyle \frac{\gamma _5\gamma ^\mu }{2m}}\left(G_{\mathrm{pv},s}\tau Q_iG_{\mathrm{pv},u}Q_f\tau \right)m_{\mathrm{bare}}^\mu .`$ Comparing with Eq. (25), note that the last two terms of this gauge-invariance-preserving current cancel the bare Kroll–Ruderman term and replace it by a dressed one, where instead of the pion charge $`Q_\pi \tau `$, there is now a dressing term $`G_{\mathrm{pv},s}\tau Q_iG_{\mathrm{pv},u}Q_f\tau `$. This dressing term expresses the pion charge in terms of the nucleon charges modified by hadronic form factors, and, in general, it will be non-zero even if the pion is uncharged. We emphasize that the transition from Eq. (32) to (42) is not unique, of course, since one may add a divergence-free current to $`j_0^\mu `$ without changing the necessary condition (32). In fact, the replacement of $`g`$ by $`\widehat{F}`$ amounts to the addition of such a transverse current (which in turn may be understood as a phenomenological way of getting a handle on the neglected transverse current $`\mathrm{\Delta }j_0^\mu +\mathrm{\Delta }j_1^\mu `$). At this stage, $`\widehat{F}`$ is completely undetermined. Below, when considering the relationship of the present results to existing tree-level approaches, we will discuss some specific choices. We also note that the terms appearing in Eq. (42) do not introduce any new singularities into the amplitude. This is easily illustrated for the example of the $`s`$-channel pole diagram $`M_s^\mu `$, viz. $`M_s^\mu `$ $`=`$ $`[F_s\tau ]{\displaystyle \frac{p/+k/+m}{sm^2}}\gamma ^\mu Q_i+M_{\mathrm{t},s}^\mu `$ (43) $`=`$ $`\gamma _5G_s{\displaystyle \frac{(2p+k)^\mu }{sm^2}}\tau Q_i+g_{\mathrm{pv}}{\displaystyle \frac{\gamma _5k/}{2m}}G_{\mathrm{pv},s}{\displaystyle \frac{(2p+k)^\mu }{sm^2}}\tau Q_i`$ (45) $`+\gamma _5G_{s,q/}{\displaystyle \frac{k^\mu \gamma ^\mu k/}{sm^2}}\tau Q_i+M_{\mathrm{t},s}^\mu ,`$ where $`G_{s,q/}`$ is given by Eq. (33) with $`p/p/^{}`$ replaced by $`q/`$; $`M_{\mathrm{t},s}^\mu `$ splits off the transverse part of the electromagnetic nucleon current given in Eq. (50) below. The decomposition given here makes it immediately obvious that the effect of adding the $`s`$-channel term, with $`p^2=m^2`$, of Eq. (42) to this expression is to replace $`G_s`$ in the first term here by $`\widehat{F}`$. The same happens also for the $`u`$\- and $`t`$-channel diagrams. In other words, apart from providing a dressed Kroll–Ruderman term, the effect of the gauge-invariance preserving current (42) is to provide a common form factor $`\widehat{F}`$ for some (but not all) of the contributions originally containing individual form factors $`G_s`$, $`G_u`$, and $`G_t`$, without changing the original singularities of the Born diagrams. Next, to construct the current $`j_1^\mu `$, we recall that the gauge-invariant nucleon and pion currents appearing in the $`u`$\- and $`t`$-channel terms of the final-state interaction contribution are given by $`J_\mathrm{n}^\mu `$ $`=`$ $`\gamma ^\mu Q_\mathrm{n}+J_{\mathrm{t},\mathrm{n}}^\mu ,`$ (47) $`J_\pi ^\mu `$ $`=`$ $`(2q_\pi k)^\mu Q_\pi +J_{\mathrm{t},\pi }^\mu ,`$ (48) respectively, with transverse pieces which follow from demanding the validity of the Ward–Takahashi identities (2), i.e., $`J_{\mathrm{t},\mathrm{n}}^\mu `$ $`=`$ $`\left(\gamma ^\mu k^\mu {\displaystyle \frac{k/}{k^2}}\right)Q_\mathrm{n}(F_11)+i{\displaystyle \frac{\sigma ^{\mu \nu }k_\nu }{2m}}\kappa F_2,`$ (50) $`J_{\mathrm{t},\pi }^\mu `$ $`=`$ $`\left[(2q_\pi k)^\mu k^\mu {\displaystyle \frac{k(2q_\pi k)}{k^2}}\right]Q_\pi (F_\pi 1),`$ (51) where $`F_1`$ and $`F_2`$ respectively are the electromagnetic Dirac and Pauli form factors of the nucleon, with $`\kappa `$ being its anomalous magnetic moment, and $`F_\pi `$ is the electromagnetic form factor of the pion. Their four-divergence is given by $`k_\mu J_\mathrm{n}^\mu `$ $`=`$ $`k_\mu \gamma ^\mu Q_\mathrm{n},`$ (54) $`k_\mu J_\pi ^\mu `$ $`=`$ $`k_\mu (2q_\pi k)^\mu Q_\pi .`$ (55) Inserting this into Eq. (22), we may then extract the gauge-invariance-preserving current as $`j_1^\mu `$ $`=`$ $`XG_0\{\gamma ^\mu Q_\mathrm{n}S_{p_\mathrm{n}k}[F_u\tau ]`$ (57) $`+(2q_\pi k)^\mu Q_\pi \mathrm{\Delta }_{q_\pi k}[F_t\tau ]+m_{\mathrm{bare}}^\mu \}.`$ Its effect on the final-state interaction part of the production amplitude (25) is seen to simply cancel the bare current and to reduce the $`u`$\- and $`t`$-channel contributions to their respective transverse pieces, i.e., $`XG_0\left(M_{\mathrm{t},u}^\mu +M_{\mathrm{t},t}^\mu \right)`$ $`=`$ $`XG_0\left\{M_u^\mu +M_t^\mu +m_{\mathrm{bare}}^\mu \right\}+j_1^\mu `$ (58) $`=`$ $`XG_0\{J_{\mathrm{t},\mathrm{n}}^\mu S_{p_\mathrm{n}k}[F_u\tau ]`$ (60) $`+J_{\mathrm{t},\pi }^\mu \mathrm{\Delta }_{q_\pi k}[F_t\tau ]\}.`$ The entire contribution from the final-state interaction, therefore, is purely transverse. ### B Relation to Ohta’s and Haberzettl’s tree-level prescriptions To make the connection with tree-level approaches, we need to switch off all final-state interactions and put $`X=0`$ in Eqs. (22) and (25). The conditions to be satisfied, therefore, are Eq. (20) and Eq. (22) in the form $$k_\mu j_1^\mu =0.$$ (61) In other words, simply putting $`j_1^\mu =0`$ satisfies all gauge-invariance constraints at the tree level. Both Ohta’s and Haberzettl’s prescriptions for preserving gauge invariance can be understood as different choices for the function $`\widehat{F}`$ in Eq. (42). Ohta’s approach, based on a particular application of the minimal substitution procedure, finds $$\widehat{F}=g_{\mathrm{ps}}G_{\mathrm{ps}}(q,p^{},p)+g_{\mathrm{pv}}\frac{p/p/^{}}{2m}G_{\mathrm{pv}}(q,p^{},p),$$ (62) where the external hadron momenta of the photoproduction current—a four-point function—appear here as the momenta of the $`\pi NN`$ vertex—a three-point function. Since the momenta satisfy $`q+p^{}=p+k`$, this mismatch corresponds to an unphysical region of the $`\pi NN`$ vertex (which leads to the problems discussed in Ref. ). Only in the infrared limit of $`k=0`$, this mismatch is resolved and then this choice prevents the current (42) from being singular at $`k=0`$. In Haberzettl’s prescription, the function $`\widehat{F}`$ is a linear combination of the three kinematical situations in which the $`\pi NN`$ vertices appear in the Born terms, i.e., \[cf. Eq. (33)\] $$\widehat{F}=a_sG_s+a_uG_u+a_tG_t,$$ (63) with coefficients constrained by $`a_s+a_u+a_t=1`$, which may be fixed according to prejudice or used as free fit parameters. In contrast to Ohta’s choice, this does not require any unphysical values for the $`\pi NN`$ form factors in practical applications, and it also has a well-behaved infrared limit. In direct comparisons to Ohta’s, this prescription is found to provide better agreement with the experimental data . ## III Summary and Discussion In summary, we have treated here the electromagnetic production current for mesons off the nucleon for both real and virtual photons. We used the constraints following from requiring the validity of the generalized Ward–Takahashi identities to construct auxiliary current pieces that ensure that gauge invariance is preserved even if—as it is invariably the case in practical applications—one does not treat the problem completely and consistently. The result for the production current $`M^\mu `$ obtained here in Eqs. (25), (42), and (57) can be summarized by $`M^\mu `$ $`=`$ $`M_s^\mu +M_u^\mu +M_t^\mu +j_{\mathrm{gip}}^\mu `$ (65) $`+XG_0\left(M_{\mathrm{t},u}^\mu +M_{\mathrm{t},t}^\mu \right),`$ where the gauge-invariance-preserving current, $$j_{\mathrm{gip}}^\mu =j_0^\mu +m_{\mathrm{bare}}^\mu ,$$ (66) is given via $`j_0^\mu `$ of Eq. (42). As far as the choice of the function $`\widehat{F}`$ appearing in $`j_0^\mu `$ is concerned, the generally better results obtained with Haberzettl’s prescription at the tree-level seem to favor the form given in Eq. (63). We emphasize, however, that other functions are possible and that, in general, as far as gauge invariance is concerned, any current $`j_0^\mu `$ that satisfies the necessary condition (32) is permitted here. The part describing the hadronic final-state interaction due to $`X`$ does no longer contain the Kroll–Ruderman contact term since it is canceled by the current of Eq. (57). Moreover, what remains is seen to be entirely transverse since the corresponding $`u`$\- and $`t`$-channel contributions, $`M_{\mathrm{t},u}^\mu `$ and $`M_{\mathrm{t},t}^\mu `$, respectively contain only the transverse electromagnetic nucleon and pion operators of Eq. (II A). For real photons, in particular, both $`F_\pi 1`$ and $`F_11`$ vanish and therefore $`M_{\mathrm{t},t}^\mu =0`$, i.e., the $`t`$-channel does not contribute at all, and the $`u`$-channel current $`M_{\mathrm{t},u}^\mu `$ is reduced to the magnetic $`\sigma ^{\mu \nu }k_\nu `$ term from Eq. (50). We emphasize that even though the gauge-invariant production current of Eq. (65) is an approximation to the full dynamics of the problem as summarized in Figs. 1 and 2, it is complete as far as the longitudinal components of the current are concerned. Any additional pieces must be transverse. Let us repeat once again that the present work was restricted to pions and nucleons merely to simplify the presentation. The concepts developed here are quite general, however. An extension to other baryons and mesons, therefore, is straightforward and easily done along the lines given here. In particular, the fact that the final-state interaction contributions are purely transverse remains true if one takes into account additional intermediate hadrons ($`\mathrm{\Delta }`$, $`\rho `$, etc.) since, just as was demonstrated here for the nucleons and pions, only the transverse parts of their respective current operators survive in the final-state interaction terms. Finally, let us point out that the present results remain equally valid whether the form factors $`F`$ and the final-state amplitude $`X`$ are obtained via some sophisticated Bethe–Salpeter-type formalism or are based on a simple phenomenological model ansatz. How these elements are obtained does not enter any of the present considerations and therefore has no bearing on the question of gauge invariance. ###### Acknowledgements. The author gratefully acknowledges discussions with S. Krewald and K. Nakayama which precipitated the present work. This work was supported in part by Grant No. DE-FG02-95ER-40907 of the U.S. Department of Energy.
warning/0003/cond-mat0003160.html	ar5iv	text	# Vortex in a 𝑑-wave Superconductor at Low Temperatures ## I INTRODUCTION Vortices in classic superconductors involve a winding of $`2\pi `$ of the order parameter phase $`\varphi (𝐫)`$ around the vortex line. Since the order parameter $`\mathrm{\Delta }_𝐤=\|\mathrm{\Delta }_𝐤\|e^{i\varphi }`$ in a $`d_{x^2y^2}`$-wave superconductor like the high-$`T_c`$ cuprates is also simply a complex scalar with single global phase $`\varphi `$, it was initially expected that the vortex state in the cuprates might be structurally quite similar to the textbook case. One remarkable difference was pointed out by Joynt: in the $`d`$-wave case, if subdominant pair potential components of different symmetry exist, corresponding order parameter components can be induced at $`T_c`$ by any probe which couples to gradients of the order parameter, even if the subdominant zero field “bare” critical temperatures are very small or zero. Thus the structure of a $`d`$-wave vortex generically involves admixtures of different symmetry order parameters. Secondly, as noted by Volovik, the traditional roles played by extended and localized quasiparticle states in classic superconductors are reversed in the $`d`$-wave case. In classic superconductors at low temperatures, Caroli-de Gennes-Matricon bound states in the vortex core dominate the electronic density of states because extended states are fully gapped and therefore unoccupied. In the $`d`$-wave case, the existence of order parameter nodes inhibits the formation of the bound states (their very existence is questionable ), and populates the extended ones, which are found to dominate thermodynamics at low temperatures and magnetic fields. Several authors have attacked the $`d`$-wave mixed state structure problem in recent years, armed with these ideas. Early studies focussed on an isolated $`d`$-wave vortex, allowing for an induced $`s`$-wave order parameter component and solving the Bogoliubov-de Gennes equations on a lattice. Within the Ginzburg-Landau (GL) theory, similar results were obtained. The $`s`$ component was shown to have opposite winding number to that of the parent $`d_{x^2y^2}`$ in the core regions. Far away from the core center, it decays as $`1/r^2`$, and its winding number becomes 3,, implying that there are four extra vortices in the $`s`$-field at large distances from the main core. These results were confirmed by numerical solutions of the Eilenberger equations by Ichioka et al. The possibility of an induced $`d_{xy}`$ component was also allowed for in Ref. , which concluded that the structure in this case was similar to that of the (induced) $`s`$-wave case, except that the induced order parameter at large distances was found to decay more rapidly, roughly as $`1/r^4`$, and have opposite winding number 5. Koyama and Tachiki pointed out, however, that if the calculation is done in a gauge invariant manner there is an additional term in the free energy not found by Ichioka et al., involving a Zeeman coupling of the field to an intrinsic orbital magnetic moment in a state with structure close to a uniform $`d_{x^2y^2}+id_{xy}`$. They furthermore showed, within a GL framework with coefficients determined by BCS weak-coupling theory, that this term is proportional to the particle-hole asymmetry of the normal metal from which the superconductor condenses, and dominates sufficiently far from the vortex core. Interest in the possibility of order-parameter mixing in the vortex state was heightened by the experimental observation of a plateau in the thermal conductivity as a function of the magnetic field $`H`$ when $`H`$ is above some critical value $`H^{}`$. Krishana et al., in particular, speculated that their observation of a sharp kink at $`H^{}`$ might be explained by the sudden onset of an out-of-phase $`d_{xy}`$ component at this critical field. The new high-field state was proposed to be fully gapped, with vanishing quasiparticle transport. There are several difficulties with this explanation, which we discuss below, but theorists were nonetheless persuaded to revisit the problem. Laughlin then pointed out, in analogy to the quantum Hall state, the peculiar nature of the time-reversal symmetry breaking $`d_{x^2y^2}+id_{xy}`$ state, which appears to be quite different from possible ground states in the case of the $`d,s`$-wave mixture. He proposed that the development of a magnetic moment coupling to the magnetic field might account for the phase transition, and put forward a $`T`$\- and $`H`$-dependent free energy functional driving the phase transition, and found a critical field $`H^{}T^2`$ similar to experiment. This special free energy functional does not have a microscopic basis and it is not clear yet whether there is such a field-induced secondary phase transition. To date, no magnetic-field induced transition has been found in relevant numerical studies of the vortex lattice. Laughlin’s argument ignored the physics of the core region, thought to be negligible, but Ramakrishnan pointed out that quite similar effects are to be expected due to a combination of superfluid velocity Doppler shifts and Andreev reflection near the vortex cores. In higher fields, he proposed that these local $`d_{x^2y^2}+id_{xy}`$ patches might overlap, causing a transition to a uniform gapped state. This scenario is similar to one proposed by Movshovich et al. in the related case of magnetic impurities in a $`d`$-wave superconductor in zero field. Finally, Balatsky has recently investigated the effect of the orbital Zeeman term near the upper critical field, and argued that the $`d`$-wave state is always unstable to a $`d_{x^2y^2}+id_{xy}`$ mixture. An unusual collective clapping mode in association with the relative phase of $`d_{xy}`$ to $`d_{x^2y^2}`$ was predicted in this superconductor. Several important questions have not been addressed in the analyses of these issues thus far and have motivated this work. First and foremost, we would like to understand whether a phase transition of the type proposed by Krishana et al. is possible. Analyses in the GL regime are not applicable, and numerical calculations are not always useful to understand competing physical effects. We have therefore developed a systematic calculational approach capable of treating carefully the relevant quasiparticle states in the presence of spatially varying superflow together with the relevant subdominant order parameter components on an equal footing in the low-temperature phase. Our theory works for $`HH_{c2}`$, in which case the vortex core region can be safely neglected. Secondly, we would like to understand the structure of the vortex state and the role of the quasiparticles as a preliminary to the yet-unsolved problem of quasiparticle transport in applied magnetic field. In addition to being inapplicable at low temperatures, the GL calculations on which most of one’s intuition for this problem is based are unable to predict magnitudes of physical effects since they are based entirely on symmetry considerations. This is particularly important in the case of the orbital Zeeman coupling in the $`d_{x^2y^2},d_{xy}`$ mixing problem. The magnitude of the induced orbital moment is a very difficult quantity to estimate properly, as one might a priori deduce by analogy to the intrinsic orbital angular momentum problem in the $`{}_{}{}^{3}He`$ A-phase. In this case, it was found that naive calculations dramatically overestimated this effect, and we show here in fact that the orbital Zeeman coupling in the current problem is quite small. In this paper, we adopt a semiclassical approach, expanding the BCS free energy in powers of the local superfluid velocity, local order parameter magnitude fluctuations, and their gradients. We thus neglect states possibly localized in the vortex core, and other quasiparticle bandstructure effects in a periodic vortex lattice discussed recently by several authors . Initially Volovik and Kopnin and Volovik proposed that a semiclassical analysis of this type should be valid only down to a scale $`(\mathrm{\Delta }_0^2/E_F)\sqrt{H/H_{c2}}`$. Recent numerical work indicated, however, that the true crossover scale is much smaller for realistic systems with $`\mathrm{\Delta }_0/E_F1`$. Such fine details of the true quantum quasiparticle band structure will also be smeared out by impurity effects. We believe, therefore, that our neglect of the vortex core and quasiparticle bandstructure will be justified for cuprate superconductors at low fields ($`HH_{c2}`$) and temperatures, and that the current analysis will thus be adequate. We begin by presenting in Sec. II the method, which involves a functional integral representation of the BCS free energy $`F`$, which we then expand in powers of slow superfluid velocity gradients and small subdominant order parameter components. Analytical results for $`F`$ in the GL regime, the low temperature regime, and, for $`s`$-wave case, an ultralow temperature regime where nonlinear superflow effects dominate, are given. This allows us in Sec. III to calculate the order parameter fluctuations directly. We then apply these results to the comparison of structure of a single isolated vortex with $`s`$ or $`d_{xy}`$ subdominant pairing at various temperatures in Sec. IV, and go on in Sec. V to discuss the prospects for observing a low temperature field-induced transition of this structure. In Sec. VI we discuss existing experiments and make some comments on the various available scenarios. In Appendix A, a detailed derivation of the free energy is presented, while Appendix B is devoted to a general calculation of the spontaneous magnetization in a $`d_{x^2y^2}+id_{xy}`$-wave superconductor. ## II FREE ENERGY We start from a two-dimensional (2D) phenomenological BCS mean-field Hamiltonian in the mixed state: $`\widehat{H}_{\mathrm{MF}}`$ $`=`$ $`{\displaystyle \underset{\sigma }{}}{\displaystyle }d^2𝐫c_\sigma ^{}(𝐫)\{{\displaystyle \frac{1}{2m}}[i\stackrel{}{}{\displaystyle \frac{e}{c}}𝐀(𝐫)]^2\mu \}c_\sigma (𝐫)+{\displaystyle }d^2𝐫d^2𝐫^{}[\mathrm{\Delta }(𝐫,𝐫^{})c_{}^{}(𝐫)c_{}^{}(𝐫^{})+h.c.]`$ (2) $`{\displaystyle d^2𝐫d^2𝐫^{}V(𝐫𝐫^{})\|b(𝐫,𝐫^{})\|^2},`$ where $`V(𝐫)=V_d\mathrm{\Phi }_d(𝐫)+V_s+V_{xy}\mathrm{\Phi }_{xy}(𝐫)`$ (3) with $`\mathrm{\Phi }_i`$ characterizing the irreducible representations $`d_{x^2y^2}`$, $`s`$, and $`d_{xy}`$, for which the lowest order basis functions over a circular Fermi surface are $`\mathrm{\Phi }_{i𝐤}=\{\begin{array}{cc}\mathrm{cos}2\phi ,\hfill & i=d_{x^2y^2}\hfill \\ 1,\hfill & i=s\hfill \\ \mathrm{sin}2\phi ,\hfill & i=d_{xy}\hfill \end{array},`$ (7) $`\mathrm{\Delta }(𝐫,𝐫^{})=V(𝐫𝐫^{})b(𝐫,𝐫^{})`$ is the pairing order parameter with $`b(𝐫,𝐫^{})=c_{}(𝐫^{})c_{}(𝐫)`$, and $`𝐀(𝐫)`$ the magnetic vector potential. Throughout the paper $`\mathrm{}=k_B=1`$ units are chosen. The magnetic field in the problem is perpendicular to the 2D plane, i.e., along the $`\widehat{z}`$ direction. We assume that $`V_d<0`$, and $`V_d,V_s`$ and $`V_{xy}`$ take such values that in the absence of the magnetic field, the superconducting state is of $`d_{x^2y^2}`$-wave symmetry with $`\mathrm{\Delta }(𝐫,𝐫^{})=\mathrm{\Delta }_0\mathrm{\Phi }_d(𝐫𝐫^{})`$. In the mixed state, $`\mathrm{\Delta }(𝐫,𝐫^{})`$ takes the form of $`\mathrm{\Delta }(𝐫,𝐫^{})=e^{i\varphi (\frac{𝐫+𝐫^{}}{2})}\stackrel{~}{\mathrm{\Delta }}(\frac{𝐫+𝐫^{}}{2},𝐫𝐫^{})`$, where $`\stackrel{~}{\mathrm{\Delta }}(𝐑,\stackrel{}{\rho })=\overline{\mathrm{\Delta }}_d(𝐑)\mathrm{\Phi }_d(\stackrel{}{\rho })+𝒟_s(𝐑)+𝒟_{xy}(𝐑)\mathrm{\Phi }_{xy}(\stackrel{}{\rho })`$, with $`\overline{\mathrm{\Delta }}_d(𝐑)=\mathrm{\Delta }_0+𝒟_d(𝐑)`$. The $`𝒟_i(𝐫)`$ are the magnetic field-induced pairing order parameter deviations from their values in zero field. The partition function of Hamiltonian (2), after making a canonical transformation to eliminate the phase field of the pairing order parameter $`\varphi (𝐫)`$, is: $`𝒵=𝒵_0\mathrm{exp}\left\{{\displaystyle \underset{n}{}}\mathrm{Tr}\mathrm{ln}\left[\widehat{M}_0+\left(\begin{array}{cc}\widehat{V}_1& \widehat{𝒟}\\ \widehat{𝒟}^{}& \widehat{V}_2\end{array}\right)\right]\right\},`$ (10) where $`𝒵_0`$ $`=`$ $`\mathrm{exp}(F_0/T),`$ (11) $`F_0`$ $`=`$ $`{\displaystyle d^2𝐫\left(\frac{\|\overline{\mathrm{\Delta }}_d(𝐫)\|^2}{V_d}+\frac{\|𝒟_s(𝐫)\|^2}{V_s}+\frac{\|𝒟_{xy}(𝐫)\|^2}{V_{xy}}\right)},`$ (12) and, in the momentum representation, $`(\widehat{M}_0)_{𝐤,𝐤^{}}=\left(\begin{array}{cc}i\omega _n+ϵ_𝐤& \mathrm{\Delta }_𝐤\\ \mathrm{\Delta }_𝐤& i\omega _nϵ_𝐤\end{array}\right)\delta _{𝐤,𝐤^{}}`$ (15) $`(\widehat{V}_1)_{𝐤,𝐤^{}}`$ $`=`$ $`{\displaystyle d^2𝐫e^{i(𝐤^{}𝐤)𝐫}\left[𝐯_s(𝐫)𝐤^{}+\frac{1}{2}mv_s^2(𝐫)\right]},`$ (16) $`(\widehat{V}_2)_{𝐤,𝐤^{}}`$ $`=`$ $`{\displaystyle d^2𝐫e^{i(𝐤^{}𝐤)𝐫}\left[𝐯_s(𝐫)𝐤\frac{1}{2}mv_s^2(𝐫)\right]},`$ (17) $`\widehat{𝒟}_{𝐤,𝐤^{}}`$ $`=`$ $`{\displaystyle \underset{i}{}}(\widehat{𝒟}_i)_{𝐤,𝐤^{}}={\displaystyle \underset{i}{}}{\displaystyle d^2𝐫e^{i(𝐤^{}𝐤)𝐫}𝒟_i(𝐫)\mathrm{\Phi }_{i\frac{𝐤+𝐤^{}}{2}}}`$ (18) with $`\omega _n=(2n+1)\pi T`$ the fermion Matsubara frequency, $`\mathrm{\Delta }_𝐤=\mathrm{\Delta }_0\mathrm{\Phi }_{d𝐤}`$, and $`𝐯_s(𝐫)={\displaystyle \frac{e𝐀(𝐫)}{mc}}{\displaystyle \frac{\varphi (𝐫)}{2m}}`$ (19) the supercurrent velocity. In writing down Eq. (10) we have used the relation $`\stackrel{}{}𝐯_s(𝐫)=0`$ which corresponds to the conservation of the supercurrent. Otherwise, $`mv_s^2(𝐫)/2`$ in Eqs. (16) and (17) has no significant effect and will be neglected hereafter. The free energy resulting from Eq. (10) is $`F`$ $`=`$ $`T\mathrm{Tr}\mathrm{ln}𝒵=F_0T\mathrm{Tr}\mathrm{ln}\widehat{M}_0+T{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m}}\mathrm{Tr}\widehat{p}^m,`$ (20) where $`\widehat{p}=\widehat{g}\left(\begin{array}{cc}\widehat{V}_1& \widehat{𝒟}\\ \widehat{𝒟}^{}& \widehat{V}_2\end{array}\right),`$ (23) with $`\widehat{g}`$ the Green-function matrix which, in the momentum representation, is $`\widehat{g}_𝐤=\left(\begin{array}{cc}g_{1𝐤}& g_{2𝐤}\\ g_{2𝐤}& g_{4𝐤}\end{array}\right)=\left(\begin{array}{cc}\frac{i\omega _n+ϵ_𝐤}{W_{n𝐤}}& \frac{\mathrm{\Delta }_𝐤}{W_{n𝐤}}\\ \frac{\mathrm{\Delta }_𝐤}{W_{n𝐤}}& \frac{i\omega _nϵ_𝐤}{W_{n𝐤}}\end{array}\right).`$ (28) Here, $`W_{n𝐤}=\omega _n^2+E_𝐤^2,E_𝐤=\sqrt{ϵ_𝐤^2+\mathrm{\Delta }_k^2}.`$ (29) The calculation of the trace of $`\widehat{p}^m`$ in Eq. (20) can be done by noting that $`\mathrm{Tr}\widehat{p}^m`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \frac{d^2𝐤_j}{(2\pi )^2}d^2𝐫_je^{i(𝐤_j𝐤_{j+1})𝐫_j}\mathrm{Tr}\left(\underset{j=1}{\overset{m}{}}\widehat{p}_{𝐤_j,𝐤_{j+1}}(𝐫_j)\right)},`$ (30) where $`𝐤_{m+1}=𝐤_1`$. In most of the bulk region not close to a vortex core, $`𝐯_s`$ as well as the $`𝐯_s`$-induced $`𝒟_i`$ are spatially slowly varying functions. Thus we are allowed to expand $`𝐯_s(𝐫_j)`$ and $`𝒟_i(𝐫_j)`$ in $`\widehat{p}_{𝐤,𝐤_1}(𝐫)`$ in Eq. (30) as power series in their derivatives , $`𝐯_s(𝐫_j)`$ $``$ $`e^{(𝐫_j𝐫)\stackrel{}{}_𝐫}𝐯_s(𝐫)=𝐯_s(𝐫)+[(𝐫_j𝐫)\stackrel{}{}_𝐫]𝐯_s(𝐫)+\mathrm{},`$ (31) $`𝒟_i(𝐫_j)`$ $``$ $`e^{(𝐫_j𝐫)\stackrel{}{}_𝐫}𝒟_i(𝐫)=𝒟_i(𝐫)+[(𝐫_j𝐫)\stackrel{}{}_𝐫]𝒟_i(𝐫)+\mathrm{},`$ (32) the first few terms of which make main contribution to Tr$`\widehat{p}^m`$. This property enables us to develop a perturbation theory to obtain the free energy. The resulting calculation is straightforward but tedious, and is summarized in Appendix A. The final result for the free energy with respect to the pairing order parameters $`\stackrel{~}{F}=d^2𝐫\stackrel{~}{f}(𝐫)`$, where the free energy density $`\stackrel{~}{f}(𝐫)`$, keeping terms up to quadratic in $`𝒟_i`$, is $`\stackrel{~}{f}(𝐫)=f_{\overline{\mathrm{\Delta }}_d}(𝐫)+f_s(𝐫)+f_{xy}(𝐫)+\delta f(𝐫),`$ (33) $`f_{\overline{\mathrm{\Delta }}_d}(𝐫)={\displaystyle \frac{\|\overline{\mathrm{\Delta }}_d(𝐫)\|^2}{V_d}}T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\mathrm{ln}\left(\frac{\stackrel{~}{W}_{n𝐤}(𝐫)}{W_{n𝐤}}\right)},`$ (34) $`f_s(𝐫)=\mathrm{\Delta }_0^2(T)N_0\{L_s^{}(T,𝐯_s)\overline{𝒟}_s^{}(𝐫)+L_s^{\prime \prime }(T,𝐯_s)\overline{𝒟}_s^{\prime \prime }(𝐫)+[c_s+\eta _s^{}(T,𝐯_s)][\overline{𝒟}_s^{}(𝐫)]^2+[c_s+\eta _s^{\prime \prime }(T,𝐯_s)][\overline{𝒟}_s^{\prime \prime }(𝐫)]^2\},`$ (35) $`f_{xy}(𝐫)=F_4(𝐫)+F^{OZ}(𝐫)+\mathrm{\Delta }_0^2(T)N_0\left\{[c_{xy}+\eta _{xy}^{}(T,𝐯_s)][\overline{𝒟}_{xy}^{}(𝐫)]^2+[c_{xy}+\eta _{xy}^{\prime \prime }(T,𝐯_s)][\overline{𝒟}_{xy}^{\prime \prime }(𝐫)]^2\right\},`$ (36) $`F^{OZ}(𝐫)=\mathrm{\Delta }_0^2(T)N_0Q^{OZ}(T){\displaystyle \frac{e}{mc}}B(𝐫)\overline{𝒟}_{xy}^{\prime \prime }(𝐫),`$ (37) $`F_4(𝐫)=\mathrm{\Delta }_0^2(T)N_0\left\{L_{xy}^{}(T,𝐯_s)\overline{𝒟}_{xy}^{}(𝐫)+L_{xy}^{\prime \prime }(T,𝐯_s)\overline{𝒟}_{xy}^{\prime \prime }(𝐫)\right\},`$ (38) where and <sup>′′</sup> indicate real and imaginary parts, respectively, $`\overline{𝒟}_i=𝒟_i/\mathrm{\Delta }_0(T)`$, $`i=s,d_{xy}`$, $`c_s=(V_sN_0)^12c_d`$, $`c_{xy}=(V_{xy}N_0)^1c_d`$, with $`c_d=(V_dN_0)^1`$, $`B(𝐫)`$ is the magnetic induction, and $`\stackrel{~}{W}_{n𝐤}(𝐫)`$, $`Q^{OZ}(T)`$, $`L_i(T,𝐯_s)`$ and $`\eta _i(T,𝐯_s)`$ are defined in Appendix A. $`f_{\overline{\mathrm{\Delta }}_d}`$ in Eq. (34) is the free energy density in association with the dominant $`d_{x^2y^2}`$ component in the absence of $`𝒟_s`$ and $`𝒟_{xy}`$. $`f_s`$ and $`f_{xy}`$ are the free energy densities for the $`s`$ and $`d_{xy}`$ components, respectively. $`\delta f`$ involves terms of high orders in $`v_s`$, $`𝒟_i`$, derivatives of $`v_s`$, and mixed $`s`$ and $`d_{xy}`$ terms. Note in Eqs. (35) and (36) the quadratic terms in $`D_i`$, $`i=s,d_{xy}`$ can be easily reformulated to coincide with the familiar GL form: $`(c_i+\frac{\eta ^{}+\eta ^{\prime \prime }}{2})\|\overline{𝒟}_i(𝐫)\|^2+\frac{\eta ^{}\eta ^{\prime \prime }}{4}[\overline{𝒟}_i(𝐫)^2+\overline{𝒟}_i^{}(𝐫)^2]`$. In addition, derivatives of $`𝒟_i`$ terms are absorbed into the powers of $`𝒟_i`$ terms by partial integration, as indicated in Appendix A. With the free energy in Eqs. (33)-(38) we are in a position to investigate the vortex state. We will first show general results for an arbitrary $`𝐯_s(𝐫)`$ distribution in Sec. III, and apply the theory to the single vortex case in Sec. IV. ## III Order parameters ### A $`d`$-wave pairing order parameter $`\overline{\mathrm{\Delta }}_d(𝐫)=\mathrm{\Delta }_0+𝒟_d(𝐫)`$ In studying the dominant $`d_{x^2y^2}`$ component, $`𝒟_s`$ and $`𝒟_{xy}`$ can be set to zero since the $`\overline{\mathrm{\Delta }}_d`$-$`𝒟_s`$ and $`\overline{\mathrm{\Delta }}_d`$-$`𝒟_{xy}`$ mixing terms appear in higher orders of $`v_s`$ or its derivatives. The gap equation $`f_{\overline{\mathrm{\Delta }}_d}/\overline{\mathrm{\Delta }}_d(𝐫)=0`$ produces $`{\displaystyle \frac{\overline{\mathrm{\Delta }}_d(𝐫)}{V_d}}=T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\frac{\overline{\mathrm{\Delta }}_d(𝐫)\mathrm{cos}^22\phi }{\stackrel{~}{W}_{n𝐤}(𝐫)}},`$ (39) where $`\stackrel{~}{W}_{n𝐤}(𝐫)`$ is defined in Eq. (A6) in Appendix A. At $`H=0`$, Eq. (39) reduces to $`1=V_dT{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\frac{\mathrm{cos}^22\phi }{W_{n𝐤}}}=V_d{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\frac{\mathrm{cos}^22\phi \mathrm{tanh}(E_𝐤/2T)}{2E_𝐤}},`$ (40) which yields the well-known $`T_c`$ formula $`c_d=(1/2)\mathrm{ln}(2e^\gamma \omega _D/\pi T_c)`$, with $`\omega _D`$ the BCS cutoff and $`\gamma 0.5772`$ the Euler constant, as well as the asymptotic behaviors of the gap maximum $$\mathrm{\Delta }_0(T)\{\begin{array}{cc}3.54[T_c(TT_c)]^{1/2},\hfill & T\mathrm{\Delta }_0(T)\hfill \\ 2.14T_c0.39(T/T_c)^3T_c,\hfill & T\mathrm{\Delta }_0(T)\hfill \end{array}.$$ (41) The presence of vortices depletes the pairing order parameter $`\mathrm{\Delta }_0`$ by $`𝒟_d(𝐫)`$. Sufficiently far from the vortex cores, either $`T`$ or $`\mathrm{\Delta }_0(T)`$ is larger than $`v_s(𝐫)k_F`$ in the whole $`T`$ regime, and thus we can treat $`v_s(𝐫)k_F`$ and $`𝒟_d(𝐫)`$ perturbatively. It is straightforward to show that the normalized real part of $`𝒟_d(𝐫)`$ is $`{\displaystyle \frac{𝒟_d^{}(𝐫)}{\mathrm{\Delta }_0(T)}}\{\begin{array}{cc}\frac{1}{3/4+ϵ_Fmv_s^2(𝐫)/\mathrm{\Delta }_0^2(T)}\frac{ϵ_Fmv_s^2(𝐫)}{\mathrm{\Delta }_0^2(T)},\hfill & \mathrm{GL}\mathrm{regime}\hfill \\ \frac{2(\mathrm{ln}\mathrm{\hspace{0.17em}2})}{1[9\zeta (3)T^34(\mathrm{ln}\mathrm{\hspace{0.17em}2})ϵ_FTmv_s^2(𝐫)]/\mathrm{\Delta }_0^3(T)}\frac{ϵ_Fmv_s^2(𝐫)}{\mathrm{\Delta }_0^2(T)}\frac{T}{\mathrm{\Delta }_0(T)},\hfill & v_sk_FT\mathrm{\Delta }_0(T)\hfill \\ \frac{_{l=\pm 1}\|\mathrm{cos}(\theta +l\frac{\pi }{4})\|^3}{3[_{l=\pm 1}\|\mathrm{cos}(\theta +l\frac{\pi }{4})\|^3][v_s(𝐫)k_F]^3/\mathrm{\Delta }_0^3(T)}\frac{ϵ_Fmv_s^2(𝐫)}{\mathrm{\Delta }_0^2(T)}\frac{v_sk_F}{\mathrm{\Delta }_0(T)},\hfill & Tv_sk_F\mathrm{\Delta }_0(T)\hfill \end{array},`$ (45) where $`ϵ_F`$ is the Fermi energy and $`\zeta (3)`$ is the Riemann function. The imaginary part of $`𝒟_d(𝐫)`$ driven by the derivatives of the supercurrent is unimportant compared with $`𝒟_d^{}(𝐫)`$ in the whole temperature region, for its driving term $`F_d^{(1)}`$ in Eq. (A16) is very small as discussed in Appendix A. ### B Field-Induced $`𝒟_s(𝐫)`$ From Eqs. (33) and (35) we obtain the gap equation for the $`s`$ component, $`f_s(𝐫)/𝒟_s(𝐫)=0`$, which gives $`\overline{𝒟}_s^{}(𝐫){\displaystyle \frac{L_s^{}(T,𝐯_s)}{2[c_s+\eta _s^{}(T,𝐯_s)]}},\overline{𝒟}_s^{\prime \prime }(𝐫){\displaystyle \frac{L_s^{\prime \prime }(T,𝐯_s)}{2[c_s+\eta _s^{\prime \prime }(T,𝐯_s)]}},`$ (46) where $`L_s^{}(T,𝐯_s)`$, $`L_s^{\prime \prime }(T,𝐯_s)`$, $`\eta _s^{}(T,𝐯_s)`$ and $`\eta _s^{\prime \prime }(T,𝐯_s)`$ are defined in Appendix A. For $`T>𝐯_s𝐤_F`$, we can plug $`L_s`$ and $`\eta _s`$ obtained in Eqs. (A18)-(A21) into Eq. (46) to get the scaling functions of $`\overline{𝒟}_s^{}(𝐫)`$ and $`\overline{𝒟}_s^{\prime \prime }(𝐫)`$, $`\overline{𝒟}_s^{}(𝐫)`$ $``$ $`G_{s1}({\displaystyle \frac{\mathrm{\Delta }_0(T)}{T}},{\displaystyle \frac{ϵ_Fmv_s^2(𝐫)}{T^2}}){\displaystyle \frac{ϵ_Fm[v_{sx}^2(𝐫)v_{sy}^2(𝐫)]}{T^2}},`$ (47) $`\overline{𝒟}_s^{\prime \prime }(𝐫)`$ $``$ $`G_{s2}({\displaystyle \frac{\mathrm{\Delta }_0(T)}{T}},{\displaystyle \frac{ϵ_Fmv_s^2(𝐫)}{T^2}}){\displaystyle \frac{ϵ_F_xv_{sx}(𝐫)}{T^2}},`$ (48) where $`G_{s1}(d,z)`$ $`=`$ $`{\displaystyle \frac{h_{1s}(d)}{c_sh_{3s}(d)+2d^2h_{2s}(d)+zh_{4s}(d)}},`$ (49) $`G_{s2}(d,z)`$ $`=`$ $`{\displaystyle \frac{2h_{2s}(d)}{c_sh_{3s}(d)+zh_{4s}(d)}},`$ (50) with $`h_{1s}`$, $`\mathrm{}`$, $`h_{4s}`$ defined in Eqs. (A22) and (A23). It is easy to see that in the GL regime, $`\overline{𝒟}_s^{\mathrm{GL}}(𝐫){\displaystyle \frac{1}{c_s+\frac{3}{4}\frac{\mathrm{\Delta }_0^2(T)}{\pi ^2T^2}+2\frac{ϵ_Fmv_s^2(𝐫)}{\pi ^2T^2}}}{\displaystyle \frac{ϵ_Fm[v_{sx}^2(𝐫)v_{sy}^2(𝐫)]}{\pi ^2T^2}},`$ (51) $`\overline{𝒟}_s^{\prime \prime \mathrm{GL}}(𝐫){\displaystyle \frac{1}{c_s\frac{\mathrm{\Delta }_0^2(T)}{4\pi ^2T^2}+2\frac{ϵ_Fmv_s^2(𝐫)}{\pi ^2T^2}}}{\displaystyle \frac{ϵ_F_xv_{sx}(𝐫)}{\pi ^2T^2}}.`$ (52) It follows that the ratio between the real and imaginary parts of $`𝒟_s`$ $`{\displaystyle \frac{\overline{𝒟}_s^{\mathrm{GL}}(𝐫)}{\overline{𝒟}_s^{\prime \prime \mathrm{GL}}(𝐫)}}{\displaystyle \frac{m[v_{sx}^2(𝐫)v_{sy}^2(𝐫)]}{_xv_{sx}(𝐫)}}`$ (53) is of order unity in most of the region outside the core where $`v_s(𝐫)1/r_0`$ and $`_xv_{sx}(𝐫)1/r_0^2`$ with $`r_0`$ the distance to the closest vortex core center. In the case of $`𝐯_s𝐤_F<T\mathrm{\Delta }_0(T)`$, $`\overline{𝒟}_s^{\mathrm{lowT}}(𝐫){\displaystyle \frac{(\mathrm{ln}\mathrm{\hspace{0.17em}2})}{c_s+\frac{1}{2}+\frac{ϵ_Fmv_s^2(𝐫)}{4T\mathrm{\Delta }_0}}}{\displaystyle \frac{T}{\mathrm{\Delta }_0}}{\displaystyle \frac{ϵ_Fm[v_{sx}^2(𝐫)v_{sy}^2(𝐫)]}{\mathrm{\Delta }_0^2}},`$ (54) $`\overline{𝒟}_s^{\prime \prime \mathrm{lowT}}(𝐫){\displaystyle \frac{12(\mathrm{ln}\mathrm{\hspace{0.17em}2})\frac{T}{\mathrm{\Delta }_0}}{c_s\frac{1}{2}+2(\mathrm{ln}\mathrm{\hspace{0.17em}2})\frac{T}{\mathrm{\Delta }_0}+\frac{ϵ_Fmv_s^2(T)}{4T\mathrm{\Delta }_0}}}{\displaystyle \frac{ϵ_F_xv_{sx}(𝐫)}{\mathrm{\Delta }_0^2}}.`$ (55) So the ratio between the real and imaginary parts becomes $`{\displaystyle \frac{\overline{𝒟}_s^{\mathrm{lowT}}(𝐫)}{\overline{𝒟}_s^{\prime \prime \mathrm{lowT}}(𝐫)}}(\mathrm{ln}\mathrm{\hspace{0.17em}2}){\displaystyle \frac{T}{\mathrm{\Delta }_0}}{\displaystyle \frac{m[v_{sx}^2(𝐫)v_{sy}^2(𝐫)]}{_xv_{sx}(𝐫)}}.`$ (56) Comparing Eqs. (53) and (56) we find an extra prefactor $`T/\mathrm{\Delta }_0`$ is acquired at low $`T`$, indicative of a suppressed real part of the $`s`$ component with decreasing $`T`$. This is an interesting observation in the present work, the consequence of which on the structure of a single vortex will be discussed in Sec. IV. For $`T<𝐯_s𝐤_F`$, which can be achieved either by lowering $`T`$ in a certain spatial position or approaching the core region at a certain $`T`$, the nonlinear effects dominate over the thermal effects and the prefactor $`T/\mathrm{\Delta }_0`$ in Eq. (56) is expected to be replaced by $`v_sk_F/\mathrm{\Delta }_0`$. So the real part is negligibly small, and the imaginary part at $`T=0`$ is $`\overline{𝒟}_s^{\prime \prime T=0}(𝐫){\displaystyle \frac{1S_\theta v_s(𝐫)k_F/\mathrm{\Delta }_0}{c_s1/2+S_\theta v_s(𝐫)k_F/2\mathrm{\Delta }_0}}{\displaystyle \frac{ϵ_F_xv_{sx}(𝐫)}{\mathrm{\Delta }_0^2}},`$ (57) where $`S_\theta =_{l=\pm 1}\left\|\mathrm{cos}(\theta +l\frac{\pi }{4})\right\|`$. It is interesting to find from Eqs. (57) and (55) that the factor $`T/\mathrm{\Delta }_0`$ at $`𝐯_s𝐤_F<T`$ is replaced by $`v_sk_F/\mathrm{\Delta }_0`$ (up to some prefactor) at $`T<𝐯_s𝐤_F`$, reflecting the nonlinear effect due to the Doppler energy shift. Its effect on the free energy and the penetration depth in the Meissner state of a $`d`$-wave superconductor was extensively discussed in the previous papers by the present authors. ### C Field-Induced $`𝒟_{xy}(𝐫)`$ As shown in Eq. (36), there are two terms, $`F^{OZ}`$ and $`F_4`$, competing in driving the induced $`d_{xy}`$ component. The first term $`F^{OZ}`$, the so-called orbital Zeeman term, can be rewritten as $$F^{OZ}(𝐫)=𝐌(𝐫)𝐁(𝐫),$$ (58) with $`𝐌(𝐫)`$ the effective magnetic moment associated with the internal orbital current of Cooper pairs in a $`d_{x^2y^2}+id_{xy}`$-wave superconductor. A general derivation of this spontaneous magnetization is presented in Appendix B. As also discussed there, $`𝐌(𝐫)`$ is proportional to the particle-hole asymmetry $`\alpha `$ (see Eq. (A44)), making the Zeeman orbital effects very small for general density of states. It is therefore worth pointing out that analyses of the $`d_{x^2y^2}`$-$`d_{xy}`$ mixing problem based solely on a symmetry-based analysis of the orbital moment may lead to unrealistic results. Besides, it is important to note the very weak temperature dependence of the coefficient $`Q^{OZ}`$ of $`F^{OZ}`$, as shown in Fig. 1. The second term, $`F_4`$, is in the order $`𝒪(v_s^4)`$ and contains driving terms for $`𝒟_{xy}^{\prime \prime }`$ as well as $`𝒟_{xy}^{}`$. This term can be significant particularly if the system has small or zero particle-hole asymmetry. Note that in the $`s`$-wave case, it is irrelevant except for very short length scales of order the core size, due to the nonvanishing leading ($`𝒪(v_s^2)`$) term. The orbital Zeeman term has been invoked by Laughlin and Balatsky as the effect driving a putative transition to a time reversal symmetry-breaking $`d_{x^2y^2}+id_{xy}`$ state induced by field. It is therefore particularly interesting to investigate the relative importance of $`F^{OZ}`$ and $`F_4`$ within the current approach. For fields $`HH_{c1}`$, the overlap of vortices leads to nearly homogeneous $`B(𝐫)H`$ in space and hence $`F^{OZ}/\|𝒟_{xy}\|`$ can be taken as a constant over the bulk. On the other hand, $`F_4/\|𝒟_{xy}\|`$, which scales as $`(\xi /r)^4`$, is strongly space dependent. It increases rapidly when approaching the vortex core, but decays into the bulk. Thus we expect a critical radius $`r^{}`$ beyond which $`F^{OZ}`$ dominates over $`F_4`$, but for $`r<r^{}`$ $`F_4`$ becomes more important and determines the structure. $`r^{}`$ can be estimated from $`F_4(r^{})/F^{OZ}(r^{})1`$. It turns out to grow with increasing $`H`$ and decreasing $`T`$. In the GL and low $`T`$ regime, we have $$r^{}(T,H)=\{\begin{array}{cc}\xi \left(\frac{ϵ_F}{eH/mc}\right)^{1/4}\xi \left(\frac{R_H}{\lambda _F}\right)^{1/2},\hfill & GL\hfill \\ \xi \left(\frac{\mathrm{\Delta }_0(T)}{T}\right)^{1/4}\left(\frac{ϵ_F}{eH/mc}\right)^{1/4}\xi \left(\frac{\mathrm{\Delta }_0(T)}{T}\right)^{1/4}\left(\frac{R_H}{\lambda _F}\right)^{1/2},\hfill & \mathrm{low}T\hfill \end{array},$$ (59) with $`\lambda _F`$ the Fermi wavelength, $`\xi =v_F/\pi \mathrm{\Delta }_0`$ the superconducting coherence length, and $`R_H=\sqrt{c/eH}`$ the average intervortex distance. We summarize the above estimates in Table I, and display schematically the competition between the two driving terms in Fig. 2. Note the “low T” results for $`F_4`$ obtain only down to temperatures above the local superfluid velocity Doppler shift $`v_sk_F`$. For lower $`T`$, the perturbation calculation for $`F_4`$ that we present in Appendix A.3 breaks down because of the zero modes in $`E_𝐤`$. Instead one has to get a full expression for the local Doppler shift and the derivatives of the supercurrents which is suitable for perturbation expansion after integrating over momenta. This is apparently too complicated to be achieved analytically and thus requires a self-consistent numerical work, which is beyond the scope of the present paper. However, based on the nonlinear results for the $`s`$ component and $`Q^{OZ}`$, we do not expect anything qualitatively new in the extremely low $`T`$ case compared with the $`𝐯_s𝐤_FT\mathrm{\Delta }_0(T)`$ case. The resulting $`d_{xy}`$-wave order parameter, for $`\xi <r<r^{}`$ where $`F_4`$ dominates, is $`\overline{𝒟}_{xy}^{}(𝐫){\displaystyle \frac{L_{xy}^{}(T,𝐯_s)}{2[c_{xy}+\eta _{xy}^{}(T,𝐯_s)]}},\overline{𝒟}_{xy}^{\prime \prime }(𝐫){\displaystyle \frac{L_{xy}^{\prime \prime }(T,𝐯_s)}{2[c_{xy}+\eta _{xy}^{\prime \prime }(T,𝐯_s)]}},`$ (60) where $`L_{xy}^{}(T,𝐯_s)`$, $`L_{xy}^{\prime \prime }(T,𝐯_s)`$, $`\eta _{xy}^{}(T,𝐯_s)`$, and $`\eta _{xy}^{\prime \prime }(T,𝐯_s)`$ are defined in Eqs. (A41), (A42), (A45) and (A46), and their asmptotic behaviors can be found in Table II. In the region $`r^{}<r<R_H`$ where $`F^{OZ}`$ dominates, the $`d_{xy}`$ component reads $`\overline{𝒟}_{xy}^{}(𝐫)0,\overline{𝒟}_{xy}^{\prime \prime }(𝐫){\displaystyle \frac{Q^{OZ}(T)}{2[c_{xy}+\eta _{xy}^{\prime \prime }(T,𝐯_s)]}}{\displaystyle \frac{eH}{mc}},`$ (61) where $`Q^{OZ}(T)`$ defined in Eq. (A43). Eq. (61) and the asymptotic behaviors of $`Q^{OZ}`$, $`\eta _{xy}^{}(T,𝐯_s)`$, and $`\eta _{xy}^{\prime \prime }(T,𝐯_s)`$ shown in Tabel II lead to $`\overline{𝒟}_{xy}^{\prime \prime }(𝐫)\{\begin{array}{cc}\frac{1}{2}\alpha (2c_d1)[c_{xy}\frac{3}{8}\frac{\mathrm{\Delta }_0^2(T)}{\pi ^2T^2}+\frac{ϵ_Fmv_s^2(𝐫)}{\pi ^2T^2}]^1\frac{1}{ϵ_F}\frac{eH}{mc}\hfill & \mathrm{\Delta }_0(T)T\hfill \\ \frac{1}{2}\alpha [\mathrm{ln}(\frac{4\omega _D}{\mathrm{\Delta }_0})\frac{1}{2}6(\mathrm{ln2})\frac{T}{\mathrm{\Delta }_0}][c_{xy}\frac{1}{2}+2(\mathrm{ln2})\frac{T}{\mathrm{\Delta }_0}+\frac{ϵ_Fmv_s^2(𝐫)}{4\mathrm{\Delta }_0T}]^1\frac{1}{ϵ_F}\frac{eH}{mc}\hfill & T\mathrm{\Delta }_0(T)\hfill \end{array}.`$ (64) It is interesting to estimate the value of $`r^{}`$ in a real material. It follows from Eq. (59) that for $`TT_c`$, $`r^{GL}/R_H\sqrt{\xi /\lambda _F}(H/H_{c2})^{1/4}`$. For high-$`T_c`$ cuprates, $`\xi /\lambda _F10`$, and thus $`r^{GL}>R_H`$ for fields $`H>0.01H_{c2}`$. In materials with larger $`\xi /\lambda _F`$, $`r^{GL}`$ becomes order of $`R_H`$ for smaller fields. Since from Eq. (59) $`r^{}`$ increases with decreasing temperatures, it seems unlikely that the orbital Zeeman free energy plays an important role in determining the local order parameter in the vortex state for fields in the Tesla range. We would like to make some further remarks on the magnetic field dependence of $`\overline{𝒟}_{xy}^{\prime \prime }`$ for $`r>r^{}`$ at low temperatures. $`\overline{𝒟}_{xy}^{\prime \prime }`$ shown in Eqs. (61) and (64) are obtained from minimizing the free energy density $`f_{xy}(𝐫)`$ in Eq. (36). This free energy density is up to quadratic in $`\overline{𝒟}_{xy}^{\prime \prime }`$, which is sufficient for low field and for generic $`c_{xy}1`$. In the case of larger field and/or special case of $`c_{xy}`$ close to or smaller than 1, one has to include the free energy density term cubed in $`\overline{𝒟}_{xy}^{\prime \prime }`$. This term can be easily found to be $`(\overline{𝒟}_{xy}^{\prime \prime })^3/3`$ with the spatial dependence of the coefficient neglected. Thus the free energy density up to cubed in $`\overline{𝒟}_{xy}^{\prime \prime }`$ for $`r>r^{}`$ at low temperatures can be written as $`\stackrel{~}{f}_{xy}=\gamma B\overline{𝒟}_{xy}^{\prime \prime }+c_0(\overline{𝒟}_{xy}^{\prime \prime })^2+{\displaystyle \frac{1}{3}}(\overline{𝒟}_{xy}^{\prime \prime })^3,`$ (65) where $`c_0=c_{xy}+\eta _{xy}^{\prime \prime }(T,𝐯_s)`$ and $`\gamma =\mathrm{\Delta }_0^2N_0Q^{OZ}e/mc`$. Minimizing $`\stackrel{~}{f}_{xy}`$ with respect to $`\overline{𝒟}_{xy}^{\prime \prime }`$ we immediately get $`\overline{𝒟}_{xy}^{\prime \prime }\sqrt{c_0^2+\gamma B/2}c_0.`$ (66) It is obvious that there is crossover of the linear-$`B`$ dependence of $`\overline{𝒟}_{xy}^{\prime \prime }`$ for $`\gamma B/2c_0^2`$ to square root of $`B`$ dependence of $`\overline{𝒟}_{xy}^{\prime \prime }`$ for $`\gamma B/2c_0^2`$. This interesting behavior is shown in Fig. (3). In a periodic vortex lattice, the supercurrent field $`𝐯_s(𝐫)`$ in the London approximation reads $`𝐯_s(𝐫)={\displaystyle \frac{\pi }{m}}{\displaystyle \underset{𝐊G}{}}{\displaystyle \frac{i\stackrel{}{}_𝐫(e^{i𝐊𝐫})}{K^2+\lambda ^2}},`$ (67) and thus there will be special symmetry points where $`𝐯_s=0`$. At these points, a careful examination shows that the coefficients $`L_{xy}^{}`$ and $`L_{xy}^{\prime \prime }`$ in $`F_4`$ (Eq. (38)) vanish, and that $`𝒟_{xy}`$ is driven entirely by $`F^{OZ}`$, as shown in Eq. (66). This coincides with the numerical result of Yasui and Kita. Up to now, we have not included the effect of the $`d_{x^2y^2}`$ order-parameter suppression $`𝒟_d`$ on the $`d_{xy}`$ component, since we have shown in Sec. III A that $`𝒟_d`$ is negligibly small in the bulk region, in which we are primarily interested. However, the mechanism for a transition to a $`d_{x^2y^2}+id_{xy}`$ state proposed by Ramakrishnan involves precisely the interplay between $`𝒟_d`$ and the supercurrent near the core, leading potentially to local $`d_{x^2y^2}+id_{xy}`$ patches which then overlap at some critical field. Motivated by this suggestion, we examine the free energy terms including this effect within our approximation. These terms, which we refer to as the order parameter suppression terms $`F^{OPS}`$, are obtained from Eq. (A9-A12) (and have been already neglected in arriving at Eq. (20)). It can easily be seen that the $`𝒟_{xy}`$ component couples to derivatives of either $`𝒟_d(𝐫)`$ or of $`𝐯_s`$ in these terms, indicative of a pure nonlocal effect as found by Ramakrishnan. However, for $`𝐯_s𝐤_F<T`$, our surprising finding is that up to leading order, terms including derivatives of $`𝒟_d(𝐫)`$ vanish, leaving $`F^{OPS}={\displaystyle d^2r\left(T\underset{n}{}\frac{d^2𝐤}{(2\pi )^2}\frac{ϵ_𝐤\mathrm{cos}^22\phi }{W_{n𝐤}^2}\right)𝒟_d(𝐫)[_𝐫\times 𝐯_s(𝐫)]_z𝒟_{xy}^{\prime \prime }(𝐫)}.`$ (68) Since the $`[_𝐫\times 𝐯_s(𝐫)]_z`$ factor is dominated by the vector potential $`𝐀(𝐫)`$ rather than $`\varphi `$ part in $`𝐯_s(𝐫)`$, it may be replaced by $`eB(𝐫)/(mc)`$. Comparing Eqs. (68) and (37), we see that $`F^{OPS}`$ in the core region is of the same order as $`F^{OZ}`$ in the bulk (in particular, it is also proportional to the particle-hole asymmetry of the normal state), and may be viewed as the leading correction to $`F^{OZ}`$, if $`\mathrm{\Delta }_0(T)`$ in $`F^{OZ}`$ is replaced by $`\overline{\mathrm{\Delta }}_d(𝐫)=\mathrm{\Delta }_0(T)+𝒟_d(𝐫)`$. It is therefore clear from our previous discussion of $`F^{OZ}`$ that $`F^{OPS}`$ also gives in fact a very small effect even near the vortex core. ## IV Single vortex The results obtained in Sec. III enable us to compare the subdominant order parameters, at various temperatures, in the presence of vortices characterized by a certain superfluid velocity field $`𝐯_s(𝐫)`$. Our purpose in this section is to test these results in the concrete case of a single isolated vortex. We use cylindrical coordinates $`𝐫=(r,\theta )`$ with the origin located at the vortex core center. The phase field of the order parameter $`\varphi (𝐫)=\theta `$ leading to $`\varphi (𝐫)=(1/2mr)\widehat{\theta }`$. In the spatial regime of interest, the magnetic field is roughly homogeneous and thus $`𝐀(𝐫)=Br\widehat{\theta }`$. For $`rR_H`$, $`\|\varphi (𝐫)/2m\|eA(𝐫)/(mc)`$, so we can neglect $`𝐀`$, and simply write the supercurrent components and their derivatives as $$v_{sx}\frac{\mathrm{sin}\theta }{2mr},v_{sy}\frac{\mathrm{cos}\theta }{2mr},_xv_{sx}(𝐫)\frac{\mathrm{sin}2\theta }{2mr^2},$$ (69) except for studying the orbital Zeeman term, in which case $`A(𝐫)`$ is no longer negligible because $`\times 𝐀`$ enters. ### A $`d_{x^2y^2}`$-$`s`$ mixing In the GL regime, the $`s`$-wave subdominant order parameter has been investigated by many authors, with substantial agreement. The existing results in this limit are easily shown to be recovered in the present theory. Inserting Eq. (69) into Eqs. (51) and (52) leads, in the generic case of $`c_s1`$, to $$\overline{𝒟}_s^{GL}(𝐫)0.6c_s^1\left(\frac{\xi }{r}\right)^2(\mathrm{cos}2\theta +2i\mathrm{sin}2\theta )=0.3c_s^1\left(\frac{\xi }{r}\right)^2(3e^{2i\theta }e^{2i\theta })$$ (70) which is consistent with the results in Ref. in the bulk asymptotic region $`r\xi `$. As a result, the $`s`$ component is of four-fold symmetry, and the relative winding of $`s`$ and $`d`$ components is uniform across the whole vortex as shown in Fig. 4. In the low $`T`$ case, as discussed in Sec. III B, $`𝒟_s^{}(𝐫)`$ is smaller than $`𝒟_s^{\prime \prime }(𝐫)`$ by a factor of (ln 2)$`T/\mathrm{\Delta }_0`$, and the vortex structure at low $`T`$ is thus expected to be qualitatively different from that in the GL regime. From Eqs. (69), (54), and (55), one finds that $$\overline{𝒟}_s^{\mathrm{lowT}}(𝐫)1.23c_s^1\left(\frac{\xi }{r}\right)^2\left[(\mathrm{ln2})\frac{T}{\mathrm{\Delta }_0}\mathrm{cos}2\theta +2i\mathrm{sin}2\theta \right].$$ (71) The magnitude and phase of $`𝒟_s^{\mathrm{lowT}}`$ are shown in Fig. 4. The relative winding of $`s`$ and $`d`$ components at low $`T`$ takes place in a very narrow region of real space near antinode directions set by max$`(T/\mathrm{\Delta }_0,v_sk_F/\mathrm{\Delta }_0)`$. In Fig. 5, we show the $`T`$ dependence of $`\overline{𝒟}_s^{\prime \prime }`$ and $`\overline{𝒟}_s^{}`$ in a spatial position in the bulk. ### B $`d_{x^2y^2}`$-$`d_{xy}`$ mixing As discussed in Sec. III C, the competition between $`F^{OZ}`$ and $`F_4`$ divides the outside-vortex-core region into two rings: the outer ring region $`r^{}<r<R_H`$ is dominated by $`F^{OZ}`$ resulting in a rigid $`d_{x^2y^2}+id_{xy}`$ superconducting state with spatially nearly constant $`𝒟_{xy}^{\prime \prime }`$ obtained in Eqs. (61) and (64); In the inner ring $`\xi <r<r^{}`$, $`F_4`$ is more important and a spatially varying $`𝒟_{xy}`$ is expected. Now we show the single vortex structure in the inner ring region. We first insert Eq. (69) into $`U(𝐯_s)`$ defined in Eqs. (A36), (A37), (A29), (A32), and (A40) to find that $`U_1^{}(𝐯_s)=8U_2^{}(𝐯_s)=32U_3^{}(𝐯_s)={\displaystyle \frac{2ϵ_F^2\mathrm{sin}4\theta }{m^2r^4}},`$ (72) $`U_1^{\prime \prime }(𝐯_s)={\displaystyle \frac{1}{96}}U_2^{\prime \prime }(𝐯_s)={\displaystyle \frac{ϵ_F^2\mathrm{cos}4\theta }{4m^2r^4}}.`$ (73) Equations (60), (72) and (73) together with the asymptotic behaviors of $`L_{xy}`$ in Tabel II imply that in the single vortex case, $`L_{xy}^{}Y^{}(T)\left({\displaystyle \frac{\xi }{r}}\right)^4\mathrm{sin}4\theta ,L_{xy}^{\prime \prime }Y^{\prime \prime }(T)\left({\displaystyle \frac{\xi }{r}}\right)^4\mathrm{cos}4\theta ,`$ (74) where $`Y^{}7.4`$ in the GL regime, and $`7.1\mathrm{\Delta }_0(T)/T`$ for $`T\mathrm{\Delta }_0`$, while $`Y^{\prime \prime }9.8`$ in the GL regime, and $`25.9\mathrm{\Delta }_0(T)/T`$ for $`T\mathrm{\Delta }_0`$. Equations (60), (74) and the asymptotic behaviors of $`\eta _{xy}`$ in Tabel II yield $`\overline{𝒟}_{xy}^{\mathrm{GL}}(𝐫)3.7\left[c_{xy}{\displaystyle \frac{\mathrm{\Delta }_0^2(T)}{8(\pi T)^2}}+{\displaystyle \frac{ϵ_Fmv_s^2(𝐫)}{(\pi T)^2}}\right]^1\left({\displaystyle \frac{\xi }{r}}\right)^4\mathrm{sin}4\theta ,`$ (75) $`\overline{𝒟}_{xy}^{\prime \prime \mathrm{GL}}(𝐫)4.9\left[c_{xy}{\displaystyle \frac{3}{8}}{\displaystyle \frac{\mathrm{\Delta }_0^2(T)}{(\pi T)^2}}+{\displaystyle \frac{ϵ_Fmv_s^2(𝐫)}{(\pi T)^2}}\right]^1\left({\displaystyle \frac{\xi }{r}}\right)^4\mathrm{cos}4\theta ,`$ (76) and hence for $`c_{xy}1`$ $`\overline{𝒟}_{xy}^{\mathrm{GL}}(𝐫)0.61ic_{xy}^1\left({\displaystyle \frac{\xi }{r}}\right)^4(e^{i4\theta }+7e^{i4\theta }).`$ (77) This result coincides with that of Ref. . At low $`T`$, we obtain $`\overline{𝒟}_{xy}^{\mathrm{low}T}(𝐫)`$ $``$ $`3.5\left[c_{xy}+{\displaystyle \frac{ϵ_Fmv_s^2(𝐫)}{4\mathrm{\Delta }_0T}}\right]^1{\displaystyle \frac{\mathrm{\Delta }_0}{T}}\left({\displaystyle \frac{\xi }{r}}\right)^4\mathrm{sin}4\theta ,`$ (78) $`\overline{𝒟}_{xy}^{\prime \prime \mathrm{low}T}(𝐫)`$ $``$ $`12.9\left[c_{xy}{\displaystyle \frac{1}{2}}+2(\mathrm{ln2}){\displaystyle \frac{T}{\mathrm{\Delta }_0}}+{\displaystyle \frac{ϵ_Fmv_s^2(𝐫)}{4\mathrm{\Delta }_0T}}\right]^1{\displaystyle \frac{\mathrm{\Delta }_0}{T}}\left({\displaystyle \frac{\xi }{r}}\right)^4\mathrm{cos}4\theta ,`$ (79) leading, for $`c_{xy}1`$, to $`\overline{𝒟}_{xy}^{\mathrm{lowT}}(𝐫)4.69ic_{xy}^1\left({\displaystyle \frac{\mathrm{\Delta }_0}{T}}\right)\left({\displaystyle \frac{\xi }{r}}\right)^4(e^{i4\theta }+1.76e^{i4\theta }).`$ (80) In Fig. 6, we show the winding of $`d_{xy}`$ component and the magnitude of the normalized $`\overline{𝒟}_{xy}`$ at high $`T`$ (GL) and low $`T`$ respectively. The eight-fold symmetry is more obvious at low $`T`$. Since many current theories postulate a homogeneous $`d_{x^2y^2}+id_{xy}`$ state without justification, it is interesting to consider the size of the spatially varying part of $`𝒟_{xy}`$ relative to its homogeneous component. The smallest relevant value of the ratio $`𝒟_{xy}^{F_4}/𝒟_{xy}^{OZ}`$, where the superscripts $`F_4`$ and $`OZ`$ indicate the relevant driving terms in Eqs. (60) and (64), respectively, is attained at $`r=R_H`$ in the physically relevant regime $`R_H<r^{}`$. The value is $`(E_F/\mathrm{\Delta }_0)^2(H/H_{c2})`$ in the GL regime, and $`(E_F/\mathrm{\Delta }_0)^2(H/H_{c2})(\mathrm{\Delta }_0/T)`$ at low temperatures; a simple estimate then shows that the spatially fluctuating component is always at least an order of magnitude larger than the homogeneous component at experimentally relevant temperatures and fields. ## V Is there magnetic field-induced pseudo-phase transition? The results presented in Sections III and IV suggest that there always exist field-induced $`s`$ and $`d_{xy}`$ components in a parent $`d_{x^2y^2}`$-wave superconducting state where the electronic interactions in the $`s`$ and $`d_{xy}`$ channels are nonzero. These subdominant components are spatially inhomogeneous and, in a general case of small $`V_s`$ and $`V_{xy}`$ compared with $`V_d`$, are negligibly small as far as any bulk physical quantity is concerned. However, there is also a special situation in which either $`V_s`$ or $`V_{xy}`$ is nearly degenerate with $`V_d`$, leading to a small $`c_s`$ or $`c_{xy}`$. In this case, the denominator(s) of $`𝒟_i,i=s,d_{xy}`$ may vanish at some critical temperature $`T_{ci}^{}`$, a singularity marking a second pseudo-phase transition into a $`d_{x^2y^2}+𝒟_i`$ state with $`𝒟_i`$ a homogeneous bulk quantity. Investigation of such a possible pseudo-phase transition is particularly interesting in association with the experimentally observed thermal-conductivity plateau as mentioned in the Introduction. Since in high-$`T_c`$ cuprates no such phase transition has been reported to be found in the absence of magnetic field, we focus on the question whether there can be magnetic field-driven phase transition. This is equivalent to searching for a nonzero $`T_{ci}^{}`$ at a finite field which vanishes at zero field. We first study the $`s`$ component. From Eqs. (51), (52), (54), and (55), we see that the denominator of $`𝒟_s^{}`$ can never be zero, and the singularity may occur in $`𝒟_s^{\prime \prime }`$. In the GL regime, the critical temperatures are $`T_{cs}^{}(𝐯_s)={\displaystyle \frac{1}{\pi }}\sqrt{c_s^1[1/42ϵ_Fmv_s^2/\mathrm{\Delta }_0^2(T_{cs}^{})]}\mathrm{\Delta }_0(T_{cs}^{}).`$ (81) It follows that $`T_{cs}^{}(𝐯_s)<T_{cs}^{}(0)`$, implying there is no instability at nonzero field leading to a homogeneous subdominant $`𝒟_i`$. At low $`T`$, $`T_{cs}^{}`$ becomes $`T_{cs}^{}(𝐯_s)={\displaystyle \frac{1}{4(\mathrm{ln}\mathrm{\hspace{0.17em}2})}}\left[12c_s{\displaystyle \frac{ϵ_Fmv_s^2}{2T_{cs}^{}\mathrm{\Delta }_0}}\right]\mathrm{\Delta }_0.`$ (82) Again, $`T_{cs}^{}(𝐯_s)<T_{cs}^{}(0)`$ is found, meaning that the magnetic field does not favor such a second-order phase transition from a $`d_{x^2y^2}`$ to $`d_{x^2y^2}+is`$ state. Accounting for nonlinear effects does not affect this conclusion. A similar analysis can be made in the $`d_{xy}`$ situation. From the asymptotic behaviors of $`\eta _{xy}^{}`$ and $`\eta _{xy}^{\prime \prime }`$ listed in Table II, we find that there is also a singularity in $`𝒟_{xy}^{}`$ in the GL regime, but its corresponding critical temperatures are smaller than that in $`𝒟_{xy}^{\prime \prime }`$. Except for this point, we reach the same conclusion as the $`s`$ case: No magnetic field-induced phase transition is found. At this stage, we may examine the prospects of a possible low-$`T`$ field-induced first-order phase transition into a time-reversal symmetry breaking state $`d_{x^2y^2}+id_{xy}`$ as suggested by Laughlin . The free energy functional of the $`d_{xy}`$ pairing order parameter assumed by Laughlin takes the form $`{\displaystyle \frac{F}{L^2}}={\displaystyle \frac{1}{6\pi }}{\displaystyle \frac{(𝒟_{xy}^{\prime \prime })^3}{\nu ^2}}{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{eB}{c}}𝒟_{xy}^{\prime \prime }\mathrm{tanh}^2\left({\displaystyle \frac{𝒟_{xy}^{\prime \prime }}{2T}}\right){\displaystyle \frac{4}{\pi }}{\displaystyle \frac{T^3}{\nu ^2}}\left\{{\displaystyle \frac{(𝒟_{xy}^{\prime \prime })^2}{2T^2}}\mathrm{ln}[1+e^{𝒟_{xy}^{\prime \prime }/T}]+{\displaystyle _{𝒟_{xy}^{\prime \prime }/T}^{\mathrm{}}}\mathrm{ln}[1+e^x]x𝑑x\right\},`$ (83) where $`\nu =\sqrt{2\mathrm{\Delta }_0/m}`$. The special $`T`$ dependences of the second and third terms lead to a weak first-order phase transition at $`T0.52\nu \sqrt{2eB/c}`$. We may check the justification of this free energy based on the present theory. The first and third terms in the right-hand side of Eq. (83) are understood to come from the local free energy $`F^{(loc)}`$ in Eq. (A5). The second one, as the coupling of the unusual magnetization to the field, is naturally related to the orbital Zeeman term $`F^{OZ}`$ in Eq. (58) with $`M(𝐫)`$ shown in Eq. (B14). But the assumed temperature dependence crucial for the first-order phase transition, is inconsistent with that of $`F^{OZ}`$, which we have shown to be a weak function of $`T`$ for $`T\mathrm{\Delta }_0`$ (see Fig. 1) and a linear function of $`𝒟_{xy}^{\prime \prime }`$ in the limit $`𝒟_{xy}^{\prime \prime }0`$. Therefore, the microscopic calculation in the present work does not confirm Laughlin’s free energy functional which formed the basis of the first-order phase transition found in his work. In the present theory, a small homogenous $`𝒟_{xy}^{\prime \prime }`$ is found for any magnetic field and temperature in the superconducting state. It is swamped by a larger (but still much smaller than $`\mathrm{\Delta }_0`$) $`𝒟_{xy}^{\prime \prime }`$ spatially fluctuating component over almost all of the vortex lattice for physically relevant fields. ## VI Conclusions In this paper, we have formulated a perturbation theory to investigate the magnetic-field induced subdominant order parameters of a clean $`d`$-wave superconductor in the presence of the gauge-invariant spatially varying supercurrent field in the mixed state. With the assumption of slowly spatially varying supercurrents and their induced $`s`$ and $`d_{xy}`$ components in the bulk sample, we are able to derive the free energy as power series in the Doppler energy shift, the derivatives of the supercurrents, and the subdominant components. The free energy is valid from $`T_c`$ down to very low temperatures, enabling us to compare the resulting $`s`$ and $`d_{xy}`$ components at low $`T`$ with the existing results in the GL regime. To leading order, the real and imaginary parts of the $`s`$ component, driven by local $`m(v_{sx}^2v_{sy}^2)`$ and the derivative of the supercurrent $`_xv_{sx}`$, respectively, were shown to have very different temperature dependences. In the GL regime, both the real and imaginary parts are in the same order. But at low $`T`$, the real part acquires an extra small prefactor $`T/\mathrm{\Delta }_0`$ and the resulting $`s`$ winding happens within a very small region near antinodal directions, leaving a rigid $`d_{x^2y^2}+is`$ state over most of the vortex. It is important to note, however, that this structure does not imply a gap in the quasiparticle spectrum, due to the small size of the $`s`$ component compared with the large Doppler shifts near the core. To leading order, the $`d_{xy}`$ component is driven by two terms of different symmetries competing over different parts of the vortex lattice. The first is the orbital Zeeman term $`F^{OZ}`$ arising from the coupling of the spontaneous magnetization to the magnetic field. Its significance is limited by its small magnitude due to the particle-hole asymmetric effects. The second driving term $`F_4`$ scales as $`(\xi /r)^4`$ for $`\xi rR_H`$. The crossover scale $`r^{}`$ divides the bulk region into inner and outer regions dominated by the two distinct physical effects. In the inner region, $`F_4`$ determines the vortex structure, and the $`d_{xy}`$ component has eight-fold symmetry and its relative phase winds 4 times that of the $`d_{x^2y^2}`$ component. In the outer region, $`F^{OZ}`$ is more important, leading to a rigid $`d_{x^2y^2}+id_{xy}`$ superconducting state, where the $`d_{xy}`$ component is spatially nearly homogeneous and a weak function of temperature. We have shown that the crossover scale is of order $`r^{}=\xi (R_H/\lambda _F)^{1/2}`$ in the GL regime, and becomes $`(\mathrm{\Delta }_0/T)^{1/4}`$ times larger at low $`T`$. Our best estimate for the high-$`T_c`$ cuprates suggests that $`r^{}`$ is greater than the intervortex separation $`R_H`$ for fields above $`0.01H_{c2}`$, so that it appears that vortex lattice structure in fields of order Tesla is governed by $`F_4`$ and the orbital Zeeman effect is irrelevant. The relative phase between the $`d_{x^2y^2}`$ and $`d_{xy}`$ components, as well as the magnitude of the induced $`d_{xy}`$, are therefore strongly space dependent everywhere in the sample. Our results have implications for several scenarios which have been proposed to create a state without quasiparticle excitations at low temperatures via the creation of a finite out-of-phase subdominant pair component. No such bulk state is found for generic values of the pair potentials $`V_s`$ and $`V_d`$. It remains possible that induced core $`d_{xy}`$ patches overlap with increasing field, as proposed by Ramakrishnan, leading to a gapped state at high fields beyond the scope of our analysis. However, analyzing the interplay between order parameter suppression $`𝒟_d(𝐫)`$ around the vortex cores and $`𝐯_s(𝐫)`$, we found it to be a very small effect on the induced $`d_{xy}`$ component; it therefore seems unlikely that the field scale where such a transition may occur can be significantly less than $`H_{c2}`$. Finally, we searched for a possible second order magnetic field-induced phase transition for special values of the coupling constants, and reached a negative conclusion. No such phase transition into a bulk $`d_{x^2y^2}+s`$ or a $`d_{x^2y^2}+d_{xy}`$ state is found unless the transition has already taken place in the absence of the field, which is apparently not the case in the high-$`T_c`$ cuprates. Examining the Laughlin free energy driving a first-order phase transition, we found that the crucial field-dependent term has an assumed temperature dependence inconsistent with the BCS theory; thus this phase transition picture is not supported in the present work, consistent with numerical results of Yasui and Kita. The apparent phase transition observed by Krishana et al. has not been clearly reproduced by other groups, and the possibility exists that the effect is due to inhomogeneously trapped flux. It is still interesting, however, to ask what kinds of intrinsic phase transitions might be possible in a $`d`$-wave superconductor. We have shown that it is unlikely that any bulk phase transition can be induced by a magnetic field, at least in the low-field regime where our approach is valid. Since our model neglects the vortex core regions, it is conceivable that vortex core transitions such as those observed in the $`{}_{}{}^{3}He`$ system might still be relevant. However, since the number of bound states in the core region is small and possibly zero for the present case, it seems unlikely that such a transition would have an important effect on the quasiparticles responsible for heat transport at low $`T`$. A final possibility, currently under investigation within the present framework, is that transitions occur in the vortex lattice structure as a function of field. ## ACKNOWLEDGMENTS The authors are grateful to N. Andrei, M. Fogelström, M. Franz, C. D. Gong, T. Kita, T. Kopp, S. Sachdev, S. H. Simon, and Y.-J. Wang for helpful communications. Partial support was provided by the A. v. Humboldt Foundation, NSF Grants No. DMR-9974396 and INT-9815833, and “Graduiertenkolleg Anwendungen der Supraleitung” of the Deutsche Forschungsgemeinschaft. ## A Derivation of Free Energy In this Appendix we derive the free energy $`\stackrel{~}{F}`$ in Eqs. (33-38) which is valid for space region where $`𝐯_s`$ varies slowly. We begin with rewriting Tr$`\widehat{p}^m`$ in Eq. (20) according to Eqs. (30) and (32), $`\mathrm{Tr}\widehat{p}^m(\mathrm{Tr}\widehat{p}^m)^{(\mathrm{loc})}+(\mathrm{Tr}\widehat{p}^m)^{(\mathrm{der})},`$ (A1) where $`(\mathrm{Tr}\widehat{p}^m)^{(\mathrm{loc})}`$ corresponds to $`𝐯_s(𝐫_j)`$ and $`𝒟(𝐫_j)`$ for all $`j=1,\mathrm{},m1`$ taking the local values $`𝐯_s(𝐫)`$ and $`𝒟(𝐫)`$, respectively, and becomes $`(\mathrm{Tr}\widehat{p}^m)^{(\mathrm{loc})}={\displaystyle \frac{d^2𝐤}{(2\pi )^2}d^2𝐫\mathrm{Tr}\left[\widehat{p}_{𝐤,𝐤}(𝐫)\right]^m},`$ (A2) and $`(\mathrm{Tr}\widehat{p}^m)^{(\mathrm{der})}`$ contains derivatives of $`𝐯_s(𝐫_j)`$ and/or $`𝒟(𝐫_j)`$. Inserting Eq. (A1) into Eq. (20) leads to $`FF^{(\mathrm{loc})}+F^{(\mathrm{der})},`$ (A3) where $`F^{(\mathrm{loc})}`$ is just the local free energy obtained in the semiclassical approximation, $`F^{(\mathrm{loc})}`$ $`=`$ $`F_0T\mathrm{Tr}\mathrm{ln}\widehat{M}_0+T{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m}}(\mathrm{Tr}\widehat{p}^m)^{(\mathrm{loc})}=F_0{\displaystyle d^2𝐫T\underset{n}{}\frac{d^2𝐤}{(2\pi )^2}\mathrm{ln}[\stackrel{~}{W}_{n𝐤}(𝐫)\eta _𝐤(𝐫)]},`$ (A4) $`=`$ $`F_0{\displaystyle d^2𝐫\frac{d^2𝐤}{(2\pi )^2}\sqrt{\stackrel{~}{E}_𝐤^2(𝐫)+\eta _𝐤(𝐫)}}{\displaystyle d^2𝐫\frac{d^2𝐤}{(2\pi )^2}\underset{l=\pm 1}{}\mathrm{ln}\left\{1+e^{[l𝐯_s(𝐫)𝐤_F+\sqrt{\stackrel{~}{E}_𝐤^2(𝐫)+\eta _𝐤(𝐫)}]/T}\right\}},`$ (A5) with $`\stackrel{~}{W}_{n𝐤}(𝐫)`$ $`=`$ $`[i\omega _n+𝐯_s(𝐫)𝐤_F]^2+\stackrel{~}{E}_𝐤^2(𝐫),`$ (A6) $`\stackrel{~}{E}_𝐤(𝐫)`$ $`=`$ $`\sqrt{ϵ_𝐤^2+\|\overline{\mathrm{\Delta }}(𝐫)\|^2\mathrm{\Phi }_{d𝐤}^2}`$ (A7) $`\eta _𝐤(𝐫)`$ $``$ $`{\displaystyle \underset{i=s,d_{xy}}{}}\{[\overline{\mathrm{\Delta }}(𝐫)𝒟_i^{}(𝐫)+h.c.]\mathrm{\Phi }_{d𝐤}\mathrm{\Phi }_{i𝐤}+\|𝒟_i(𝐫)\|^2\mathrm{\Phi }_{i𝐤}^2\}.`$ (A8) Expanding $`F^{(\mathrm{loc})}`$ in power series in $`𝒟_i`$ will give linear-$`\overline{𝒟}_i^{}`$ term as well as quadratic of $`𝒟_i`$ terms. The prefactors are not universal in the whole temperature regimes, as will be discussed and shown below. $`F^{(\mathrm{der})}`$ in Eq. (A3) reads $`F^{(\mathrm{der})}=T{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m}}(\mathrm{Tr}\widehat{p}^m)^{(\mathrm{der})}={\displaystyle \underset{j}{}}F^{(j)},`$ (A9) where we expanded $`(\mathrm{Tr}\widehat{p}^m)^{(\mathrm{der})}`$ as power series in the $`j`$th derivatives of $`𝐯_s`$ or $`𝒟_i`$ with respect to $`𝐫`$: $`(\mathrm{Tr}\widehat{p}^m)^{(der)}=_j(\mathrm{Tr}\widehat{p}^m)^{(j)}`$. $`F^{(\mathrm{der})}`$ reflects nonlocal couplings of the subdominant order parameters to the supercurrent fields. We examine the formal leading term, $`F^{(1)}`$ which includs $`_𝐫v_s(𝐫)`$ or $`_𝐫𝒟_i(𝐫)`$, in Eq. (A9), $`F^{(1)}={\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m}}(\mathrm{Tr}\widehat{p}^m)^{(1)}={\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m}}(\rho _{m1}+\rho _{m2}+\rho _{m3}),`$ (A10) $`\rho _{m1}=i{\displaystyle \underset{\nu =0}{\overset{m2}{}}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}d^2𝐫\mathrm{Tr}\left\{[_𝐫\widehat{p}_{𝐤,𝐤}(𝐫)]\widehat{p}_{𝐤,𝐤}^{m2\nu }(𝐫)[_𝐪\widehat{p}_{𝐪,𝐤}(𝐫)]_{𝐪=𝐤}\widehat{p}_{𝐤,𝐤}^\nu (𝐫)\right\}},`$ (A11) $`\rho _{m2}=i{\displaystyle \underset{\nu =0}{\overset{m2}{}}}\nu {\displaystyle \frac{d^2𝐤}{(2\pi )^2}d^2𝐫\mathrm{Tr}\left\{[_𝐫\widehat{p}_{𝐤,𝐤}(𝐫)]\widehat{p}_{𝐤,𝐤}^{m2\nu }(𝐫)[_𝐤\widehat{p}_{𝐪,𝐤}(𝐫)]\widehat{p}_{𝐤,𝐤}^\nu (𝐫)\right\}},`$ (A12) $`\rho _{m3}=i(m1){\displaystyle \frac{d^2𝐤}{(2\pi )^2}d^2𝐫\mathrm{Tr}\left\{[(_𝐫_𝐪)\widehat{p}_{𝐤,𝐪}(𝐫)]_{𝐪=𝐤}\widehat{p}_{𝐤,𝐤}^{m1}(𝐫)\right\}}.`$ (A13) $`F^{(1)}`$ can be expanded as power series in $`𝒟_i`$. Note in terms including $`_𝐫𝒟_i(𝐫)`$ the derivative can be transfered to that of $`𝐯_s`$ by partial integral. The resulting linear-$`𝒟_i`$ term in $`F^{(1)}`$ is $`F_𝒟^{(1)}=_{i=d,s,d_{xy}}F_i^{(1)}`$, where $`F_i^{(1)}`$ $``$ $`i{\displaystyle d^2𝐫T\underset{n}{}\frac{d^2𝐤}{(2\pi )^2}\mathrm{\Phi }_{i𝐤}\{_𝐫[𝐯_s(𝐫)𝐤_F]\}\underset{m=2}{\overset{\mathrm{}}{}}[𝐯_s(𝐫)𝐤_F]^{m2}\underset{\nu =1}{\overset{m1}{}}\nu \mathrm{Tr}\left\{(_𝐤\widehat{g}_𝐤)\widehat{g}_𝐤^\nu \widehat{𝒟}_i(𝐫)\widehat{g}_𝐤^{m\nu 1}\right\}}`$ (A14) $`=`$ $`{\displaystyle d^2𝐫𝒟_i^{\prime \prime }(𝐫)T\underset{n}{}\frac{d^2𝐤}{(2\pi )^2}\mathrm{\Phi }_{i𝐤}2\mathrm{\Delta }_𝐤[_𝐤ϵ_𝐤_𝐫][𝐯_s(𝐫)𝐤_F]Z_{n𝐤}(𝐯_s)}`$ (A16) $`{\displaystyle d^2𝐫𝒟_i^{\prime \prime }(𝐫)T\underset{n}{}\frac{d^2𝐤}{(2\pi )^2}\mathrm{\Phi }_{i𝐤}2ϵ_𝐤[_𝐤\mathrm{\Delta }_𝐤_𝐫][𝐯_s(𝐫)𝐤_F]Z_{n𝐤}(𝐯_s)},`$ with $`\widehat{𝒟}_i(𝐫)=\left(\begin{array}{cc}0\hfill & 𝒟_i(𝐫)\hfill \\ 𝒟_i^{}(𝐫)\hfill & 0\hfill \end{array}\right)`$ and $`Z_{n𝐤}(𝐯_s)={\displaystyle \frac{\left\{[𝐯_s(𝐫)k_F]^2W_{n𝐤}\right\}^24\omega _n^2[𝐯_s(𝐫)k_F]^2}{\left\{[(𝐯_s(𝐫)k_F)^2W_{n𝐤}]^2+4\omega _n^2[𝐯_s(𝐫)k_F]^2\right\}^2}}.`$ (A17) It is easy to see that $`F_d^{(1)}`$ is negligibly small, so $`F_{\overline{\mathrm{\Delta }}}`$ in Eq. (34) is simply obtained from $`F^{(loc)}`$ in Eq. (A4) by setting $`𝒟_s,𝒟_{xy}=0`$. We argue that $`F^{(loc)}`$ and $`F_s^{(1)}`$ in Eqs. (A4) and (A16) are sufficient for studying the $`s`$ component up to leading orders. Expanding $`F^{(loc)}`$ as power series in $`𝒟_s`$ and $`v_sk_F`$, we will see that the driving term for $`𝒟_s^{}`$ is in leading order of $`m(v_{sx}^2v_{sy}^2)`$, which, in the spatial regime of interest $`\xi <r<R_H`$, scales as $`1/r^2`$. $`F_s^{(1)}`$ gives a driving term for $`𝒟_s^{\prime \prime }`$ scaling as $`_xv_{sx}1/r^2`$. Clearly, terms contained in $`_{l>1}F^{(l)}`$ are of higher orders of $`v_s`$ and derivatives of $`v_s`$. We first show results for $`s`$ component. The $`d_{xy}`$ situation is more complicated and will be discussed later. ### 1 Scaling expressions for $`L_s(T,𝐯_s)`$ and $`\eta _s(T,𝐯_s)`$ at $`T>𝐯_s𝐤_F`$ In this temperature regime, we can expand $`\stackrel{~}{W}_{n𝐤}(𝐫)`$ in Eqs. (A4) and (A16) as power series in $`v_sk_F`$, $`𝒟_i^{}(𝐫)`$, and $`𝒟_i^{\prime \prime }(𝐫)`$ to find explicit expressions for $`L_s(T,𝐯_s)`$ and $`\eta _s(T,𝐯_s)`$ in Eq. (35), $`L_s^{}(T,𝐯_s)=N_0^1T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\frac{2\mathrm{cos}^22\phi (4\omega _n^2W_{n𝐤})}{W_{n𝐤}^3}ϵ_Fm[v_{sx}^2v_{sy}^2]}=2{\displaystyle \frac{ϵ_Fm[v_{sx}^2v_{sy}^2]}{T^2}}h_{1s}\left(d\right),`$ (A18) $`L_s^{\prime \prime }(T,𝐯_s)=T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\frac{4\mathrm{cos}^22\phi }{W_{n𝐤}^2}ϵ_F_xv_{sx}}=4{\displaystyle \frac{ϵ_F_xv_{sx}}{T^2}}h_{2s}\left(d\right),`$ (A19) $`\eta _s^{}(T,𝐯_s)=\eta _s^{\prime \prime }(T,𝐫)+2d^2h_{2s}(d),`$ (A20) $`\eta _s^{\prime \prime }(T,𝐯_s)=h_{3s}(d)+{\displaystyle \frac{ϵ_Fmv_s^2(𝐫)}{T^2}}h_{4s}(d),`$ (A21) where Eq. (40) was used, $`d=\mathrm{\Delta }_0(T)/T`$, and $`h_{1i}(d)={\displaystyle \frac{8}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑x{\displaystyle _0^{\pi /4}}𝑑\phi \mathrm{cos}^22\phi \mathrm{\Phi }_{i𝐤}^2{\displaystyle \frac{1}{2\eta }}{\displaystyle \frac{d^2f(\eta )}{d\eta ^2}},h_{2i}(d)={\displaystyle \frac{8}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑x{\displaystyle _0^{\pi /4}}𝑑\phi \mathrm{cos}^22\phi \mathrm{\Phi }_{i𝐤}^2{\displaystyle \frac{t(\eta )}{4\eta ^3}},`$ (A22) $`h_{3i}(d)={\displaystyle \frac{8}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑x{\displaystyle _0^{\pi /4}}𝑑\phi (12\mathrm{cos}^22\phi )\mathrm{\Phi }_{i𝐤}^2{\displaystyle \frac{\mathrm{tanh}(\eta /2)}{2\eta }},h_{4i}(d)={\displaystyle \frac{8}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑x{\displaystyle _0^{\pi /4}}𝑑\phi \mathrm{\Phi }_{i𝐤}^2{\displaystyle \frac{1}{2\eta }}{\displaystyle \frac{d^2f(\eta )}{d\eta ^2}},`$ (A23) with $`\eta =\sqrt{x^2+d^2\mathrm{cos}^22\phi }`$, $`f(\eta )=1/(e^\eta +1)`$, and $`t(\eta )=12f(\eta )+2\eta df(\eta )/d\eta `$. ### 2 $`L_s(T,𝐯_s)`$ and $`\eta _s(T,𝐯_s)`$ at $`T=0`$ At extremely low $`T`$, the perturbation expansions of the free energy as power series in $`v_s(𝐫)k_F`$ and $`𝒟_i^{}(𝐫)`$ by expanding $`\stackrel{~}{W}_{n𝐤}(𝐫)`$ before integrating over $`𝐤`$ breaks down due to the existence of the zero modes in $`E_𝐤`$. The correct approach is to expand the free energy after doing the integrals over $`𝐤`$. For simplicity, we only show the $`T=0`$ results. It is easy to find that the driving term for $`𝒟_s^{}`$ is negligibly small compared with that for $`𝒟_s^{\prime \prime }`$. We may simply set $`𝒟_s^{}=0`$. The local free energy in Eq. (A4) becomes $`F_{T=0}^{(loc)}=F_0{\displaystyle d^2𝐫\frac{d^2𝐤}{(2\pi )^2}\sqrt{\stackrel{~}{E}_𝐤^2(𝐫)+\eta _𝐤(𝐫)}}+2{\displaystyle d^2𝐫\frac{d^2𝐤}{(2\pi )^2}Y_𝐤(𝐫)\mathrm{\Theta }(Y_𝐤(𝐫))}`$ (A24) where $`Y_𝐤(𝐫)=\sqrt{\stackrel{~}{E}_𝐤^2(𝐫)+[𝒟_s^{\prime \prime }(𝐫)]^2}𝐯_s(𝐫)𝐤_F`$ and $`\mathrm{\Theta }(x)`$ is the Heaviside step function. We obtain $`F_{T=0}^{(loc)}=\mathrm{\Delta }F+d^2𝐫[c_s+\eta _s^{\prime \prime }(T=0,𝐯_s)][\overline{𝒟}_s^{\prime \prime }(𝐫)]^2`$, where $`\mathrm{\Delta }F`$ is the local free energy at $`𝒟_s=0`$, and $`\eta _s^{\prime \prime }(T=0,𝐯_s){\displaystyle \frac{1}{2}}+S_\theta {\displaystyle \frac{v_s(𝐫)k_F}{2\mathrm{\Delta }_0}},`$ (A25) with $`S_\theta =_{l=\pm 1}\left\|\mathrm{cos}(\theta +l\frac{\pi }{4})\right\|`$. From $`F_s^{(1)}`$ in Eq. (A16) we get the driving term for $`𝒟_s^{\prime \prime }`$, from which we find $`L_s^{\prime \prime }(T=0,𝐯_s){\displaystyle \frac{2ϵ_Fm[v_{sx}^2v_{sy}^2]}{\mathrm{\Delta }_0^2}}\left[1S_\theta {\displaystyle \frac{v_s(𝐫)k_F}{\mathrm{\Delta }_0}}\right].`$ (A26) ### 3 Free energy with respect to $`d_{xy}`$ component We first analyze the driving terms at $`T>𝐯_s𝐤_F`$. The leading-order linear-$`𝒟_{xy}`$ term from $`F^{(loc)}`$ in Eq. (A4) turns out to be $`F_4^{(loc)}{\displaystyle d^2𝐫\mathrm{\Delta }_0^2(T)N_0Q_3^{}(T)U_3^{}(𝐯_s)\overline{𝒟}_{xy}^{}},`$ (A27) where $`Q_3^{}(T)`$ $`=`$ $`N_0^1T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\mathrm{sin}^24\phi \left\{\frac{5}{2W_{n𝐤}^3}+\frac{10E_𝐤^2}{W_{n𝐤}^4}\frac{8E_𝐤^4}{W_{n𝐤}^5}\right\}},`$ (A28) $`U_3^{}(𝐯_s)`$ $`=`$ $`4ϵ_F^2m^2\left[v_{sx}^2(𝐫)v_{sy}^2(𝐫)\right]v_{sx}(𝐫)v_{sy}(𝐫).`$ (A29) For $`\xi <r<R_H`$, $`F_4^{(loc)}`$ is of order $`1/r^4`$. As for $`F_{xy}^{(1)}`$ in Eq. (A16), the two terms in the right-hand side are qualitatively different. We leave the discussion on the second term for a bit later. The first term is nonzero in the leading order of $`(mv_s^2)_xv_{sx}`$ which scales as $`1/r^4`$ too. For $`T>v_sk_F`$, this term is $`F_{xy}^{(11)}{\displaystyle d^2𝐫\mathrm{\Delta }_0^2(T)N_0Q_1^{\prime \prime }(T)U_1^{\prime \prime }(𝐯_s)\overline{𝒟}_{xy}^{\prime \prime }},`$ (A30) where $`Q_1^{\prime \prime }(T)`$ $`=`$ $`N_0^1T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\mathrm{sin}^24\phi \left\{\frac{5}{W_{n𝐤}^3}+\frac{6E_𝐤^2}{W_{n𝐤}^4}\right\}},`$ (A31) $`U_1^{\prime \prime }(𝐯_s)`$ $`=`$ $`ϵ_F^2\left\{m\left[v_{sx}^2(𝐫)v_{sy}^2(𝐫)\right]\left[_xv_{sy}(𝐫)+_yv_{sx}(𝐫)\right]+4mv_{sx}(𝐫)v_{sy}(𝐫)_xv_{sx}(𝐫)\right\}.`$ (A32) A simple analysis shows that there are terms coming from $`F^{(2)}`$ and $`F^{(3)}`$ in Eq. (A9) which are of the same order. After some algebra, we find that the terms of order $`1/r^4`$ in $`F^{(2)}`$ are $`F_{xy}^{(2)}`$ $``$ $`{\displaystyle d^2𝐫\mathrm{\Delta }_0^2(T)N_0\left\{Q_1^{}(T)U_1^{}(𝐯_s)+Q_2^{}(T)U_2^{}(𝐯_s)\right\}\overline{𝒟}_{xy}^{}},`$ (A33) where $`Q_1^{}(T)`$ $`=`$ $`N_0^1T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\mathrm{sin}^24\phi \left\{\frac{3}{W_{n𝐤}^3}\frac{4(E_𝐤^2+3ϵ_𝐤^2)}{W_{n𝐤}^4}+\frac{16ϵ_𝐤^2E_𝐤^2}{W_{n𝐤}^5}\right\}},`$ (A34) $`Q_2^{}(T)`$ $`=`$ $`N_0^1T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\mathrm{sin}^24\phi \left\{\frac{2}{W_{n𝐤}^3}+\frac{24ϵ_𝐤^2}{W_{n𝐤}^4}\frac{32ϵ_𝐤^2E_𝐤^2}{W_{n𝐤}^5}\right\}},`$ (A35) $`U_1^{}(𝐯_s)`$ $`=`$ $`ϵ_F^2\left\{v_{sx}(𝐫)\left[_x^2v_{sy}(𝐫)+3_x_yv_{sx}(𝐫)\right]+v_{sy}(𝐫)\left[_y^2v_{sx}(𝐫)+3_x^2v_{sx}(𝐫)\right]\right\},`$ (A36) $`U_2^{}(𝐯_s)`$ $`=`$ $`ϵ_F^2[_xv_{sx}(𝐫)][_xv_{sy}(𝐫)+_yv_{sx}(𝐫)],`$ (A37) and in $`F^{(3)}`$ are $`F_{xy}^{(3)}`$ $``$ $`{\displaystyle d^2𝐫\mathrm{\Delta }_0^2(T)N_0Q_2^{\prime \prime }(T)U_2^{\prime \prime }(𝐯_s)\overline{𝒟}_{xy}^{\prime \prime }},`$ (A38) where $`Q_2^{\prime \prime }(T)`$ $`=`$ $`N_0^1T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\mathrm{sin}^24\phi \left\{\frac{1}{2W_{n𝐤}^3}\frac{2ϵ_𝐤^2}{W_{n𝐤}^4}\right\}},`$ (A39) $`U_2^{\prime \prime }(𝐯_s)`$ $`=`$ $`ϵ_F^2\left[{\displaystyle \frac{1}{m}}_x^3v_{sy}(𝐫){\displaystyle \frac{1}{m}}_y^3v_{sx}(𝐫)+{\displaystyle \frac{6}{m}}_x^2_yv_{sx}(𝐫)\right].`$ (A40) The sum of $`F_4^{(loc)},F_{xy}^{(11)},F_{xy}^{(2)}`$, and $`F_{xy}^{(3)}`$ leads to $`F_4(𝐫)`$ term in Eq. (38), with $`L_{xy}^{}(T,𝐯_s)`$ $`=`$ $`Q_1^{}(T)U_1^{}(𝐯_s)+Q_2^{}(T)U_2^{}(𝐯_s)+Q_3^{}(T)U_3^{}(𝐯_s),`$ (A41) $`L_{xy}^{\prime \prime }(T,𝐯_s)`$ $`=`$ $`Q_1^{\prime \prime }(T)U_1^{\prime \prime }(𝐯_s)+Q_2^{\prime \prime }(T)U_2^{\prime \prime }(𝐯_s).`$ (A42) The second term on the right-hand side of Eq. (A16) is nonvanishing by noting that $`[_𝐤\mathrm{\Delta }_𝐤_𝐫][𝐯_s(𝐫)k_F]=\mathrm{\Delta }_0\mathrm{sin}2\phi [_𝐫\times 𝐯_s(𝐫)]_z`$ which picks up the magnetic vector potential $`𝐀(𝐫)`$ in $`𝐯_s(𝐫)`$. This term is nothing but $`F^{OZ}`$ in Eq. (37), with $`Q^{OZ}(T)=N_0^1T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}2ϵ_𝐤\mathrm{sin}^22\phi Z_{n𝐤}(𝐯_s)}N_0^1T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\frac{2ϵ_𝐤\mathrm{sin}^22\phi }{W_{n𝐤}^2}}.`$ (A43) $`F^{OZ}`$ is very different from $`F_4`$. Discussions about comparative importance of these two terms can be found in Sec. III.C. The integrand of $`Q^{OZ}`$ includes odd power in $`ϵ_𝐤`$, implying nonzero contribution only in a particle-hole asymmetric system. In order to proceed further, we adopt the following density of states near the Fermi surface in the normal state, $`N(ϵ)N_0+ϵ{\displaystyle \frac{dN(ϵ)}{dϵ}}\|_{ϵ=0}N_0(1+\alpha {\displaystyle \frac{ϵ}{ϵ_F}}),`$ (A44) to take into account both the particle-hole symmetric and asymmetric contributions. Here $`\alpha `$ is of order unity, and $`ϵ`$ is typically order of $`T`$ or $`\mathrm{\Delta }_0`$, implying that the second term on the right-hand side of Eq. (A44) may be negligible when the contribution from the particle-hole symmetric mode does not vanish. The quadratic free energy terms in Eq. (36) are extracted from $`F^{(loc)}`$, where $`\eta _{xy}^{}(T,𝐯_s)=\eta _{xy}^{\prime \prime }(T,𝐯_s)+2d^2h_{2dxy}(d),`$ (A45) $`\eta _{xy}^{\prime \prime }(T,𝐯_s)=h_{3xy}(d)+{\displaystyle \frac{ϵ_Fmv_{sx}^2(𝐫)}{T^2}}h_{4xy}(d),`$ (A46) with $`h_{idxy}`$ defined in Eqs. (A22) and (A23). ## B Spontaneous magnetization in the $`d_{x^2y^2}+id_{xy}`$ state To proceed with a general derivation of the spontaneous magnetization, we rewrite the free energy in Eq. (20) as $`F=F_0T{\displaystyle \underset{n}{}}\mathrm{Tr}\mathrm{ln}\left(\widehat{G}^1\right)T{\displaystyle \underset{n}{}}\mathrm{Tr}\mathrm{ln}\left[1\widehat{G}\left(\begin{array}{cc}\widehat{V}_1& 0\\ 0& \widehat{V}_2\end{array}\right)\right],`$ (B3) where $`\widehat{G}=(\widehat{M}_0+\widehat{D})^1`$ with $`\widehat{D}`$ the field-induced off-diagonal pairing order parameter matrix. The orbital Zeeman term coming from the linear-$`\widehat{V}`$ term is $`F^{OZ}`$ $``$ $`T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤d^2𝐤_1}{(2\pi )^4}d^2𝐫d^2𝐫_1e^{i(𝐤_1𝐤)(𝐫_1𝐫)}\left[G_{11𝐤,𝐤_1}(𝐫)+G_{22𝐤,𝐤_1}(𝐫)\right]\left[(𝐫_1𝐫)_𝐫\right][𝐯_s(𝐫)𝐤_F]}`$ (B4) $``$ $`T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}d^2𝐫\mathrm{lim}_{𝐤_1𝐤}\left\{i_𝐤\left[G_{11𝐤,𝐤_1}(𝐫)+G_{22𝐤,𝐤_1}(𝐫)\right]_𝐫\right\}[𝐯_s(𝐫)𝐤_F]}`$ (B5) $``$ $`{\displaystyle d^2𝐫\frac{e}{2mc}T\underset{n}{}\frac{d^2𝐤}{(2\pi )^2}\mathrm{lim}_{𝐤_1𝐤}\left\{(i𝐤\times _𝐤)\left[G_{11𝐤,𝐤_1}(𝐫)+G_{22𝐤,𝐤_1}(𝐫)\right]\right\}𝐁(𝐫)}.`$ (B6) Comparing Eq. (B6) with Eq. (58) we see that the magnetization is just the expectation value of the magnetic moment operator $`(e/2mc)(i𝐤\times _𝐤)`$, $`𝐌(𝐫)={\displaystyle \frac{e}{2mc}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\mathrm{Tr}\widehat{d}_𝐤^{}(i𝐤\times _𝐤)\widehat{d}_𝐤}={\displaystyle \frac{e}{2mc}}T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\mathrm{lim}_{𝐤_1𝐤}(i𝐤\times _𝐤)\mathrm{Tr}\widehat{G}_{𝐤,𝐤_1}(𝐫)},`$ (B7) with $`\widehat{d}_𝐤^{}=(c_𝐤,c_𝐤^{})`$. It is easy to check that $`𝐌=0`$ for a pure $`d_{x^2y^2}`$\- ($`d_{xy}`$-) wave superconductor, and also for a $`d_{x^2y^2}+is`$-wave one in the order of linear $`s`$. For a $`d_{x^2y^2}+id_{xy}`$-wave superconductor with $`d_{xy}`$ component perturbatively small compared to the $`d_{x^2y^2}`$ component, we use the Dyson equation, $`\widehat{G}_{𝐤,𝐤_1}\widehat{g}_𝐤\delta (𝐤𝐤_1)+\widehat{g}_𝐤\left(\begin{array}{cc}0& 𝒟_{xy}(𝐫)\mathrm{\Phi }_{xy\frac{𝐤+𝐤_1}{2}}\\ 𝒟_{xy}^{}(𝐫)\mathrm{\Phi }_{xy\frac{𝐤+𝐤_1}{2}}& 0\end{array}\right)\widehat{g}_{𝐤_1},`$ (B10) where $`\widehat{g}_𝐤`$ is defined in Eq. (28). Inserting Eq. (B10) into Eq. (B7) yields $`𝐌(𝐫)`$ $`=`$ $`{\displaystyle \frac{e}{2mc}}iT{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\mathrm{\Phi }_{xy𝐤}\left[𝒟_{xy}(𝐫)g_{4𝐤}+𝒟_{xy}^{}(𝐫)g_{1𝐤}\right](𝐤\times _𝐤)g_{2𝐤}}`$ (B11) $`=`$ $`{\displaystyle \frac{e}{2mc}}T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}𝒟_{xy}^{\prime \prime }(𝐫)\frac{4\mathrm{\Delta }_0ϵ_𝐤\mathrm{sin}^22\phi }{W_{n𝐤}^2}},`$ (B12) which, up to the leading order of $`D_{xy}`$ and $`v_s`$, is consistent with the second term on the right-hand side of Eq. (A16). A particle-hole asymmetric system is required to obtain a nonvanishing $`𝐌(𝐫)`$. To understand the physics, we note $`𝐌(𝐫)`$ in Eq. (B11) is in fact $`𝐌(𝐫)`$ $`=`$ $`{\displaystyle \frac{e}{2mc}}{\displaystyle \underset{n,𝐤}{}}2𝒟_{xy}^{\prime \prime }\left\{(u_𝐤^2v_𝐤^2){\displaystyle \frac{1}{i\omega _n+E_𝐤}}(u_𝐤^2v_𝐤^2){\displaystyle \frac{1}{i\omega _nE_𝐤}}\right\}(𝐤\times _𝐤)g_{2𝐤}.`$ (B13) with $`u_𝐤^2=(1+ϵ_k/E_𝐤)/2`$ and $`v_𝐤^2=(1ϵ_k/E_𝐤)/2`$ measuring the particle and hole populations, respectively. Eq. (B13) can be interpreted as a magnetic moment contributed from number currents of both particles and holes flowing in the opposite directions, which cancel in a particle-hole symmetric system. The spontaneous magnetization $`M(𝐫)`$ as a full expression of $`\overline{𝒟}_{xy}^{\prime \prime }`$ can also be obtained from Eq. (B7). After some straightforward algebra we get $`𝐌(𝐫){\displaystyle \frac{eT}{2mc}}{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2𝐤}{(2\pi )^2}\frac{4\mathrm{\Delta }_0(T)𝒟_{xy}^{\prime \prime }(𝐫)ϵ_𝐤\mathrm{\Phi }_{xy𝐤}^2}{W_{n𝐤}\{W_{n𝐤}+[𝒟_{xy}^{\prime \prime }(𝐫)]^2\mathrm{\Phi }_{xy𝐤}^2\}}}.`$ (B14) $`𝐌(𝐫)`$ in Eq. (B14) is important for the direct comparison with Laughlin’s free energy as discussed in Sec. V.
warning/0003/hep-th0003225.html	ar5iv	text	# 1 Introduction ## 1 Introduction The dilaton $`\varphi `$ and the axion $`\chi `$ are scalar fields predicted by superstring theory . Both of them arise in a natural way in the massless spectrum of 10 dimensional (10d) type IIB superstring theory and its lower dimensional compactifications. In the language of 4d gauge theory of field strength $`F_{\mu \nu }`$ and its dual $`\stackrel{~}{F}_{\mu \nu }`$, the $`\varphi `$ and $`\chi `$ have very special couplings. The dilaton $`\varphi `$ couples to gauge fields through a term $`exp(\varphi )F^2`$ and the axion $`\chi `$ couples to the topological term. The $`\varphi `$ and $`\chi `$ fields play a central role in superstring dualities , F-theory compactifications and in the derivation of the exact results in 4d N=2 supersymmetric gauge theories . Recently it was observed in that a string inspired coupling of a dilaton $`\varphi `$ to 4d $`SU(N_c)`$ gauge fields $`A_\mu =T^aA_\mu ^a`$, with $`T^a`$ the $`(N_c^21)`$ $`SU(N_c)`$ generators, yields a phenomenologically interesting potential V(r) for the quark-quark interactions. Following , this potential is obtained as follows: First we start from the following model for the scalar field-gluon coupling $$L(\varphi ,A)=\frac{1}{4G(\varphi )}F_{\mu \nu }^aF_a^{\mu \nu }\frac{1}{2}(_\mu \varphi )^2+W(\varphi )+J_\mu ^aA_a^\mu .$$ (1) Then we choose $`G(\varphi )`$, the coupling of the scalar field $`\varphi `$ to the $`SU(N_c)`$ field strength $`F_{\mu \nu }`$, and the interacting lagrangian $`W(\varphi )`$ as: $$G(\varphi )=const.+\frac{f^2}{\varphi ^2}$$ $$W(\varphi )=\frac{1}{2}m^2\varphi ^2$$ (2) where the parameter $`f`$ is a scale characterizing the strength of the scalar-gluon coupling and m is the mass of the scalar field $`\varphi `$. Next we consider the equations of motion of the fields $`A_\mu `$ and $`\varphi `$ and solve them for static point like color source of current density $`J_a^\mu =\rho _a\eta ^{\mu 0}`$. After some straightforward algebra, we find that the Dick quark interaction potential $`V_D(r)`$ is given by: $$V_D(r)=\frac{1}{r}f\sqrt{\frac{N_c}{2(N_c1)}}\mathrm{ln}[exp(2mr)1]$$ (3) Eq.(3) is very remarkable since for large values of $`r`$ it leads to a confining potential $`V_D(r)2fm\sqrt{\frac{N_c}{2(N_c1)}}r`$. In this regard, we will show in this paper that for a general gluon-dilaton coupling $`G(\varphi )`$, the quark interactions potential $`V(r)`$ reads as: $$V(r)=𝑑r\frac{G[\varphi (r)]}{r^2}.$$ (4) Such form of the potential is very attractive. On one hand it extends the usual Coulomb formula $`V_c1/r`$ which is recovered from eq.(4) by taking $`G=1`$. Moreover for $`Gr^2`$, which by the way corresponds to a coupling $`G(\varphi )\varphi ^2`$, and $`W(\varphi )=\frac{m^2}{2}\varphi ^2`$, $`m0`$, eq.(4) yields a linearly increasing interquark potential $`Vr`$ having the good behaviour to describe the $`SU(N_c)`$ quarks confinement . On the other hand eq.(4) may be also used to describe other non perturbative effects associated with higher dimension quark and gluon vacuum condensates. Following , see also , one may extract interesting phenomenological informations on the dilaton-gluon coupling $`G[\varphi ]`$ by comparing eq.(4) to the Bian-Huang-Shen’s potential $`V_{BHS}(r)`$ namely: $$V_{BHS}(r)\frac{1}{r}\underset{n0}{}C_nr^n$$ (5) where $`C_n`$’s are related to the quark and gluon vacuum condensates. In fact one can do better if one can put the coupling $`G(\varphi )`$ in the form $`G[\varphi (r)]`$. In this case one can predict the type of vacuum condensates of the $`SU(N_c)`$ gauge theory which contributes to the quark-quark interaction potential . Thus, although, the derivation of the formula (4) for the interquark potential from eq.(1) is by itself an important result, there remain however other steps to overcome before one can exploit eq.(4). As mentioned above, a crucial step is to determine what type of couplings $`G(\varphi )`$ which can be put in the form $`G[\varphi (r)]`$. In other words for what couplings $`G(\varphi )`$, one can solve the equation of motion of the scalar field $`\varphi `$. This is a technical problem without solving it one cannot get V(r). An other step which remains to clarify is to show how the effective model (1) may be got from a more fundamental theory. If this is possible, one may for instance justify the mass scale $`f`$ introduced by hand in eqs.(2,3). One might also get some informations on the axion field couplings and more generally on the moduli of 10d superstrings compactified on six dimensional compact manifolds and especially type IIB on Calabi-Yau Threefolds . In trying to explore eq.(4), we have observed some remarkable facts among which we quote the three following : 1) The functional $`G[\varphi (r)]`$, and then the potential V(r) of eq.(4) may be obtained from the following one dimensional lagrangian $`L_D`$ $$L_D=\frac{1}{2}(y^{})^2+r^2W(y/r)+\frac{\alpha }{2r^2}G(y/r)$$ (6) where $`y=r\varphi `$, $`y^{}=(\frac{dy}{dr})`$and $`\alpha =\frac{g^2}{16\pi ^2}\frac{N_c1}{2N}`$ and where $`g`$ is the gluon coupling constant. In particular $`L_D`$ reads, for $`W(\varphi )`$ and $`G(\varphi )`$ like in eq(2),as: $$2L_D=(y^{})^2+m^2y^2+\frac{\mu ^2}{y^2}$$ (7) where $`\mu =\alpha f^2`$. 2) Eq.(7) has a striking resemblance with the following harmonic superspace lagrangian $`L_{EH}`$ used in in the derivation of the 4d Eguchi-Hanson metric $$2L_{EH}=(D^{++}\omega )^2+m_{}^{++}{}_{}{}^{2}\omega ^2+\frac{\mu _{}^{++}{}_{}{}^{2}}{\omega ^2}.$$ (8) In this equation, $`\omega `$ is an analytic harmonic superspace (hs) superfield taken to be dimensionless, $`D^{++}`$ is the hs covariant derivative and $`m^{++}`$ and $`\mu ^{++}`$ are coupling constants. More details on hs tools will be described in section 3. Much more precision can be found in . For the moment note only the formal analogy between y, dy/dr, m and $`\mu `$ of eq.(7) with $`\omega `$, $`(D^{++}\omega )`$, $`\mu ^{++}`$ and $`m^{++}`$ respectively. Both of models eqs.(7) and (8), involve hermitian fields with a self interacting potential proportional to the inverse of the square of the scalar field variable. 3) The Dick potential (3) is viable only for non zero mass dilaton field exactly as in 4d N=2 supersymmetric theories where the scalar potential is proportional to the mass eigenvalues of the central charges of the 4d N=2 superalgebra . Recall by the way that in 4d N=2 supersymmetric QFT, mass terms are generated by central charges. We shall see in sections 3 and 4 that this formal analogy between the Dick model (1) and 4d N=2 QFT’s is much deeper since it allows us to derive a new model containing eq.(1) and where the symmetries behind the solvability of Dick equations as well as the couplings of both the dilaton and axion fields are manifest. The aim of this paper is to generalize the Dick model (1) by exploiting the formal analogy with 4d N=2 supersymmetric theories formulated in hs and using known 4d N=2 exact results. In addition to the derivation of new model exhibiting a U(1) gauge invariance, we give an interpretation of the mass scale $`f`$, introduced by hand in eq.(2), as a Kahler moduli of a blown up SU(2) singularity of Calabi-Yau threefold of type II superstring compactifications. The appearance of the local U(1) symmetry in the analysis of eqs.(1-3) has a quite interesting consequence as it offers a possibility to incorporate in the game the axion field $`\chi `$ couplings. Recall that in Dick model as formulated in , the role of the topological field $`\chi `$ is ignored. We shall show in section 4 how this field can be incorporated by going in a general gauge other than $`\varphi =\varphi ^{}`$. The presentation of this paper is as follows: In section 2, we formulate the Dick problem as a one dimensional field theory. In section 3 we give general solutions including those of . In section 4, we review briefly the building of the Eguchi Hanson hyperkahler metric in harmonic superspace. In section 5 we use the formal analogy between the Eguchi Hanson model and our one dimensional field theoretical formulation of the Dick problem to determine the dilaton couplings, the axion ones and interprete the mass scale $`f`$ as a kind of Fayet-Iliopoulos coupling. Our conclusion is given in section 6. ## 2 The Dick model as a one dimensional field theory. Following ref., the analysis of the Coulomb problem of the theory (1) is based on considering a point like static color source which in its rest frame is described by a current $`J_a^\mu =g\delta (r)C_a\eta _0^\mu `$ where $`C_a`$ is the expectation value of the $`SU(N_c)`$ generator for a normalized spinor in the color space. These $`C_a`$’s satisfy the algebraic identity $$\underset{a=1}{\overset{N^21}{}}C_a^2=\frac{(N_c1)}{2N_c}$$ (9) The next step is to use the residual SO(3) space symmetry, which remains after setting $`J_a^\mu =\rho _a\eta _0^\mu `$, to rewrite the equations of motion $$[D_\mu ,G^1(\varphi )F^{\mu \nu }]=J^\nu $$ $$_\mu ^\mu \varphi =\frac{W}{\varphi }\frac{1}{4}F_{\mu \nu }^aF_a^{\mu \nu }\frac{G^1(\varphi )}{\varphi }$$ (10) into a simple form. Indeed setting $`F_a^{0i}=\frac{gC_a}{4\pi }_iV`$; $`\alpha =\frac{g^2}{16\pi }\frac{(N_c1)}{2N_c}`$ one finds after some easy algebra: $$\frac{dV}{dr}=r^2G[\varphi ](a)$$ (11) $$\mathrm{\Delta }\varphi =\frac{W}{\varphi }+\frac{\alpha }{r^4}\frac{G(\varphi )}{\varphi }(b)$$ Note that eqs.(11) have four unknown field quantities; the field $`\varphi `$, the interacting color potential V(r), the dilaton-gluon coupling $`G(\varphi )`$ and the $`\varphi `$ potential $`W`$. To solve eqs.(11) one has to fix two of them. For example choosing $`2W=m\varphi ^2`$and $`G(\varphi )`$ as in equation (3), one finds : $$\varphi =\varphi _D(r)=r^1[\frac{\alpha f}{m}(1exp(2mr)]^{1/2}$$ (12) $$V_D(r)=\frac{1}{r}f\sqrt{\frac{N_c}{2(N_c1)}}ln[exp(2mr)1].$$ (13) In general given $`G(\varphi )`$ and $`W(\varphi )`$, the color potential V(r) can be exactly determined up on solving one equation namely equation (11-b). For later use, let us introduce the new dimensionless field $`y=r\varphi `$ and take the spherical coordinate frame $`(r,\theta ,\phi )`$ to rewrite the lagrangian (1) as: $$L=\frac{r^2}{2G(\varphi )}F_{0r}^aF_a^{0r}\frac{r^2}{2}_r\varphi ^r\varphi +r^2W(\varphi )+F_{0r}^a\rho _a.$$ (14) In deriving eq.(14), we have used the stationarity of the color source, the SO(3) symmetry and the identity $`\mathrm{\Delta }(1/r)=\delta (r)`$. Putting this equality back into eq.(14) and using the change of variable $`y=r\varphi `$ together with the convention notations $`y^{}=_ry`$; $`^r\varphi =r^2_r\varphi `$ as well as eq.(9), one gets the lagrangian form eq.(6). Consequently the coupling $`G(\varphi )`$ of eq.(1) appears as a part of interacting potential of the one dimensional field theory eq.(6). From this point of view, the finding of the interquark potential V(r) is equivalent to solve the equation of motion $$y^{\prime \prime }\frac{L_D}{y^{}}+y^{}\frac{L_D}{y}+_r^{exp}L_D=0.$$ (15) ## 3 Solving the Dick model First of all observe that the lagrangian (6) including the Dick model (7) is a particular one dimensional field theory of lagrangian $$L=\frac{1}{2}(y^{})^2U(y,r)$$ (16) where $`U(y,r)`$ is a priori an arbitrary potential. Though simple, this theory is not easy to solve except in some special cases. A class of solvable models is given by potentials of the form : $$U(y)=\lambda ^2y^{2(n+p)}+\gamma ^2y^{2(qn)}+\delta y^k$$ (17) where $`n,p,q`$ and $`k`$ are numbers and $`\lambda ^2`$ , $`\gamma ^2`$ and $`\delta `$ are coupling constants scaling as $`(lenght)^2`$. The next thing to note is that eq.(17) has no explicit dependence in $`r`$ and consequently the following identity usually hold : $$y^2=U+c$$ (18) where c is a constant. Actually eq.(18) is just an integral of motion which may be solved under some assumptions. Indeed by making appropriate choices of the coupling $`\lambda `$ as well as the integral constant c , one may linearise y’ in eq.(18) as follows : $$y^{}=U_1+U_2.$$ (19) Once the linearisation in y’ is achieved and the terms $`U_1`$ and $`U_2`$ are identified, we can show that the solutions of eq.(18) are classified by the product $`U_1U_2`$ and the ratio $`U_1/U_2`$. In what follows we discuss briefly some interesting examples. For convenience let us rewrite eq.(18) as: $$y^2=w_0+w_1+C_0$$ (20) where $`w_0=m^2y^2`$ and $`w_1`$ is the interaction term which we take for the moment to be the Dick interaction that is $`w_1=c_1^2y^2`$, where $`c_1`$ is a coupling constant. Starting from eq.(20), it is not difficult to see that there are two possibilities to put it in the form (19): ### 3.1 First possibility: the Dick solution This corresponds to take $`w_0=U_1^2`$ that is $`U_1=my`$ and $`U_2=c_1y^1`$. Putting back into eq.(19) one gets the Dick solution given by eqs.(12,13). ### 3.2 Second possibility: New solutions In this case the mass term is related to the product $`U_1U_2`$ as : $$U_1U_2=\pm \frac{1}{2}m^2y^2$$ (21) Eq.(21) cannot however determine $`U_1`$ and $`U_2`$ independently as in general the following realizations are all of them candidates, $$U_1=\lambda y^{n+p};U_2=\gamma y^{qn}$$ (22) where the integers $`p`$ and $`q`$ are such that $`p+q=2`$ and where $`\lambda \gamma =\pm m^2`$. A remarkable example corresponds to take $`p+q=1`$. In this case we distinguish two solutions according to the sign of the product of $`\lambda \gamma `$ . For $`\lambda \gamma =+m^2`$, the solution is $$y(r)=[\frac{1}{\lambda }tan(\frac{nmr}{\sqrt{2}}+const)]^{\frac{1}{n}}.$$ (23) For $`\lambda \gamma =m^2`$ , we have: $$y(r)=[\frac{1}{\lambda }tanh(\frac{nmr}{\sqrt{2}}+const)]^{\frac{1}{n}}.$$ (24) The solutions (23) and (24) have quite interesting features inherited essentially from the features of tan and tanh functions. We remark that for n=0 the solution is: $$y(r)=const.exp(\frac{\lambda +\gamma }{\sqrt{2}}r).$$ (25) In the end of this section, it should be noted that one can go beyond the above mentioned solutions which are just special case of general models involving interactions classified by the following constraint equations $$U_1.U_2y^k$$ (26) where $`U_1`$ and $`U_2`$ are as in eq.(19) and $`k`$ is an integer. For $`k=0`$, one gets the Dick model and for $`k=2`$ one has solutions described in subsection 3-2 . For general values of $`k`$, one has to know moreover the ratio $`U_1/U_2`$ in order to work out solutions. For the example where $$U_1=\lambda y$$ $$U_2=\gamma y^{k1};kinteger$$ (27) one can check, after some straightforward algebra, that the solution of y is just a generalization of eq.(12) that is $$y_k(r)=[r\varphi _D]^{\frac{2}{(2k)}}.$$ (28) For $`k=0`$, one discovers the solution (12). ## 4 The Eguchi Hanson hs model To start recall that Eguchi Hanson metric is a vacuum solution of the self dual euclidean four dimensional gravity. It is a Ricci flat hyperkahler metric having an $`SU(2)\times U(1)`$ isometry. There are different, but equivalent, ways of writing this metric. A remarkable way of expressing this metric is that using a local coordinate system exhibiting manifestly the $`SU(2)\times U(1)`$ symmetry. The element of length $`ds^2`$ reads as: $$ds^2=g_{iajb}df_1^{ia}df_1^{jb}+k_{iajb}df_2^{ia}df_2^{jb}+h_{iajb}df_2^{ia}df_1^{jb}$$ (29) where the metric factors are given by: $$g_{iajb}=ϵ_{ab}ϵ_{ij}\frac{4f_{2ia}f_{2jb}}{f_1^{kc}f_{1kc}+f_2^{kc}f_{2kc}}(a)$$ $$k_{iajb}=ϵ_{ab}ϵ_{ij}\frac{4f_{1ia}f_{1jb}}{f_1^{kc}f_{1kc}+f_2^{kc}f_{2kc}}(b)$$ (30) $$h_{iajb}=\frac{4f_{1ia}f_{2jb}}{f_1^{kc}f_{1kc}+f_2^{kc}f_{2kc}}c)$$ together with the SU(2) isovector constraint $$ϵ_{ab}(f_1^{ia}f_2^{jb}+f_1^{ja}f_2^{ib})\lambda ^{ij}=0.$$ (31) A tricky way to derive this metric is to use results of 4d N=2 supersymmetric non linear $`\sigma `$ models. In the harmonic superspace approach where 4d N=2 supersymmetry is manifest, the field theoretical model giving the family of Eguchi Hanson metrics reads in the superfield language as : $$S[\omega ]=\frac{1}{2k^2}𝑑z^{(4)}𝑑u[(D^{++}\omega )^2m^{++2}\omega ^2\frac{\lambda ^{++2}}{\omega ^2}]$$ (32) In this equation $`\omega =\omega (x_A,\theta ^+,\overline{\theta }^+,u)`$, is an analytic hs superfield taken to be dimensionless. $`D^{++}=(u^{+i}\frac{}{u^i}2\theta ^+\sigma ^m\overline{\theta }^+_m)`$ is the hs covariant derivative; $`dz^4`$ is the analytic superspace measure with U(1) Cartan charge (-4) and the couplings $`m^{++}`$ and $`\lambda ^{++}`$ are given by $$m^{++}=u_i^+u_j^+m^{ij};\lambda ^{++}=u_i^+u_j^+\lambda ^{ij}$$ (33) where $`u_i^+`$ and $`u_i^{}`$ are the harmonic variables parameterizing the $`SU(2)/U(1)S^2`$ sphere. We shall not use here after these hs tools, we are only interested in the formal analogy with the Dick problem. This is why we shall give here after only the necessary material. For more details on the HS method and the derivation of the Eguchi Hanson metric see . Note also that the Eguchi Hanson metric with $`SU(2)\times U(1)`$ isometry corresponds to $`m^{++}=0`$. Metrics with $`m^{++}0`$ have a $`U(1)\times U(1)`$ symmetry and fall in the family of multicenter metrics . Let us take $`m^{++}=0`$ and sketch the main steps in putting eq.(32) in the form (29-31). In fact there are two possible paths one may follow: First, a direct method which consists to start from the superfield equation of motion of the hermitian hs superfield $`\omega `$: $$D^{++2}\omega =\frac{\lambda ^{++2}}{\omega ^3}$$ (34) and use the $`\theta `$expansion of the superfield $`\omega `$, that is $$\omega =\varphi +\theta ^{+2}M^{(2)}+\overline{\theta }^{+2}\overline{N}^{(2)}+\theta ^+^m\overline{\theta }^+B_m^{(2)}+\theta ^{+2}\overline{\theta }^{+2}P^{(4)};$$ (35) where we have ignored fermions. Then fix N =2 supersymmetry partially on shell by eliminating the auxiliary fields $`P^{(4)}`$ and $`B_m^{(2)}`$. The relevant equations are those corresponding to the projection of eq.(34) along the $`\theta ^+=0`$ and $`\theta ^+\sigma ^m\overline{\theta }^+`$ directions, i.e. $$^{++2}\varphi =\lambda ^{++2}/\varphi ^3$$ $$^{++}B_m^2=2(_m\frac{3}{2}\frac{\lambda ^{++2}}{\varphi ^4}B_m^{(2)})\varphi .$$ (36) The next thing to do is to find the explicit dependence of $`\varphi `$ and $`B_m^{(2)}`$ in harmonic variables $`u_i^\pm `$ by solving eqs.(36). Then put the solution into eq.(32) once the integrations with respect to $`\theta ^+`$ and $`\overline{\theta }^+`$ are performed. In other words put the solutions $`\varphi =\varphi (u_i^\pm )`$ ,$`B_m^{(2)}=B_m^{(2)}(u_i^\pm )`$ into the following component field action: $$S[\omega ]\frac{1}{k^2}𝑑x^4𝑑u[^{++}B_m^{(2)}^m\varphi +^mB_m^{(2)}^{++}\varphi ].$$ (37) The last step is to integrate with respect to harmonic variables. Once this is done, we get the bosonic part of the 4d N=2 supersymmetric non linear $`\sigma `$ model from which one can read the Eguchi Hanson metric in the $`\omega `$ representation. The second method, which interest us here, is indirect but it has the merit of being based on hs superfield theory exhibiting manifestly the $`SU(2)\times U(1)`$ symmetry. The main steps of this approach are as follows: 1. Instead of working with a real superfield $`\omega `$, we take a complex superfield $`\omega `$ : $`\overline{\omega }\omega `$ . 2. Modify the action eq.(32) as: $$S[\omega ]\frac{1}{2k^2}𝑑z^{(4)}𝑑u[\|(D^{++}+iV^{++})\omega \|^2+\lambda ^{++}V^{++}],$$ (38) where $`V^{++}`$ is a U(1) gauge superfield. Eq.(38) is invariant under the following gauge transformations of parameter $`\lambda `$. $$\omega ^{}=\mathrm{exp}(i\lambda )\omega ;V^{++}=V^{++}+D^{++}\lambda .$$ (39) Note that $`V^{++}`$ has no kinetic term. It is an auxiliary superfield which can be eliminated through its equation of motion namely $$2V^{++}=\frac{1}{\omega \overline{\omega }}[i(\overline{\omega }D^{++}\omega \omega D^{++}\overline{\omega })\lambda ^{++}].$$ (40) For the special case where $`\omega `$ is real; $`\overline{\omega }=\omega `$ , eq.(40) reduces to $$V^{++}=\lambda ^{++}/\omega ^2$$ (41) and consequently the action (38) coincides with eq.(32). Note by the way that the term $`\lambda ^{++}V^{++}`$ is a Fayet-Iliopoulos (FI) coupling. 3. Rewrite eq.(38) in an equivalent form by using O(2) notations i.e. express the complex supefield $`\omega =\omega _1+i\omega _2`$ as an O(2) doublet $`(\omega _1,\omega _2)`$ and introducing two other auxiliary superfield $`F_1^{++}`$ and $`F_2^{++}`$ $$S[\omega _1,\omega _2,F_1^{++},F_2^{++}]=\frac{1}{2k^2}dz^{(4)}du[(F_1^{++})^2+2F_1^{++}D^{++}\omega _1+(12)$$ (42) $$V^{++}(\omega _1F_2^{++}\omega _2F_1^{++}+\lambda ^{++})]$$ Eliminating $`V^{++}`$,$`F_1^{++}`$ and $`F_2^{++}`$ and choosing the gauge $`\omega _2=0`$ we reproduce the action $`S_{EH}`$ (32) with $`m^{++}=0`$. The second order action(42) is interesting since it has a manifest SU(2) invariance rotating $`\omega _i`$ and $`F_i^{++}`$ . To make this invariance more explicit we make the following change for both $`(\omega _1,F_1^{++})`$ and $`(\omega _2,F_2^{++})`$. $$\omega =U_a^{}q^{+a};F^{++}=U_a^+q^{+a}$$ $$q^{+a}=ϵ^{ab}q_b^+;q_a^+=(q^+,\overline{q}^+);ϵ^{12}=1.$$ (43) Thus for both $`\omega _1`$ and$`\omega _2`$, we have $`\omega =U_a^{}q_I^{+a}`$ with $`I=1,2`$ and so on. Putting back into eq.(42), we get the following action, $$S=\frac{1}{2k^2}dz^{(4)}du[\overline{q}_1^+D^{++}q_1^++\overline{q}_2^+D^{++}q_2^+++V^{++}(\overline{q}_1^+q_2^++\overline{q}_2^+q_1^++\lambda ^{++})].$$ (44) This action has the invariance under the following groups: i) O(2) gauge group acting as $$\delta q_I^+=ϵ_{IJ}\lambda q_J^+;\delta V^{++}=D^{++}\lambda ;\overline{\lambda }=\lambda .$$ (45) ii) U(1) subgroup of the rigid SU(2) automorphism group of supersymmetry that leaves $`\lambda ^{++}`$ invariant. iii) The SU(2) Pauli Cursey symmetry rotating $`q_I^+`$ and $`ϵ_{IJ}\overline{q}_J^+`$. Now starting from the last form of the E.H. action (42) and solving the $`\theta ^+=0`$ and $`\theta ^+\sigma ^m\overline{\theta }^+`$components of the equations of motion: $$D^{++}q_Iϵ_{IJ}V^{++}q_J^+=0$$ $$ϵ^{IJ}\overline{q}_I^+q_J^++\lambda ^{++}=0$$ (46) in the Wess-Zumino gauge, one gets by following the same lines described for the direct method, the E.H. metric (29-31). ## 5 The Dick model revisited In section 2 we have learnt that the Dick problem may be formulated as a one dimensional field theory of lagrangian $`L_D`$ given by eq.(16) namely: $$2L_D=(\frac{dy}{dr})^2m^2y^2\frac{\mu ^2}{y^2}$$ (47) In section 3 we have showed that hyperkhaler metrics of the Eguchi Hanson family can be derived from the following 4d N=2 supersymmetric model, $$2L_{EH}=(D^{++}\omega )^2m_{}^{++}{}_{}{}^{2}\omega ^2\frac{\mu _{}^{++}{}_{}{}^{2}}{\omega ^2}$$ (48) In section 4 we have seen that this lagrangian is equivalent to the following first order one, once the auxiliary U(1) gauge superfield $`V^{++}`$ is eliminated through its equation of motion: $$2L_{EH}^{}=\|D^{++}\omega \|^2m_{}^{++}{}_{}{}^{2}\omega \overline{\omega }V^{++}(\overline{\omega }D^{++}\omega \omega D^{++}\overline{\omega }\mu ^{++})+V_{}^{++}{}_{}{}^{2}\omega \overline{\omega }$$ (49) This form of $`L_{EH}`$ may be also transformed into two other forms as shown on eq.(42) and eq.(44). The difference between $`L_{EH}`$ and $`L_{EH}^{}`$ is that in eq.(48) $`\omega `$ is hermitian whereas in eq.(49) $`\omega `$ is complex. As we have seen, we can go from eq.(49) to eq.(48) either by constraining the superfield to be real, that is : $$\omega _2=0,$$ (50) Or equivalently by keeping $`\omega _20`$ and working in the Wess-Zumino gauge: $$D^{++}V^{++}=0$$ (51) Eq.(51) turns on to be helpful in the derivation of the Eguchi Hanson metric. Now using the formal analogy between eqs.(47) and (48), it is not difficult to see that the $`L_D`$ lagrangian may be also formulated in terms of the auxiliary fields $`F`$ and $`\overline{F}`$ as: $$L_2=F\overline{F}+\overline{F}y+F\overline{y}+V[\xi +i(y\overline{F}\overline{y}F)]$$ (52) where V is a one dimensional U(1) gauge field and $`\xi `$ is a 1d constant vector breaking explicitly invariance under space translations. Note that in this formulation, the two scalars $`y_1`$ and $`y_2`$ of the complex field $`y=\frac{1}{\sqrt{2}}(y_1+iy_2)`$ represent respectively the dilaton $`\varphi `$ and the axion $`\chi `$ in agreement with the requirement of F-theory and 10d type IIB superstring and 4d N supersymmetric gauge theory. Eliminating the auxiliary fields F and through their equations of motion namely $$F=(+iV)y=y(a)$$ (53) $$\overline{F}=(+iV)\overline{y}=\overline{y}(a)$$ one obtains the following first order lagrangian $`L_1`$ $$L_1=\|y\|^2+m^2y\overline{y}+\xi V.$$ (54) Moreover eliminating the auxiliary U(1) gauge field V through its equation of motion $$V=\frac{1}{2y\overline{y}}[\xi +i(\overline{y}y+y\overline{y})],$$ (55) one gets the one dimensional field theory of the dilaton-axion system extending eq.(7) which may be recovered from eqs.(54,55) by going to the gauge fixing $`y_2=\chi =0`$. However to exhibit the effect of the axion field $`\chi `$, one has to keep $`\chi 0`$ and imposes a constraint on the gauge field $`V`$ that we write as follows: $$C(V,V)=0$$ (56) Using this constraint, the first order lagrangian is no longer invariant under the change $`ye^{i\psi }y`$ and $`VV+\psi `$, where $`\psi `$ is the U(1) gauge parameter, but as a counterpart one can work out a non trivial solution for the axion field $`\chi =\chi (r)`$ by solving the conjugate where the field V should be substituted in equations of motion (55) $`^2y=m^2y`$ and its complex by the value $`V_0(r)`$ verifying the constraint eq.(56) and satisfying the identity $`(\overline{y}yV_0)=0`$. In the end of this study we would like to note that the term $`\xi V`$ appearing in eq.(52) plays a similar role as the Fayet Iliopoulos term $`m^{++}V^{++}`$ of the 4d N=2 supersymmetric Eguchi Hanson model (32) . Thus the mass scale $`f`$ introduced by hand in the Dick model may be viewed, under some assumptions, as the scale of breaking of the U(1) symmetry rotating the dilaton and axion fields. Recall by the way that in general supersymmetric gauge theories with a U(1) gauge invariance, the FI term is generally used to break supersymmetry and/or gauge invariance. The FI couplings are Kahler moduli of Calabi-Yau threefolds on which 10d type II superstrings are compactified on, and their magnitude are of order of the Calabi-Yau compactification scale. ## 6 Conclusion Inspired from the dilaton-gluon coupling in superstring theory, Dick built a field theoretical model having the remarkable property of leading to a confining quark-quark interaction potential. The model is mainly a 4d $`SU(N_c)`$ gauge theory coupled to a massive scalar field $`\varphi `$ of lagrangian (1) and a dilaton-gluon coupling $`G(\varphi )=1+f^2/\varphi ^2`$, where $`f`$ is a mass scale introduced by hand. The parameter $`f`$ may be compared with the mass scale of the $`\sigma `$ model of mesonic theory . The confining phase of Dick model is parameterized by non zero mass for the dilaton and non vanishing $`f`$. In trying to analyze the potential $`V_D(r)`$ we have observed that the Dick problem has a perfect formal analogy with the problem of building the Eguchi-Hanson metric in 4d N=2 supersymmetric harmonic superspace. This formal similarity appears at several levels. In the introduction we have quoted some of these striking analogies. For example, vanishing masses for both Dick and Eguchi-Hanson scalar fields lead to trivial potentials. An other example is that the mass scale $`f`$ which is introduced by hand and interpreted as a compactification scale by Dick plays a similar role as the 4d N=2 FI coupling appearing in the Eguchi Hanson model (32). Recall by the way that now it is well established that the FI couplings are of order of the compactification scale since they are just the Kahler moduli of the Calabi Yau threefold on which the 10d type II superstring are compactified on. To understand the thing behind the striking similarity between the Dick problem and the Eguchi Hanson one, we have reformulated the Dick problem as a one dimensional field theory. As a consequence we have found a general formula for the inter-quark potential V(r) which of course depend on the nature of the dilaton-gluon $`G[\varphi ]`$ as shown on eq.(4). The beauty of this formula is not only because it extends the Coulomb and Dick theory but also because it can be compared with known parameterizations of the confinement especially the contribution of the quark and gluon vacuum condensates. From this point of view, the Dick model as we have formulated it, may be viewed as a phenomenological theory modeling the non-perturbative contributions responsible of confinement. In this regard a more explicit analysis will be presented in . Moreover, having at hand the 1d field theoretical formulation of the Dick model and the analogy with 4d N=2 Eguchi Hanson model, we have shown how the axion field may be incorporated in the game in agreement with the requirement of F-theory and 10d type IIB superstrings according to which the dilaton and the axion form a complex field. In our formulation the dilaton-axion model is represented by a 1d U(1) gauge theory of lagrangian (52). The $`U(1)SO(2)`$ symmetry rotates the dilaton and the axion fields and allows to interpret the Dick mass scale as a kind of FI coupling. This research work has been supported by the program PARS number 372-98 CNR. References 1. \] M. Green, J. Schwartz and E. Witten, Superstring Theory , Cambrige University Press (1987). 2. \] C. Vafa, Lectures on String and Dualities HUTP-79/A009, hep-th/9702201. 3. \] A. Sen, Unification of string dualities , hep-th/9609176. 4. \] C. Vafa, Evidence for F-theory , HUTP-96/A004, hep-th/9602022. 5. \] N. Seiberg and E. Witten , hep-th/9602022; Nucl. Phys. B426(1994)19 and Nucl. Phys. B431(1994) 484. 6. \] R. Dick, ”Confinment from a massive scalar in QCD”, LMU-TPW-98/06, to appear in Eur. Phys. Journal. C6 (1999). Proceedings of the Miniworkshop on ”Recent Developements in Quantum Field Theory”, UFR-PHE, Rabat (October 1998) to appear. 7. \] R. Dick, ”Implication of dilaton in gauge theory and cosmology”, hep-th/9701047, LMU-TPW-97/10. Phys. lett. B397 (1997)193; Phys. lett. B409 (1997)321. 8. \] T. Huang, Z.Huang, Phys. Rev. D39 (1989)1213. 9. \] J.G. Bian, T. Huang and Q.X.Shen, The quarkonium potential within short and intermediate range in QCD , BIHEP-TH-92-38 and references therein. 10. \] S. Katz, P. Mayr and C. Vafa, Mirror summetry and exact solution of 4d N=2 gauge theories I , HUTP-97/A025, OSU-M.97-5, IASSNS-HEP 97/65, hep-th/9706110. 11. \] A. Galperin, E. Ivanov, V. Ogievetsky and P.K. Townsend, Class. Quantum .Grav. 3 (1986) 625. 12. \] A. Galperin, E. Ivanov, S. Kalitzin, V. Ogievetsky and E. Sokatchev, Class. Quantum . Grav. 1 (1985) 469. 13. \] J. Wess and J. Bagger, Supersymmetry and supergravity , Princeton Press (1983). S. Ferrara and C. Savoy, Supergravity 81 (Cambridge U.P; Cambridge 1982). 14. \] E.A. Gibbons and P.J. Ruback, Commun. Math. Phys. 115 (1988) 267. 15. \] G.W. Gibbons, D. Olivier, P.J. Ruback and G. Valent, Nucl . Phys. B296(1988)679. E. Sahraoui and E.H. Saidi, Class. Quantum. Grav. 16 (1999) 1. 16. \] K. Wehrberger, Phys. Rep. 5-6 (1993) 273. 17. \] M. Chabab, R. Markazi and E.H. Saidi, (in progress). 18. \]M.A. Shifman, A.I. Vainshtein and V.I. Zakharov, Nucl. Phys. B147 (1979) 385.
warning/0003/cond-mat0003018.html	ar5iv	text	# Quasiparticle Localization Transition in Dirty Superconductors ## I Introduction The problem of Anderson localization in normal disordered electronic systems has been studied for years, and has continued to pose intriguing puzzles . Recently, the question of whether an analogous transition can occur in the behaviour of low-energy quasiparticle excitations about a superconducting ground state has come to light (see for e.g. Ref..) In fact, the stability of two very distinct superconducting phases which may be characterized by their transport properties - one with extended quasiparticle states at the Fermi energy, and the other with localized states- has been established. Here, it is the transition point between these two phases, and its critical properties, that forms the subject of our attention. While one can afford to ask of the superconducting systems the same questions that have been addressed in normal systems, the former displays refreshingly new physics that shows both conceptual and qualitative differences from the latter. To begin with, quasiparticle excitations in superconducting systems do not conserve charge, and thus one cannot study the transition through charge dynamics. This absence of $`U(1)`$ symmetry in the Bogoliubov deGennes Hamiltonian, which we employ to describe the excitations in question, has marked consequences. As emphasized by Ref., it gives rise to a host of new universality classes ,with critical exponents whose values are significantly different from those of their normal partners. In particular, the density of states(DOS) exhibits astonishing features. In this paper, we explore these surprises in the context of superconducting systems that respect spin-rotational ($`SU(2)`$) invariance and time-reversal (T) symmetry. The issue of a delocalization-localization transition within the superconducting state has been a recurrent theme (see for e.g. Ref). But it is only in recent years that analyses of both microscopic and continuum models have paid focused attention from a variety of different avenues, and delved into the prospect of making this transition realizable in physical systems. To repeat the example offered of the dirty d-wave superconductor, consider a system of impure superconducting sheets with d-wave pairing coupled to one another. At the nodal points, one has low-energy quasiparticle excitations, and in fact, due to disorder, one even expects a finite DOS at the Fermi energy. For low interplane coupling strength or high impurity concentration, one would expect these states to be localized, and upon increasing the coupling or lowering the disorder, one could conceive of accessing a critical point beyond which these states become extended. Here, we study such a transition on more generic grounds by exploiting the field-theoretic set-up offered by Ref., and through numerical analyses, both of which serve to bring out the novel features of superconducting systems quite dramatically. We review the framework used to describe quasiparticle excitations in the absence of interactions. We expand on some details of the two phases, the ’thermal metal’ and the ’thermal insulator’. Then, in order to keep our paper self-contained, we elaborate on the field-theoretic and numerical methods. We proceed to discuss the field-theoretic and numerical predictions for the localization length exponent $`\nu `$ associated with the thermal metal-thermal insulator transition, and show that they both confirm the existence of a new universality class with $`\nu <\nu _n`$, where $`\nu _n`$ is the corresponding exponent for normal systems. We then study the unique properties of the DOS at criticality. Finally, we mention the characteristics of transitions in superconducting systems besides those that respect $`SU(2)`$ and T symmetry, and possibilities for experiment. ## II Models and Symmetries ### A The Superconducting Hamiltonian Quasiparticle excitations about the superconducting ground state are well described within the framework of the Bogoliubov deGennes(BdG) Hamiltonian. For the spin-singlet paired superconductor, which we focus on here, the BdG Hamiltonian has the form $`_0=H_1+H_2,`$ (1) $`H_1`$ (2) $`=`$ $`{\displaystyle d^dxc_\sigma ^{}(x)\left(\frac{\left(i\mathrm{}\stackrel{}{}\frac{e}{c}\stackrel{}{A}(x)\right)^2}{2m}E_F+V(x)\right)c_\sigma (x)},`$ (3) $`H_2`$ (4) $`=`$ $`{\displaystyle d^dxd^dx^{}(c_{}^{}(x)\mathrm{\Delta }(x,x^{})c_{}^{}(x^{})+c_{}(x)\mathrm{\Delta }^{}(x,x^{})c_{}(x^{}))},`$ (5) where $`c^{}`$ and $`c`$ are electron creation and annhilation operators respectively, $`m`$ the mass , $`E_F`$ the Fermi Energy , $`V(x)`$ a random potential describing impurities in the system, and $`A(x)`$ a vector potential describing any external magnetic field. The lattice version of $`_0`$, which is more tractable for numerics, has the form $$_{0L}=\underset{ij}{}[t_{ij}\underset{\alpha }{}c_{i\alpha }^{}c_{j\alpha }+(\mathrm{\Delta }_{ij}c_i^{}c_j^{}+h.c)].$$ (6) From Hermiticity, one requires the condition $`t_{ij}=t_{ji}^{}`$, while spin rotation invariance requires $`\mathrm{\Delta }_{ij}=\mathrm{\Delta }_{ji}`$. The BdG Hamiltonian contains anamolous terms that reflect the fact that the excitations do not conserve charge, and hence break the associated $`U(1)`$ symmetry. As a result, the BdG Hamiltonian lives in an ’extended particle-hole space’ with twice the degrees of freedom of normal electronic systems. More importantly, the absence of $`U(1)`$ charge conservation endows the BdG Hamiltonian with symmetry properties which are completely different from normal systems. It is this difference that plays the key role in giving rise to properties that are unique to superconducting systems. The symmetries of the BdG Hamiltonian have recently been studied in the context of mesoscopic systems and random matrix theory. Within the class of models described in Ref, the Hamiltonians that we study in this paper have SU(2) and T symmetry; they describe systems that are singlet-paired and have no spin-orbit scattering, thus having spin-rotational invariance, and in addition, they have time-reversal invariance. Since the quasiparticles conserve both spin and energy, one can explicitly write Eq.6 in terms of conserved quantities by defining a new set of fermionic $`d`$operators: $$d_i=c_i,d_i=c_i^{}.$$ (7) Thus one makes a particle-hole transformation on the spin-down operators, but leaves the spin-up operators unchanged. The Hamiltonian of Eq.6 now takes the form $$_L=\underset{ij}{}d_i^{}\left(\begin{array}{cc}t_{ij}& \mathrm{\Delta }_{ij}\\ \mathrm{\Delta }_{ij}^{}& t_{ij}^{}\end{array}\right)d_j\underset{ij}{}d_i^{}H_{ij}d_j,$$ (8) where time reversal invariance requires $`t_{ij}`$ and $`\mathrm{\Delta }_{ij}`$ be real. Spin rotational invariance now requires $$\sigma ^yH_{ij}\sigma ^y=H_{ij}^{},$$ (9) where $`\sigma ^y`$ is the standard Pauli matrix. Spin conservation along the z-direction is evident from the fact that the physical spin is related to the number operator for the ’d-particles’ by $$S_i^z=\frac{\mathrm{}}{2}\left(d_i^{}d_i1\right),$$ (10) and that $`_L`$ of Eq.8 conserves particle number. The Hamiltonian $`_L`$, in principle, may be diagonalized by solving the eigenvalue equation $$\underset{i}{}H_{ij}\left[\begin{array}{c}u(j)\\ v(j)\end{array}\right]=E\left[\begin{array}{c}u(i)\\ v(i)\end{array}\right],$$ (11) where $`u`$ and $`v`$ describe the wave-function amplitudes at each site. Given Eq.11, one can construct the following state $$i\sigma ^y\left[\begin{array}{c}u(j)^{}\\ v(j)^{}\end{array}\right]=\left[\begin{array}{c}v^{}(i)\\ u^{}(i)\end{array}\right],$$ (12) with eigenvalue $`E`$; $`SU(2)`$ invariance requires a symmetric dispersion about the Fermi energy, and eigenvalues of the BdG Hamiltonian come in pairs $`(E,E)`$. We now draw attention to the physical situation described by the wavefunctions of Eq.11, and the relation they bear to the phase transition of interest. To begin with, in gapless superconductors, which we focus on here, one finds states at and about the Fermi energy. These states have a profound impact on thermodynamic and transport properties, as discussed previously. Moreover, the transport properties depend crucially on the nature of the eigenstates of Eq.11 at $`E=0`$, which in turn is determined by the spatial configuration of the random potential $`V(x)`$ and the gap-function $`\mathrm{\Delta }(x)`$. In particular, depending on whether these eigenstates are extended or localized, one can conceive of two very different superconducting phases - the ’thermal metal’ which is capable of transporting energy, and the ’thermal insulator’ which cannot conduct energy over large length scales. In fact, it has been shown that both phases are stable in three dimensions. Thus, one can characterize the two phases by the thermal conductivity $`\kappa `$. In systems with SU(2) symmetry, since the same quasiparticles that carry energy also carry spin, the spin conductivity $`\sigma _s`$, may also be used to describe the two phases, and is related to $`\kappa `$ via an analog of the Weidemann-Franz law in the thermal metal: $$\frac{\kappa }{T\sigma _s}=const.$$ (13) Note however, that as the quasiparticles do not carry well defined charge, the Weidemann-Franz law breaks down with regards to electrical conductivity. The two distinct superconducting phases at hand are in analogy with, but quite different from, the metallic and insulating phases in normal systems. We study the critical properties of the transition between the two phases using the formalism and physical set-up described below. ### B A Field-theoretic Formulation As with the case of Anderson localization, the problem of quasiparticle localization within the superconducting state can be described within a field-theoretic framework. To briefly review its key features, in previous work, starting with the Hamiltonian of Eq.6, T.Senthil.et.al derived a useful field theoretic action from which one can extract a rich variety of properties of superconducting systems. Coupling the fermionic degrees of freedom in Eq.6 to an infinitesimal Zeeman field $`\eta `$ (which acts as a chemical potential for the d-quasiparticles), and employing the replica method to calculate disorder averages, they obtained an effective action. Fluctuations about the saddle-point of this action were captured near two spatial dimensions by a non-linear sigma model($`NL\sigma M`$) treatment, yielding for the final form of the action, $$S_{NL\sigma M}=d^dx[\frac{1}{2t}Tr(\stackrel{}{U}.\stackrel{}{U^{}})\eta Tr(U+U^{})].$$ (14) Here, ’t’ is a dimensionless coupling constant that has the physical interpretation of inverse spin conductance, i.e. $`\frac{1}{t}=\frac{\pi }{2}\sigma _s`$. $`U(x)`$ is a matrix with symplectic $`Sp(2n)`$ group structure, where ’n’ is the number of replicas. The action $`S_{NL\sigma M}`$ given above is quite different from the analogous action describing normal systems. The field theory , referred to as the ’the principle chiral $`Sp(2n)`$ model’, has its first term invariant under the global ’rotation’ $`UA^{}UB`$, where $`A^{},BSp(2n)`$ , thus possessesing $`Sp(2n)\times Sp(2n)`$ symmetry. The second term in the action reduces the symmetry to an $`Sp(2n)`$ symmetry as it only allows invariance under $`UA^{}UA`$. A knowledge of this symmetry structure proves to be very useful in deriving critical properties. As done for electrical conductivity in normal systems, we employ a scaling theory for the inverse spin conductivity $`t`$, and analyze the critical point seperating the thermal metal and the thermal insulator. Specifically, we extract the localization length exponent and the unusual singular behaviour of the DOS at the Fermi energy. ### C Hamiltonian for Numerics The predictive power of the effective field theory of the previous section lies in the fact that some results derived from it are universal to all Hamiltonians satisfying the appropriate symmetries. In practice, we find that in obtaining the localization length exponent, an analog of the tight-binding Anderson model commonly used for numerics shows crisp data with much less noise than other models that we have studied. We focus on the same model for the density of states since the localisation length numerics enables us to identify the critical point quite accurately. With reference to Eq.6, the couplings take the form $$t_{ij}=\{\begin{array}{cc}\hfill \frac{1}{\sqrt{2}},& ij,n.n.\hfill \\ \hfill V_{it},& i=j\hfill \\ \hfill 0,& otherwise\hfill \end{array}$$ (15) $$\mathrm{\Delta }_{ij}=\{\begin{array}{cc}\hfill \frac{1}{\sqrt{2}},& ij,n.n.\hfill \\ \hfill V_{i\mathrm{\Delta }},& i=j\hfill \\ \hfill 0,& otherwise\hfill \end{array}$$ (16) where n.n denotes nearest neighbours, and $`V_{it}`$ and $`V_{i\mathrm{\Delta }}`$ are on-site random variables chosen from a uniform probability distribution ranging from $`W`$ to $`+W`$. We work with a three-dimensional cubic lattice described by Eq.6, and having real couplings with the specific form $$_L=\frac{1}{\sqrt{2}}[t\sigma ^z+\mathrm{\Delta }\sigma ^x].$$ (17) Here, the $`\sigma `$’s denote Pauli matrices, and $`t`$ and $`\mathrm{\Delta }`$ are matrices with off-diagonal terms set to unity, and diagonal terms taking on values $`\sqrt{2}V_{it}`$ and $`\sqrt{2}V_{i\mathrm{\Delta }}`$. ### D The Transfer Matrix While we obtain the DOS from the Hamiltonian of Eq.17 by the straightforward process of exact diagonalization, we extract the localization length exponent by means of a transfer matrix formulation that caters to Eq.17 and tremendously reduces the dimensions of the matrices involved in numerical work. To describe its principle, consider a quasi 1-dimensional strip in d-dimensions with cross-sectional area $`L^{d1}`$, and in-going and out-going states at either end of this long strip as shown in Fig.1. One might formally obtain the scattering matrix S for the in-going and out-going states using the definition $$S=\underset{T\mathrm{}}{lim}\mathrm{exp}\frac{iHT}{\mathrm{}}.$$ (18) The scattering matrix, in the specific basis of in-going and out-going states can be written in terms of reflection and transmission matrices ’r’ and ’t’ respectively: $$S=\left(\begin{array}{cc}r& t^{}\\ t& r^{}\end{array}\right).$$ (19) Given this form of the scattering matrix, one can easily derive the transfer matrix T which has the property $$\stackrel{}{\psi _L}=T\stackrel{}{\psi _R}.$$ (20) The symmetries of the BdG Hamiltonian imply that the transfer matrix too has very specific symmetry properties which distinguish it from those of normal systems. Of late, these symmetries have been explored on group theoretic grounds. The transfer matrix T can be constructed by multiplying a set of transfer matrices, each appropriate for a slice of the strip shown in Fig.1: $$\stackrel{}{\psi }_{i+1,R}=T_i\stackrel{}{\psi }_{i,R},$$ (21) $$T=\underset{i=1}{\overset{N}{}}T_i.$$ (22) The form of the transfer matrix that we use for numerical calculation does not make the symmetry of the BdG Hamiltonian manifest, but it is tailored specifically for a tight-binding Hamiltonian such as the one described in Eq.6. We begin by writing the Schrodinger Eq.11 as a difference equation $$A_i\stackrel{}{D}_i+B_{i,i+1}\stackrel{}{D}_{i+1}+B_{i,i1}\stackrel{}{D}_{i1}=E\stackrel{}{D}_i,$$ (23) where $`\stackrel{}{D}_i`$ denotes the wave-function amplitudes on each slice in the spinful d-quasiparticle eigenbasis, and $`A_i`$ and $`B_{i,i+1}`$ are $`2L^2\times 2L^2`$ matrices of the form $$A_i=\left[\begin{array}{cc}t_{ii}& \mathrm{\Delta }_{ii}\\ \mathrm{\Delta }_{ii}& t_{ii}\end{array}\right],$$ (24) $$B_{i,i+1}=\left[\begin{array}{cc}t_{i,i+1}& \mathrm{\Delta }_{i,i+1}\\ \mathrm{\Delta }_{i,i+1}& t_{i,i+1}\end{array}\right],$$ (25) where the $`t_{ij}`$ and $`\mathrm{\Delta }_{ij}`$’s are now matrices coupling slices ’i’ and ’j’ in the manner described by Eq.23. To obtain the transfer matrix, we rewrite Eq.23 as $$\stackrel{}{D}_{i+1}=B_{i,i+1}^1(EA_i)\stackrel{}{D}_iB_{i,i+1}^1B_{i1,i}\stackrel{}{D}_{i1}.$$ (26) For our model, the inter-slice coupling has the simple form $$B_{i,i+1}=\frac{1}{\sqrt{2}}[I\sigma ^x+I\sigma ^z],$$ (27) which satisfies the special property $$B_{i,i+1}=B_{i,i+1}^1.$$ (28) Using Eq.26, we are now in a position to define a transfer matrix as follows: $$\left(\begin{array}{c}\stackrel{}{D}_{i+1}\\ \stackrel{}{D}_i\end{array}\right)=T_i\left(\begin{array}{c}\stackrel{}{D}_i\\ \stackrel{}{D}_{i1}\end{array}\right),$$ (29) where $`T_i`$ has the relatively simple form $$T_i=\left(\begin{array}{cc}B_{i,i+1}(EA_i)& I\\ I& 0\end{array}\right),$$ (30) and the multiplicative property of Eq.22. With this transfer matrix at hand, which very closely resembles standard ones used for the Anderson model, we can extract the localization length for different values of energy E and disorder W in a manner completely analogous to the numerical treatment of the Anderson model. As we are interested in the behaviour of states at the Fermi energy, we set the energy E, to zero. The procedure for extracting the localization length is quite standard, and has been elaborated on in great depth in many works. But to briefly outline the method, one begins with an orthonormal basis of vectors $`\widehat{o_i}(0)`$ in the space of the transfer matrix, representing the right-most states in Fig.1. One then assumes that for a given disorder strength W, and width L, there exists a set of eigenmodes $`\widehat{w_i}(L,W)`$ that describes typical eigenmodes for a quasi 1-dimensional system of N slices, where ’N’ is large enough to represent the average behaviour of the random disorder, and that these modes decay or grow in magnitude as $`\mathrm{exp}(\pm \gamma _iN)`$ upon multiplication with the corresponding transfer matrix $`T=_{i=1}^NT_i`$. One further assumes that when the set of basis vectors $`\widehat{o_i}(0)`$ , each of which may be represented as a linear combination of the $`\widehat{w_i}(L,W)`$,is multiplied by T, the resulting vectors $`\stackrel{}{v_i}(1)`$ (where the dependence on L and W is implicit), are each composed of the appropriate linear combination of modes $`\widehat{w_i}(L,W)`$ now weighted by the corresponding growth factors $`\mathrm{exp}(\pm \gamma _iN)`$. The length ’N’ is numerically restricted by the exponentially growing magnitude of the vectors $`\stackrel{}{v}`$. Since all vectors have now grown fastest along the fastest growing mode, say $`\widehat{w_1}(L,W)`$, the magnitude and direction of any one of these modes, say $`\stackrel{}{v_1}(1)`$, are $`\mathrm{exp}\gamma _1N`$ and along $`\widehat{w_1}(L,W)`$ respectively. Projecting out the component $`\stackrel{}{v_1}(1)`$ from the next vector $`\stackrel{}{v_2}(1)`$ gives a resultant vector $`\stackrel{}{o_2}(1)`$ whose magnitude is roughly $`\mathrm{exp}\gamma _2N`$. Thus, by such a ’Gram-Schmidt’ orthogonalization procedure for the whole set of vectors $`\stackrel{}{v_i}(1)`$, one obtains a set of orthogonal vectors $`\stackrel{}{o_i}(1)`$ with associated Lyapunov exponents $$\gamma _i(1)=\frac{\mathrm{ln}\|\stackrel{}{o_i}(1)\|}{N},$$ (31) which give the characteristic inverse localization lengths associated with each mode. To reduce computational effort, we consider just the positive Lyapunov exponents corresponding to exponentially growing states. As we are restricted in our length size N, to obtain a fair estimate of the typical $`\gamma _i`$’s, we repeat the procedure of transfer matrix multiplication, now using as our initial basis vectors the normalized set $`\widehat{o_i}(1)`$ which more or less point along the ideal basis $`\widehat{w_i}(L,W)`$. An average value $`\frac{1}{M}_{j=1}^M\gamma _i(j)`$ obtained from M such iterations provides the desired estimate of the ideal $`\gamma _i`$’s. We associate the characteristic localization length $`\lambda (L,W)`$ with the slowest decaying mode, and thus with the inverse of the smallest positive Lyapunov exponent $`\gamma _{min}(L,W)`$. In the quasi 1-dimensional case, all modes are exponentially decaying or growing (corresponding to in-going or out-going states respectively) since even the slightest disorder is enough to localize states. But in the 3-dimensional limit, where the 2-dimensional cross-sectional area becomes large, we know that the modes ought to experience a transition from extended to localized behaviour as a function of disorder. To determine the critical disorder strength $`W_c`$ for this transition, the nature of the modes in 3-dimensions and the localization length for an infinite size system of given disorder strength, one can use a finite size scaling analysis of the quasi 1-dimensional system. The scaling function that we will use to do so is the dimensionless parameter $$\mathrm{\Lambda }(L,W)=\frac{\lambda (L,W)}{L}.$$ (32) We now turn to the critical properties of the phase transition between the thermal insulator and thermal metal. ## III Critical Properties As in disordered electronic systems, we have seen in Sec.II B that the thermal metal, the thermal insulator ,and the critical point seprating them may be characterized by their transport properties. The replica field theory of Eq.14, with its dimensionless coupling t, provides a powerful means of studying this transition. An analysis of the action in Eq.14 shows that in $`2+ϵ`$ dimensions, where $`ϵ=1`$ for our system, an unstable fixed point $`t_c`$ describes the critical point between the thermal metal and thermal insulator. One can study the scaling behaviour of t with system size L(see Fig.2) by deriving a form for the scaling function $`\beta (t)=\frac{dt}{d\mathrm{ln}L}`$, which we present explicitly in the next section. Near $`t_c`$, for $`t>t_c`$, the coupling $`t`$ grows larger with L, and thus exhibits a stable thermal insulator, while for $`t<t_c`$, a smaller and smaller value of $`t`$ with increasing length scale signals a thermal metal. ### A Localization Length Exponent The $`\beta `$ function for the $`Sp(2n)`$ sigma model of the action of Eq.14 can be found in Ref., and it is given to cubic order(’two’ loop) in coupling t by $$\beta (t;Sp(2n))=ϵt+(2n+1)+\frac{1}{2}(2n+1)^2t^3+𝒪(t^4),$$ (33) where $`ϵ=d2`$, and n denotes the number of replicas. In the limit $`n0`$, at the critical point where the $`\beta `$ function vanishes, its derivative gives the inverse localization length exponent: $$\frac{1}{\nu }=ϵ+\frac{ϵ^2}{2}+𝒪(ϵ^3).$$ (34) In contrast, normal systems with time reversal symmetry and spin-rotational invariance may be identified with the $`Sp(4n)/Sp(2n)\times Sp(2n)`$ model of Ref., and its associated $`\beta _n`$ function has the form $$\beta _n=ϵt+(4n+1)t^2+(8n^2+2n)t^3+𝒪(t^4).$$ (35) The localization length exponent derived from Eq.35 has the value $$\frac{1}{\nu _n}=ϵ+𝒪(ϵ^3),$$ (36) quite different from the value of $`\nu `$ in Eq.34. The numerical evidence to follow supports the field theoretic result that the value of the localization length exponent $`\nu `$ for the superconducting system with T and $`SU(2)`$ is considerably lower than the analogous exponent $`\nu _n`$ for normal systems. #### 1 Numerical Treatment The standard numerical technique that we employ for extracting the localization length exponent $`\nu `$, shows that in three dimensions it takes on the value $`1.15\pm 0.15`$. The finite-size scaling technique can be summarized as follows: scaling arguments require that the only relevant length scale in the system be the localization length $`\xi (W)`$ of the infinite sized system, and thus we have $$\frac{\lambda (L,W)}{L}=\mathrm{\Lambda }(L,W)=h(\frac{\xi (W)}{L}),$$ (37) where h is a scaling function yet to be determined. Close to the critical point $`W_c`$, we have the localization length $`\xi `$ of the infinite system behaving as $$\xi \|WW_c\|^\nu ,$$ (38) which means that the argument ’$`x`$’ of $`h(x)`$ in Eq.37 blows up at the critical point. However, $`\mathrm{\Lambda }`$ is well-behaved and finite; this is only possible if we have $$\underset{x\mathrm{}}{lim}h(x)=const.,$$ (39) where the constant refers to independence with respect to L for large L. Thus, the critical value $`\mathrm{\Lambda }_c`$ is common to all sufficiently large system sizes. Using Eq.38, we rewrite Eq.37 as $$\mathrm{ln}\mathrm{\Lambda }(L,W)=f[L^{\frac{1}{\nu }}(WW_c)].$$ (40) Linearizing the function ’f’ about the critical fixed point $`(W_c,\mathrm{ln}\mathrm{\Lambda }_c)`$ yields $$\mathrm{ln}\mathrm{\Lambda }(L,W)=\mathrm{ln}\mathrm{\Lambda }_c+A(WW_c)L^{1/\nu }.$$ (41) To procure the value of $`\nu `$, we use an iterative procedure which is equivalent to the widely used procedure of Ref. of performing a least square fit to obtain actual values of $`\xi (W)`$. We rewrite Eq.41 as $`\mathrm{ln}\mathrm{\Lambda }`$ $`=`$ $`AL^{1/\nu }W(\mathrm{ln}\mathrm{\Lambda }_c+AL^{1/\nu }W_c)`$ $`=`$ $`a(L)Wb(L,\mathrm{\Lambda }_c,W_c),`$ assume an initial value for the critical point $`(W_c,\mathrm{ln}\mathrm{\Lambda }_c)`$ from Fig.4, obtain the functions $`a(L)`$ and $`b(L,\mathrm{\Lambda }_c)`$ by curve-fitting, extract $`A`$ and $`\nu `$, determine the value of the critical point thus obtained, and repeat the procedure till convergence is achieved (which happens rather quickly).If scaling is valid, we can collapse our data onto curves of $`\mathrm{ln}\mathrm{\Lambda }(W,L)`$ vs $`L^{1/\nu }(WW_c)`$. In our transfer matrix calculations, we choose systems whose cross-sectional areas have linear dimensions of $`L=4,6,8,10`$, and our transfer matrices have dimensions $`4L^2\times 4L^2`$ with the given values of $`L`$. We choose the number of transfer matrices to be multiplied together by ensuring that the basis vectors do not grow upto a magnitude greater than $`10^5`$ upon being multiplied by the set of transfer matrices. We utilize a total of 2000 slices in the quasi 1-dimensional system for each value of L and W. Fig.3 shows the plots for $`\mathrm{\Lambda }(L,W)`$ as a function of disorder for $`L=4,6,8,10`$. For fixed disorder W, an increasing $`\mathrm{\Lambda }(L,W)`$ with increasing system size $`L`$ indicates the extended regime, while a decreasing $`\mathrm{\Lambda }(L,W)`$ shows that the system is in the localized regime. In comparison to the Anderson model for normal systems, our simulations require a much smaller number of transfer matrices for relatively noise-free data, and we believe that this really is a consequence of the relatively low critical disorder strength. Fig.4 shows the data for the iterative procedure which gives the value for the localization length exponent $$\nu =1.15\pm 0.15.$$ (42) Finally, Fig.5 indeed demonstrates clean data collapse close to the critical point. #### 2 Summary of Results Both field theory and numerics concur with the fact that the localization length exponent in the superconducting systems is significantly lower than that of their normal partners, clearly indicating a new universality class. An $`ϵ`$ expansion in $`d=2+ϵ`$ dimensions of the action in Eq.14 shows that the superconducting system has a localization length exponent $`\nu `$ of $`(ϵ+\frac{ϵ^2}{2}+𝒪(ϵ^3))^1`$ in contrast to a $`\nu _n`$ of $`(ϵ+𝒪(ϵ^3))^1`$ for normal systems. The field theory would thus predict $`\nu =\frac{2}{3}`$ versus $`\nu _n=1`$ in 3-dimensions. In comparison, one obtains the numerical estimate $`\nu =1.15\pm 0.15`$ versus $`\nu _n=1.54\pm 0.08`$ in 3-dimensions. We must remark that the system sizes and computing power utilized in our numerical studies were relatively low compared to the current cutting edge procedures. As a lot of work has gone into refining techniques with regard to normal systems(for e.g. Ref), it is well worth employing them to study this novel phase transition and analogous ones in superconducting systems with other symmetries. ### B Density of States The quasiparticle DOS in dirty superconducting systems exhibits some of the most stunning differences from normal systems. In normal systems, both in the Anderson metal and the Anderson insulator, i.e., in the absence of interactions, the DOS remains a smooth continuous function across the Fermi energy. In contrast, in gapless superconductors that respect $`SU(2)`$, well within the thermal metal, quantum interference effects cause a singularity at the Fermi energy which manifests itself as a $`\sqrt{E}`$ cusp in 3-dimensional systems. Deep in the thermal insulator, the density of states exhibits a power-law that vanishes at the Fermi energy with the form $`\rho \|E\|^\alpha `$, where $`\alpha =1`$ for systems possessing time-reversal invariance. About the critical point, the DOS once again shows power-law singularities, the details of which we discuss below. The curious form of the DOS (shown in Fig.6) has profound impact on thermodynamic properties, and in particular, manifests itself in quantities such as specific heat and spin susceptibility. #### 1 Discussion of Critical Behaviour The field theoretic action of Eq.14 not only offers concrete predictions for the DOS, if one were to use Wegner’s analogy with magnetic systems, it provides an intuitive picture for the behaviour about criticality. To elaborate, the quasiparticle DOS at the Fermi energy, which also gives a measure of the magnetization, acts as the ’order parameter’ of the field theory. It is given by $$\rho =\underset{n0}{lim}\frac{\rho _0}{4n}<Tr(U^{}+U)>,$$ (43) where $`\rho _0`$ is the bare DOS, and n the number of replicas. The field $`\eta `$, which has units of energy E, couples to the DOS in the action of Eq.14, and might be equated with the magnetic field in the magnetic analog. With a little indulgence, one can go further with parallels between the field theory and the magnetic systems, as first suggested for normal systems: \| F.T. of $`S_{NL\sigma M}`$ \| Magnetic Systems \| \| --- \| --- \| \| Distance from criticality \| Reduced Temperature \| \| $`\mathrm{\Delta }=\frac{WW_C}{W_C}`$ \| t \| \| DOS, $`\rho `$ \| Magnetization, m \| \| Energy, E \| Magnetic Field, h \| \| $`\stackrel{~}{\chi }=\frac{d\rho }{dE}`$ \| Magnetic Susceptibility,$`\chi `$ \| In normal systems, the analogy is clouded by the fact that the DOS is a continuous function of energy and disorder respectively. One can reconcile with this if the DOS obeys a power-law form with exponent zero, and in fact, one can show this to be the case on field theoretic grounds. But in superconducting systems, as we shall see, the analogy goes through in quite a striking manner, with a whole slew of nontrivial critical exponents: \| $`\rho (\mathrm{\Delta },E=0)\|\mathrm{\Delta }\|^\beta `$ \| \| --- \| \| $`\rho (\mathrm{\Delta }=0,E)\|E\|^{\frac{1}{\delta }}`$ \| \| $`\stackrel{~}{\chi }\|\mathrm{\Delta }\|^\gamma `$ \| \| $`\xi \|\mathrm{\Delta }\|^\nu `$ \| where we have used the notation of the tables above,and $`\xi `$ is the localization length and $`\nu `$ the associated exponent described in the previous section. In order to derive expressions for $`\beta `$, $`\delta `$ and $`\nu `$, we start with the ’free energy density’ f, obtained from the action of Eq.14 : $$f=\frac{1}{L^d}\underset{n0}{lim}\frac{ln𝒵_n}{n},$$ (44) $$𝒵_n=\overline{Z^n}=𝑑\stackrel{~}{U}e^{S_{NL\sigma M}},$$ (45) where $`d\stackrel{~}{U}`$,the integral volume element, takes into account the symplectic group structure of the matrices in the action $`S_{NL\sigma M}`$, n denotes the number of replicas, $`Z`$ denotes the partition function of single system, and the overbar above $`Z^n`$ refers to the average over disorder. Near criticality, $`f_s`$, the singular part of the free energy density, is expected to scale as follows: $$f_s(\mathrm{\Delta },E)=\xi ^d\stackrel{~}{f}(\xi ^yE),$$ (46) where ’y’ describes the scaling form of E, and Eq.38 gives the behaviour of the correlation length $`\xi \|\mathrm{\Delta }\|^\nu `$. Differentiating the free energy with respect to E results in the following form for the DOS: $$\rho (\mathrm{\Delta },E)=\xi ^{d+y}F^\pm (\xi ^yE),$$ (47) where $`F^+`$ corresponds to behaviour for $`\mathrm{\Delta }>0`$, and $`F^{}`$ for $`\mathrm{\Delta }<0`$. To obtain $`\beta `$, we set $`E=0`$, and compare the form of the resulting order parameter $`\rho (\mathrm{\Delta },E=0)`$ in the above table, yielding $$\beta =\nu (dy).$$ (48) Here, we require $`F^{}(0)`$ to be finite, and $`F^+(0)=0`$. To extract $`\delta `$, we impose the physical constraint that $`\rho `$ be well-behaved and finite at criticality. This requires that $`\rho (\mathrm{\Delta }0,E)`$ be independant of the diverging correlation length, and thus yields $$\frac{1}{\delta }=\frac{d}{y}1.$$ (49) Finally, taking a derivative of $`\rho `$ in Eq.47 with respect to E gives us the following expression for $`\gamma `$: $$\gamma =\nu (2yd).$$ (50) #### 2 Results To obtain estimates of critical exponents from field theoretic results in $`2+ϵ`$ dimensions, we use the value $`(ϵ+\frac{ϵ^2}{2})^1+𝒪(ϵ^3)`$ obtained for $`\nu `$ in the previous section, and the value $`\frac{1}{\delta }=\frac{ϵ}{4}+𝒪(ϵ^3)`$ from Ref.. Eq.48-50 then enable us to determine the critical exponents $`\beta `$ and $`\gamma `$ via the relationship $`\frac{4d}{ϵ+4}`$. Specifically, in the case of 3-dimensions, substituting the value $`ϵ=1`$, we obtain the rough estimates $`y=12/5`$, $`\delta =4`$, $`\nu =2/3`$, $`\beta =2/5`$ and $`\gamma =6/5`$. Shifting our focus to numerical results, the method of exact diagonalization reveals that the superconducting system at hand does indeed show the novel singular behaviour in the DOS at the Fermi energy. In the data shown below, we have once more modelled the superconducting Hamiltonian after Eq.8 using periodic boundary conditions. Systems of linear dimension ’L’ have required matrices of dimension $`2L^3X2L^3`$, and we have explored system sizes with linear dimensions $`L=4,6,8`$. Fig.7 shows the progression of the behaviour of the DOS with increasing disorder strength. As seen in the last panel of Fig.7, the DOS shows a power-law behaviour of the form $`\rho \|E\|`$ about the Fermi energy , $`E_F`$, consistent with expectations for the thermal insulator. Fig.8 shows a zoom of the DOS about $`E_F`$ for disorder close to the critical strength $`W_C`$, for which we have an estimate from the localization length study of the previous section. One can easily discern that the DOS plummets down quite markedly, and does indeed exhibit singular power-law behaviour. A plot of the DOS at the Fermi energy(Fig.9) shows that even relatively small system sizes provide numerical confirmation of the fact that $`\rho (E=0)`$ acts as the order parameter for the field theory of Eq.14; the DOS at $`E_F`$ is finite for low disorder, and it slowly drops to zero beyond a critical disorder strength. As discussed in the previous section, one would in fact expect the DOS for an infinite sized system to behave as $`\rho (\mathrm{\Delta },E=0)\|\mathrm{\Delta }\|^\beta `$, where $`\mathrm{\Delta }`$ is the distance from criticality within the thermal metal. Scaling arguments for extracting $`\beta `$ require $$\rho _L(\mathrm{\Delta },E=0)=\|\mathrm{\Delta }\|^\beta Y(L\mathrm{\Delta }^\nu ),\mathrm{\Delta }<0,$$ (51) where $`\rho _L`$ is the DOS associated with a system of linear dimension L, and $`Y`$ is a scaling function. One can rewrite the above equation in a form more conducive to numerics as follows: $$\rho _L(\mathrm{\Delta },E=0)=L^{\frac{\beta }{\nu }}\stackrel{~}{Y}(\mathrm{\Delta }L^{\frac{1}{\nu }}),$$ (52) where $`\stackrel{~}{Y}`$ is yet another scaling function with limiting behaviour $`\stackrel{~}{Y}(x\mathrm{})=\|x\|^\beta `$, reproducing the required dependence of $`\rho (\mathrm{\Delta },E=0)`$ on $`\mathrm{\Delta }`$ for infinite system size. Fig.10 exhibits the plots of $`\rho _L`$ as a function of disorder for different system sizes L, and we make use of this data to procure the value of $`\beta `$ in Eq.52. To extract $`\beta `$, we perform a fit taking $`W_c`$, $`\nu `$ and $`\beta `$ as variable parameters. Exploiting the universal nature of the function $`\stackrel{~}{Y}`$ in Eq.52, we find the appropriate values of $`\rho `$ obtained by linear interpolation for a given set of system sizes and fixed argument in $`\stackrel{~}{Y}`$, and plot these on a log-log scale versus system size; the slope for a linear fit of such a set of points determines $`\beta `$. The actual value of $`\beta `$ is obtained by performing the above procedure for different values of the argument of $`\stackrel{~}{Y}`$ and taking the average of the $`\beta `$’s thus obtained. The set of values $`\beta =0.15`$, $`W_c=4.67`$ and $`\nu =1.25`$ result in the data collapse shown in Fig.11. In comparison, as remarked at the beginning of this section, the field theoretic result predicts that $`\beta =0.4`$. Once more, as in the case of the localization length exponent, we comment on numerical accuracy; other simulations using exact diagonalization, for instance, those catering to specific physical situations, have used larger sytems sizes and number of realizations which would be well-worth employed here. However, the above numerics conveys quite clearly that the DOS at criticality exhibits a power law suppression about the Fermi energy, and that $`\rho (E=0,W)`$ acts as an ’order parameter’ with a non-trivial exponent $`\beta `$ in surprising contrast to normal systems which have $`\rho (E,W_c)`$ smooth about $`E_F`$, and a vanishing exponent $`\beta `$. ## IV Other Systems We have studied the thermal metal-thermal insulator transition in superconducting systems with $`SU(2)`$ and $`T`$, and discussed the dirty d-wave superconductor as a possible physical realization. Superconducting systems with other symmetries too promise such a transition. The properties and phases of superconducting systems with spin-rotational invariance, but no time reversal invariance are rather similar to our case which preserves both symmetries. The thermal metal and the thermal insulator are both stable only in 3-dimensions, and the transition cannot occur in lower dimensions where quasi-particle excitations are always found to be localized at large enough length scales. Due to the absence of time-reversal symmetry, the Hamiltonian of Eq.6 now has imaginary couplings. As described in detail in previous work, in this symmetry category, the pinned vortex state of a type II s-wave superconductor appears to be a fine candidate for exhibiting the thermal metal-thermal insulator transition. Low energy quasiparticles bound to the core of the vortices can tunnel from one vortex to another, and as the magnetic field is increased, the density and tunneling strength also increase. It is conjectured that there could exist a critical magnetic field $`H_{c4}`$, within the vortex phase at which the low energy states can permeate through the medium to form extended states. Field theoretic methods have analyzed the properties of these systems, and in particular, have shown that in parallel to the case with T, the density of states at the critical point has the power-law behaviour $`\rho \|E\|^{ϵ/2}`$, where $`E`$ is the energy and $`ϵ=d2`$. The presence of spin-orbit scattering or of triplet pairing introduces new ingredients. It breaks spin-rotational symmetry, and as in the case of normal systems, field theoretic arguments predict the presence of a delocalization-localization transition not only in 3-dimensions but also in 2-dimensions.In Ref., T.Senthil et.al. cast their Hamiltonian in terms of Majorana fermions, and their formulation, among other things, is highly conducive to numerical work. They discuss possibilities for experiment, and mention heavy fermion systems where spin-orbit scattering is prominent, as a physical realization. Finally, the tantalizing prospect of a Hall effect in superconducting systems has been explored in systems with other symmetries as well,i.e., those with SU(2). Ref shows that a superconductor with unconventional $`d_{x^2y^2}+id_{xy}`$ pairing symmetry is capable of exhibiting a phase with non-vanishing spin and thermal Hall conductances. Indeed, as in normal systems, sophisticated methods such as employing the network model and supersymmetric spin chains, have shed light on these systems. ## V Experiment The transport properties of normal systems have been probed in great depth, and now we see that superconducting systems could potentially offer an equally rich range of experiments in the thermal metal-thermal insulator transition. In the previous section, we have mentioned a variety of experimental candidates for study, such as the dirty d-wave, the type II s-wave, heavy fermion and other superconducting systems. These systems must share the feature of gapless superconductivity; one requires states at the Fermi energy since it is these states that determine transport properties, and distinguish the thermal insulator from the thermal metal. Associated with each system, a tunable parameter such as disorder or magnetic field ought to be able to access the phases. As was previously discussed in detail, the type II s-wave superconductor in the vortex phase offers promise as a likely candidate for observing the novel transition since in principle, one need only tune the magnetic field, and generally the vortex phase exists over a large range of field. As seen earlier, in contrast to normal systems, the density of states shows singular behaviour about the Fermi energy for both phases and at the critical point. We saw that for systems with T and SU(2), it obeys a power-law behaviour of the form $`\rho (E)\|E\|^\alpha `$, where $`\alpha =\frac{1}{2}`$ well within the thermal metal, $`\alpha =1`$ deep in the thermal insulator, and field theory predicts $`\alpha =\frac{ϵ}{4}`$ at the transition, with $`ϵ=1`$ for 3-dimensional systems. This singular behaviour ought to be reflected in thermodynamic quantities such as specific heat and spin susceptibility, and in tunneling experiments. In particular, the temperature dependence of the thermodynamic quantities would have a form $`CT^{1+\alpha }`$ for the specific heat, and $`\chi T^\alpha `$ for the spin susceptibility. Recent experiments of cuprate superconductors have observed a suppression of the specific heat close to the Fermi energy. However, measuring the differing behaviours to detemine the phase might prove tricky, especially since we have neglected various effects such as interactions which could come into play. For systems that preserve spin-rotational invariance, quasiparticle excitations about the superconducting ground state not only conserve energy, but also spin; the spin-conductance $`\sigma _s`$ could be employed to determine whether the system inhabits the thermal metal or thermal insulator. In the thermal metallic phase, a magnetic field gradient would cause the spinful quasiparticle to diffuse across the sample, while in the thermal insulator, they would be unable to conduct spin. A variety of refined spin-injection techniques have been developed in semi-conductors to measure spin dynamics (see for e.g. Ref.), but by no means would it be a simple task to cater these experiments to superconductors. We believe that by far, thermal conductivity measurements would offer most promise in probing the thermal metal-thermal insulator transition. For all superconducting systems with their differing symmetries, the thermal conductivity $`\kappa `$, distinguishes the two phases in that the ratio $`\kappa /T`$ tends to a finite constant in the thermal metal, and to zero as $`T0`$ in the thermal insulator. Along the lines of previous experiments, it would be extremely interesting to observe the transition by applying a small thermal gradient across a superconducting sample and measuring $`\kappa /T`$ as a function of a tunable parameter. We thank T.Senthil and I.A. Gruzberg for many an illuminating conversation, and D. Whysong for indispensible advice on numerical work. This research was supported by NSF Grants DMR-97-04005, DMR95-28578 and PHY94-07194.
warning/0003/nucl-th0003070.html	ar5iv	text	# References Neutron matter with a model interaction M. Ya. Amusia<sup>a,b</sup>, V. R. Shaginyan<sup>c,d</sup> <sup>1</sup><sup>1</sup>1E–mail: vrshag@thd.pnpi.spb.ru <sup>a</sup>The Racah Institute of Physics, Hebrew University, Jerusalem 91904, Israel; <sup>b</sup>A.F. Ioffe Physical-Technical Institute, 194021 St.Petersburg, Russia; <sup>c</sup>Petersburg Institute of Nuclear Physics, 188350 Gatchina, Russia; <sup>d</sup> Department of Physics, University of Washington, Seattle, WA 98195, USA ## Abstract An infinite system of neutrons interacting by a model pair potential is considered. We investigate a case when this potential is sufficiently strong attractive, so that its scattering length $`a`$ tends to infinity, $`a\mathrm{}`$. It appeared, that if the structure of the potential is simple enough, including no finite parameters, reliable evidences can be presented that such a system is completely unstable at any finite density. The incompressibility as a function of the density is negative, reaching zero value when the density tends to zero. If the potential contains a sufficiently strong repulsive core then the system possesses an equilibrium density. The main features of a theory describing such systems are considered. PACS. 24.10.Cn Many-body theory - 21.10.Dr Binding energies There exists a well known problem in many body physics related to the description of the ground state properties of an infinite system composed of interacting fermions. In general, this description is based usually on tedious numerical calculations, particularly when the interaction is rather strong. The well known exceptions from this situation, when it is possible to calculate the ground state properties analytically, are the Random Phase Approximation (RPA) for a high density electron gas and the low density approximation for dilute gases . In both cases the kinetic energy $`T_k`$ is much bigger then the interaction energy $`E_{int}`$ of the system. This allows to apply some kind of a perturbation theory. In the case of homogeneous electron liquid it turns out that the analytical RPA-like description is also possible not only at very high but medium densities when $`T_kE_{int}`$ . Similar extension of the range of validity is impossible in the case of fermion systems at low densities $`\rho `$: there the gas approximation is not applicable if $`T_kE_{int}`$. In the cases when the pair interaction is attractive and sufficiently strong, the system can have a quasi equilibrium or equilibrium states in which $`T_kE_{int}`$. On the other hand, these states are preceded by special points with density $`\rho `$ values at which the incompressibility $`K(\rho )`$ of the system tends to zero. Thus, if it would be possible to predict the existence of such points then in principle it would become possible to conclude that the system has at least a quasi equilibrium state. In this Short note we address the above mentioned problem and consider the ground state properties of the infinitely extended multi-fermion system. We demonstrate that it can be done analytically provided that the pair interaction between fermions is characterized only by the scattering length $`a\mathrm{}`$. One can say in this case that the scattering length is the dominant parameter of the problem under consideration. Such an investigation is of great importance since it can be applied to fermion systems interacting via potentials with not only infinite, but also sufficiently large $`a`$. For instance, the scattering length $`a`$ of neutron-neutron interaction is about $`20`$ fm, thus greatly surpassing the radius of the interaction $`r_0`$. On the other hand, it is possible now to prepare artificially a system composed of Fermi atoms interacting by an artificially constructed potential with almost any desirable scattering length, similarly to that how it is done for the trapped Bose gases, see e.g. . An experimental study, performed on such Fermi-system would be of great importance presenting new information on the behavior of dilute gases and the gas-liquid phase transition. Let us consider the interaction of two isolated particles. We assume that this interaction is of finite radius $`r_0`$, which is small, so that $`p_Fr_01`$ ($`p_F=(3\pi ^2\rho )^{1/3}`$ is the Fermi momentum), but its strength is such that the scattering length is negative and infinitely big, $`a\mathrm{}`$. We assume also, that the density $`\rho `$ of the system under consideration is homogeneous. As it will be demonstrated below, in such a case the system is located in the vicinity of a phase transition, transforming it into a strongly correlated one. Therefore, the problem of calculating its ground state properties has to be treated for the most part qualitatively. Let us start considering general properties of a Fermi system with some attractive two-particle bare interaction $`V(r)`$, which is sufficiently weak to create a two-particle bound state. Assume, that the scattering length $`a`$ corresponding to this potential is negative and finite. The ground state energy density $`E(\rho )`$ can be approximated by Skyrme-like expression , $$E(\rho )=\frac{3p_F^2}{10M}\rho +t_0\rho ^2+t_3\rho ^{7/3},$$ (1) Here $`M`$ is the particle mass, and $`\rho `$ is their density in the system. The first term of eq. (1) is the kinetic energy $`T_k`$, while the second and the third are related to the interaction energy $`E_{int}`$ determined by the potential $`V(r)`$. The second term which is proportional to $`t_0`$ gives a proper description in the gas limit. The third term provides the behavior of $`E(\rho )`$ at higher densities, including that of the equilibrium density. Such structure of $`E(\rho )`$ appears if the interaction is sufficiently attractive so that $`t_0<0`$, and $`t_3>0`$. Note, that eq. (1) presents at least a qualitative description of the system under consideration giving a rather simple and reasonable presentation of the function $`E(\rho )`$. A more precise picture of the energy dependence upon density can be obtained using more sophisticated functionals for the ground state energy . If the potential $`V`$ is of short range and purely attractive, then in the Hartree-Fock approximation the ground state energy $`E_{HF}`$ is given by the following expression $$E_{HF}(\rho )=\frac{3p_F^2}{10M}\rho +t_{HF}\rho ^2,$$ (2) where the parameter $`t_0=t_{HF}`$ , being negative, is entirely determined by the potential $`V(r)`$. For instance, in the case of a short range $`\delta `$-type interaction one has $`t_{HF}=v_0/4`$, with $`v_0`$ being the corresponding strength of the potential. Eq. (2) shows that at small densities $`E_{HF}>0`$ due to the kinetic energy term, but at sufficiently high densities $`\rho \mathrm{}`$ the Hartree-Fock energy becomes dominating, leading to the collapse of the system, with $`E_{HF}\mathrm{}`$. Keeping in mind that the Hartree-Fock approximation gives the upper limit to the binding energy $`E_{HF}E`$, one can conclude that the system does not have, in this case, an equilibrium density $`\rho _e`$ and energy $`E_e`$ since $`E_e\mathrm{}`$ when $`\rho \mathrm{}`$ . Note, that for a given and finite total number of particles $`N`$, the HF energy is not going to infinity and the system collapses into a small volume with the radius $`r_0`$, with the density $`\rho N/r_0^3`$. It is evident that the function $`E(\rho )`$ is positive at small densities, if the parameter $`t_0`$ is finite. Therefore, it must have at least one maximum at the density $`\rho _m`$ before it becomes negative, on the way to $`E\mathrm{}`$. If the potential $`V(r)`$ includes a repulsive core at sufficiently short distances, then $`t_3>0`$ . As a result, the system has an equilibrium density and energy, $`\rho _e`$ and $`E_e`$, respectively, determined by the repulsive core strength and its radius $`r_cr_0`$. Now let us apply eq. (1) to demonstrate the most important features of the system under consideration: a) when $`\rho 0`$ the third term on r.h.s. in eq. (1) can be omitted. The kinetic energy is relatively very big, $`T_kE_{int}`$, and $`t_0a`$, with $`a<0`$ being the scattering length. In that case we have a dilute Fermi gas with positive pressure $`P`$ and incompressibility $`K`$, the latter being determined by the equation, see e.g. , $$K(\rho )=\rho ^2\frac{dE^2(\rho )}{d\rho ^2}.$$ (3) b) on the way to higher densities, which can be achieved by applying an external pressure, the system reaches the density $`\rho _{c1}<\rho _m`$ at which the incompressibility is equal to zero, $`K(\rho _{c1})=0`$. Remembering that at the maximum the second derivative is negative, one can conclude, as it is seen from eq. (3), that $`K(\rho _m)<0`$. In the range $`\rho _{c2}\rho \rho _{c1}`$ the incompressibility is negative, $`K<0`$, and as a result the system becomes totally unstable. In fact, in this density range all calculations of the ground state energy are meaningless since such a system cannot exist and thus be observed experimentally ; c) at some point $`\rho =\rho _{c2}>\rho _m`$ the contribution due to the repulsive core becomes sufficiently strong to prevent the further collapse of the system. The incompressibility attains $`K=0`$ at $`\rho _{c2}<\rho _e`$, being positive at the higher densities. Finally, the system becomes stable at $`\rho >\rho _{c2}`$, reaching equilibrium density at $`\rho _e`$ with equilibrium energy equal to $`E_e`$. It is obvious that $`K(\rho _e)>0`$ being proportional to the second derivative at the minimum, see eq. (3). It should be kept in mind that in this density domain, $`\rho \rho _{c2}`$, the function $`E(\rho )`$ is determined by the repulsive part of the potential which makes $`t_3>0`$. As it was mentioned above, without this component of $`V(r)`$ the system’s energy would infinitely increase, $`E_{HF}\mathrm{}`$, with density growth, $`\rho \mathrm{}`$, thus inevitably collapsing. One could expect in principle the existence of metastable states at $`\rho >\rho _{c1}`$ if the potential $`V(r),`$ even being pure attractive, would have a complicated structure. It can be said that there could exist parameters of $`V(r)`$, which are able to open the possibility for the metastable states to be formed. On the other hand, a system of fermions interacting via a short-range, finite scattering-length, $`\delta `$-type potential $`V_\delta `$ , can be stable only in the dilute gas regime. While at the densities $`\rho \rho _{c1}`$ the incompressibility $`K`$ becomes negative, the system collapses. Indeed, the potential $`V_\delta `$ has no structure to ensure any metastable states at the densities $`\rho \rho _{c1}`$. As a result, one can write down a dimensionless expression for the ground state energy as a function of the only variable $`z=p_Fa`$ , $$\alpha E(z)=z^5(1+\beta (z)),$$ (4) with $`\alpha =10\pi ^2Ma^5`$. In the low density limit, $`\|ap_F\|1`$ and when the interaction has the radius $`r_0`$, eq. (4) reads , $$\alpha E(z)=z^5\left[1+\frac{10}{9\pi }z+\frac{4}{21\pi ^2}(112\mathrm{ln}2)z^2+\left(\frac{r_0}{a}\right)^3z^3\gamma (\frac{r_0}{a},z)+\mathrm{}\right].$$ (5) Here the function $`\gamma (y,z)`$ is of the order of one, $`\gamma (y,z)1`$. It is seen from eq. (5) that as soon as the scattering length becomes big enough, $`\|a\|r_0,`$ one can omit the contribution coming from the function $`\gamma `$ and neglect all the term proportional to $`(r_0/a)^3`$. Then eq. (5) reduces to eq. (4). Thus, in the case $`\|a\|\mathrm{}`$ we can use eq. (4) to determine the ground state energy $`E`$. Eq. (4) is valid up to the density $`\rho _{c1}`$ which is a singular point of the function $`\beta (z)`$, since beyond this point $`K<0`$, and the system is completely unstable. On the other hand, there is no physical reasons to have another irregular point in the region $`0\rho \rho _{c1}`$. We note, that eq. (5) is as well valid up to its own density $`\rho _{c1}^{^{}}\rho _{c1}`$, provided $`\|a\|r_0`$. Using eq. (3) for the incompressibility, one can calculate the position of the point $`z_{c1}`$ where $`K=0`$. Denoting the corresponding $`z`$ as $`z_{c1}=c_0`$, where $`c_0`$ is a dimensionless number, one is led to the conclusion that $`\rho _{c1}\|a\|^3`$ provided $`a`$ is sufficiently large to be the only dominating parameter. The system has only one stable region at small densities $`\rho \rho _{c1}`$ which decreases and even vanishes as soon as $`a\mathrm{}`$. One could expect that $`\|c_0\|\mathrm{}`$ as soon as $`a`$ becomes the dominant parameter so that the above given expression for $`\beta (z)`$ is valid in the whole domain $`\|z\|\mathrm{}`$. On the other hand, there exists another singular point $`z_{c2}`$ in the function $`E(\rho )`$ and the position of this point which corresponds to $`\rho _{c2}`$ depends mainly on the parameters such as $`r_0,`$ core radius $`r_c`$ of $`V(r)`$ rather than on the scattering length $`a`$. As to the function $`\beta (z),`$ by definition it does not contain any information about $`\rho _{c2}`$. Therefore, in order to take into account the existence of $`\rho _{c2}`$, one has to recognize that $`c_0`$ is finite, and the densities $`\rho _{c1}`$ and $`\rho _{c2}`$ are different. As a result, the function $`\beta `$ is determined in fact only in the region $`\|z\|\|z_{c1}\|`$. Now let us consider the calculation of the ground state energy $`E`$ of the system when the density approaches $`\rho \rho _{c1}\|a\|^3`$ from the low density side. As a rule, points $`\rho _{c1}`$ and $`\rho _{c2}`$ are missed in calculations because of the lack of the self consistency , which relates the linear response function of system with its incompressibility $`K`$, $$\chi (q0,i\omega 0)=\left(\frac{d^2E}{d\rho ^2}\right)^1.$$ (6) As we shall see below, these points can give important contributions to the ground state energy. To see it we express the energy of a system in the following form (see e.g. ), $$E(\rho )=T_k+E_H\frac{1}{2}\left[\chi (q,i\omega ,g)+2\pi \rho \delta (\omega )\right]v(q)\frac{d𝐪d\omega dg}{g(2\pi )^4},$$ (7) where $`\chi (q,i\omega ,g)`$ is the linear response function on the imaginary axis and $`v(q)=gV(q)`$, with $`V(q)`$ being the Fourier image of $`V(r)`$. The integration over $`\omega `$ goes from $`\mathrm{}`$ to $`+\mathrm{}`$, while the integration over the coupling constant $`g`$ runs from zero to the real value of the coupling constant, i.e. to $`g=1`$. At the point $`\rho =\rho _{c1}`$ the linear response function has a pole at the origin of coordinates $`q=0,\omega =0`$ due to eq. (6). At the densities $`\rho >\rho _{c1}`$ the function $`\chi (q,i\omega )`$ has poles at finite values of the momentum $`q`$ and frequencies $`i\omega `$. This prevents the integration over $`i\omega `$, making the integral in eq. (7) divergent. Thus, we conclude that it is the contribution of these poles that reflects the system’s instability in the density range $`\rho _{c1}\rho \rho _{c2}`$. Note, that violations of eq. (6) lead to serious errors in the calculation of the ground state energy. Eq. (7) can be rewritten, explicitly accounting for the effective interparticle interaction $`R(q,i\omega ,g)`$, (see e.g. ), in the following form $$E(\rho )=T_k+E_H\frac{1}{2}\left[\frac{\chi _0(q,i\omega )}{1R(q,i\omega ,g)\chi _0(q,i\omega )}+2\pi \rho \delta (\omega )\right]v(q)\frac{d𝐪d\omega dg}{g(2\pi )^4}.$$ (8) Here $`\chi _0`$ is the linear response function of noninteracting particles, while $`\chi `$ is given by the following equation $$\chi (q,\omega )=\frac{\chi _0(q,\omega )}{1R(q,\omega )\chi _0(q,\omega )}.$$ (9) It is seen from eqs. (6) and (9) that the denominator $`(1R\chi _0)`$ vanishes at $`\rho \rho _{c1}`$ while the radius of correlation tends to infinity . Thus, it is impossible to present the denominator as a power series in $`R\chi _0`$ approximating the expansion by the finite number of terms. This result is quite obvious since $`\rho _{c1}`$ is a singular point in the function $`E(\rho )`$ which makes it impossible to expand that function in the vicinity of this point. Therefore, one should try to satisfy eq. (6) in order to get proper results for the ground state calculations in the vicinity of the instability points. Such an approach was suggested in and is based on the exact functional equation for the effective interaction $`R(q,\omega ,g)`$, $$R(q,\omega ,g_0)=g_0v(q)\frac{1}{2}\frac{\delta ^2}{\delta \rho ^2(q,\omega )}\frac{\chi _0(k,iw)}{1R(k,iw,g)\chi _0(k,iw)}v(q)\frac{d𝐤dwdg}{g(2\pi )^4}.$$ (10) As a result, the linear response function $`\chi `$ given by eq. (9) automatically satisfies eq. (6) . Our preliminary calculations , based on eq. (9) and applied to the case when the scattering length is sufficiently large but finite, confirm the result that $`\rho _{c1}\|a\|^3`$. Let us suppose for a while that the bare potential is pure attractive. Then, the interval of the densities $`[0,\rho _{c1}]`$ within which the system is stable vanishes with the growth of $`\|a\|`$. As a result, in the limit $`a=\mathrm{}`$ the incompressibility becomes negative $`K0`$, making the considered system completely unstable at any density. Thus, the point at which $`a=\mathrm{}`$ is the only point of the system’s instability at all the densities. As soon as the scattering length deviates from its infinite value, that is $`+\mathrm{}>a>\mathrm{}`$ the system comes back to its stable state at list in the range of the density values $`\rho <\rho _{c1}\|a\|^3`$. It is of interest to understand whether it is possible to prove by e.g. numerical calculations, that $`\rho _{c1}\|a\|^3`$ when $`a\mathrm{}`$. From our point of view, at least at this moment, the answer is “no”. We are dealing with a system located in the vicinity of a phase transition, which transforms it into a strongly correlated one. As a result, it is hard to believe that the numerical calculations could be reliable. On the other hand, it is not really necessary to carry out numerical calculations if the problem allows a qualitative analysis. It was argued above, that there exists the only parameter to characterize the system which is the scattering length $`a`$. In fact the scattering length determines only the specific point $`\rho _{c1}`$ at which the incompressibility vanishes, separating the region of a dilute gas from the region of the system’s instability. As soon as $`a\mathrm{}`$ this last and the only parameters vanishes, driving the point $`\rho _{c1}\|a\|^3`$ of the curve $`E(\rho )`$ to the origin of coordinates. Thereafter, the system becomes unstable at all densities. And vice versa, as soon as the scattering length becomes finite the system is stable at list within the interval $`\rho \rho _{c1}\|a\|^3`$. Note, that as it follows from our consideration, any Fermi system possesses an equilibrium density and energy if the bare particle-particle interaction contains a repulsive core and its attractive part is strong enough, so that $`a\mathrm{}`$. Indeed, at sufficiently small densities the ground state energy is negative (since the incompressibility $`K0)`$ and the system will collapse until the core stops the density growth. Therefore, the minimal value of the ground state energy must be negative when the repulsive core will enter the play to prevent the system from the further collapse. It is worth to remark, that superfluid correlations cannot stop the system squeezing, since their contribution to the ground state energy being negative increases in the absolute value with the growth of the density. A liquid similar to the model one considered in this paper exists in Nature. This is liquid $`{}_{}{}^{3}He`$. If a helium dimer exists, its bound energy does not exceed $`10^4`$ meV while the ground state energy of helium liquid is about $`210^1`$ meV per atom . Because of this huge difference in binding energies, it is evident that there is no essential contribution coming from the binding energy of the dimer to the ground state energy of the liquid. In fact, the numerical calculations show that the pair potential is rather weak to produce the dimer $`He_2`$ . Thus, one can reliably consider an infinite homogeneous system of Helium atoms as consisting of particles interacting via pair potential, characterized by a very big but finite scattering length $`\|a\|r_0`$. Let us make also the following additional remark. It seems quite probable that the neutron-neutron scattering length ($`a20`$ fm) is sufficiently large to permit the neutron matter to have an equilibrium energy and density . Therefore, calculations of a neutron matter satisfying eq. (7) are quite desirable. In summary, the homogeneous system of interacting fermions was considered. It was shown that when the scattering length $`a`$ is negative and sufficiently large the fermion matter becomes a strongly correlated system at the densities $`\rho \|a\|^3`$. Therefore, the consideration of such a system is connected to a number of problems which yet persist and have to be resolved. At the same time, the qualitative consideration presented above gives strong evidences that the point $`\rho _{c1}`$ at which the incompressibility vanishes is defined by $`\rho _{c1}\|a\|^3`$ provided the scattering length is the dominant parameter of the problem. Thus, a homogeneous system composed of fermions, interacting via a pure attractive potential, at $`a\mathrm{}`$ is completely unstable at all the densities, with the incompressibility as a function of the density being always negative. As soon as the density $`\rho `$ goes to zero the incompressibility goes to zero as well. We thank G.F. Bertsch for attracting our attention to the discussed above many body problem which he has in fact suggested . One of us (VRS) is grateful to A. Bulgac for valuable discussions, and to the Department of Physics of the University of Washington where part of this work was done, for hospitality. This research was funded in part by INTAS under Grant No. INTAS-OPEN-97-603.
warning/0003/astro-ph0003118.html	ar5iv	text	# Determination of masses and radii of the massive eclipsing binary system HV 2543 in the Large Magellanic Cloud ## 1 Introduction The analysis of light curves of eclipsing binaries, in addition to radial velocity data, provides fundamental knowledge about the masses and physical dimensions of the stars. Studying massive binaries in the Magellanic Clouds we can learn about the evolution of these systems at metallicities lower than that of our galaxy. The Harvard variable HV 2543 ($`\alpha =5^\mathrm{h}27^\mathrm{m}27^\mathrm{s}`$, $`\delta =67\mathrm{°}11\mathrm{}54\mathrm{}`$, J2000) is a hot binary star in the LMC. It was catalogued as an OB star by Sanduleak (sk (1969)), who assigned it the identification –67$`\mathrm{°}`$117. The eclipsing nature of this binary was discovered by Gaposhkin (gap70 (1970)), who published a photographic light curve fitting a period of 4.829052 days (see also Payne-Gaposhkin pg (1971)). Photoelectric photometry was performed by Isserstedt (isser (1979)) who found $`V=12.92`$, $`(BV)=0.18`$ and $`(UB)=1.03`$. Radial velocity orbit was obtained by Niemela & Bassino (virpi (1994)) who derived physical parameters of the binary components and concluded that HV 2543 was a semidetached system with the less massive component filling its equipotential Roche surface. On the basis of their spectroscopic data, they classified this system as O8V:+O9III. Smith Neubig & Bruhweiler (neubig (1999)) have published an UV spectral classification of LMC OB stars based on IUE data, assigning to HV 2543 the type O9III. In this paper we present a CCD $`V`$ light curve for HV 2543. By means of the combined analysis of these data and previously published radial velocities, we derive new values of the fundamental parameters of this system. The paper is organized as follows: in Sect. 2 we describe the observations, reductions and calibration steps. In Sect. 3 we describe the photometric results and light curve fitting. In Sect. 4 we discuss the results and in Sect. 5 we present our conclusions. ## 2 Data Acquisition and Reductions ### 2.1 Observations The CCD images here analysed were acquired with the 2.15-m telescope at CASLEO, during three runs in 1995, 1997 and 1998. About five frames of the HV 2543 field were obtained each night, except during eclipses, when the observation was more intensive. During each observing night, a series of $`1015`$ twilight flats and bias frames was also obtained. In addition, for the 1997 and 1998 runs, $`1015`$ dome flats were also obtained each night. Tables 1 and 2 list details of the observations and Fig. 1 shows a finding chart for HV 2543. ### 2.2 Reductions and Photometry Reductions were performed using IRAF routines. For each night, the frames were reduced in the standard way (overscan and bias corrections, flat fielding) using properly combined bias and flat frames. The illumination differences between sky-flats and dome-flats were found to be less than 0.3 % at the edges of the frames, and negligible at the center. Whenever high quality sky-flats (that is, zero clouds) were available, these were preferred to dome-flats. For each observing run, the following steps were performed: 1. A master image was made by aligning and combining the highest quality frames (best seeing and darkest sky). On this combined image, a profile fitting photometry was performed using a stand-alone version of DAOPHOT II (Stetson pbs87 (1987), Stetson pbs91 (1991)). The resulting output was used as master coordinate list in the following steps. 2. Aperture photometry was performed on each individual frame using a circular diaphragm of $`9.5\mathrm{}`$ diameter. The sky background was determined locally for each star using a sky annulus of radii 9.76 and $`20.33\mathrm{}`$. 3. A list of instrumental magnitudes for each individual frame and for each star in the master list was constructed. Five faint neighbour stars were detected in the combined frame within a radius of $`9\mathrm{}`$ around HV 2543, being the brightest of these some 5 mag fainter than the binary. Light contamination from these stars slightly contributes to the measured magnitude of HV 2543 (differences of the order of 0.014 mag). Consequently, an adequate flux correction was applied to the aperture photometry. ### 2.3 Photometric calibration No single objects were chosen as comparison and check stars. Instead, zero point corrections were computed for each frame using a group of $``$15 stars, covering widely the range of magnitudes of HV 2543. To put all the observations in an unique instrumental system, we proceeded as follows: 1. A synthetic reference system was defined using the mean instrumental magnitudes of the calibration stars. 2. For each frame, a zero point offset was computed and the magnitudes were transformed to the synthetic reference system. 3. After this transformation, the RMS residuals were computed for each star. Only the stars with the smallest residuals ($`<0.01`$) were conserved as calibration stars. 4. We repeated these steps with the new, cleaned calibration star list. Since all nights were not really photometric, this procedure allowed us to have a direct estimate of the internal errors of each image. The above enumerated steps were performed independently for each observing run. We proceeded in this way because the stars available as calibrators are not the same for each run, due to pointing differences. On the other hand, some stars remain perfectly stable during a run, but show luminosity variations from one year to another, being then inadequate as calibrators between different runs. After this first stage of calibration was complete, a list of robust estimates of instrumental stellar magnitudes and their dispersions was constructed for each run. After cross identification between the three lists, the most stable stars were selected to serve as calibrators between the different runs. During the 1997 run, 19 frames through the Johnson’s $`B`$ band were also acquired. These frames were used to determine the transformations to the standard $`V`$ system. The final transformations to the standard system were performed by means of aperture photometry of 33 standard stars – ranging from $`(BV)=.004`$ to $`2.192`$ – from Landolt (lan (1992)) acquired during the photometric nights. We used the following transformation equations: $`b=B+b_1+b_2X+b_3(BV)`$ $`v=V+v_1+v_2X+v_3(BV)`$ We adopted the mean values of the extinction coefficients given by Minniti et al. (minni (1989)) and obtained the following values for the other coefficients: Nov. 19, 1997: $`b_1=4.2826\pm 0.0105`$ $`b_3=0.0208\pm 0.0115`$ $`v_1=3.4446\pm 0.0051`$ $`v_3=0.0456\pm 0.0055`$ Nov. 20, 1997: $`b_1=4.2986\pm 0.0076`$ $`b_3=0.03116\pm 0.0086`$ $`v_1=3.4502\pm 0.0057`$ $`v_3=0.0535\pm 0.0067`$ The RMS residuals of the $`V`$ transformations were 0.017 mag. We first shifted our unified instrumental system to that corresponding to a $`V`$ frame acquired immediately before a series of B exposures, on Nov. 19, 1997. Then, we applied the standard transformations derived for that night. The Nov. 20 transformations were used only to check the quality of the first ones. Our photometry of HV 2543 is presented in Table 3. In successive columns the heliocentric julian day, standard V magnitude, internal errors, seeing and airmass are given. The errors ($`\sigma _\mathrm{i}`$) account for the internal photometric errors estimated by DAOPHOT and those of the transformation to an unified instrumental system. They do not include the error of the transformation to the standard system. ## 3 Results ### 3.1 Ephemeris We have observed one primary minimum during 1997, at HJD$`=\mathrm{2\hspace{0.17em}450\hspace{0.17em}768}\stackrel{d}{.}802\pm 0.001`$. Combining this value with that given by Payne-Gaposchkin (pg (1971)) we obtained a new ephemerides for HV 2543: $`E_0=\mathrm{2\hspace{0.17em}450\hspace{0.17em}768}\stackrel{d}{.}802\pm 0.001`$ $`P=4\stackrel{d}{.}\mathrm{829\hspace{0.17em}046}\pm \mathrm{0.000\hspace{0.17em}004}`$ We have used this ephemeris in calculating the phases for the radial velocity and photometric data. ### 3.2 Light and radial velocity curve analysis The $`V`$ light curve of HV 2543 is typical of near-contact binaries, with different depths of the minima and continuous variations of light in out-of-eclipse portions. Also noticeable is the O’Connell effect, i.e., the small difference in magnitude between both maxima (see Davidge & Milone damilo (1984)). We solved simultaneously the light and radial velocity curves using the Wilson & Devinney code (hereafter WD), which is very well suited for the study of close binary systems (Wilson & Devinney wd (1971), Wilson wil (1990)). Unit weight was given to photometric CCD observations since they all were collected with the same instrumental configuration and have comparable quality. A unit weight was also given to most of the radial velocity points except for those near eclipses or indicated by Niemela & Bassino (virpi (1994)) as less confident. We also corrected the data points corresponding to the phase $`\varphi =0.05`$, that are inverted (i.e., the O8 velocity corresponds to the O9 and vice-versa) in Niemela & Bassino (virpi (1994)). The relation between photometric and spectroscopic weights was given through the values of the mean standard deviations (sigma) of the data which were estimated to be 17 and 12 km s<sup>-1</sup> for the radial velocities of the primary and secondary components, respectively, and 0.02 mag for the light points. Only the ratios between these values are significant. The radii that result from our preliminary light curve analysis indicate that the O8 component is more alike an O8III star. Chlebowski & Garmany (chg (1991)) give a temperature of 36000 K for an O8III star. Given that Schmidt-Kaler (1982) gives a lower temperature for such spectral type, we also performed an analysis assigning a 34700 K temperature for the primary. Standard values of bolometric albedos, $`A=1.0`$ (Rucinski ruc (1969)), and gravity darkening coefficients, $`g=1.0`$ (Lucy luc (1976)), for radiative envelopes were used. Linear limb-darkening coefficients were determined from tables by van Hamme (ham (1993)). These parameters were not adjusted. The adjustable parameters in our computations were: $`a`$ (semimajor axis), $`V_\gamma `$ (systemic radial velocity), $`i`$ (orbital inclination), $`T_2`$ (temperature of the secondary component), $`q`$ (mass-ratio), $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ (modified potential of both components) and $`L_1`$ (luminosity of the primary). The first approaches to fit the observations were made by means of the light curve (LC) program until an acceptable fit to the $`V`$ light curve was obtained. Then we proceeded with the differential corrections (DC) code. We first adjusted the parameters $`a`$, $`V_\gamma `$ and $`q`$ using only the out-of-eclipses radial velocity data. Thereupon, we left those parameters fixed and fitted $`i`$, $`T_2`$, $`\mathrm{\Omega }_1`$, $`\mathrm{\Omega }_2`$ and $`L_1`$, using only the light curve. Then we re-computed the first set of parameters leaving fixed the second set, and iterated this procedure until the corrections were smaller than their errors. A better solution was achieved allowing the fit of $`q`$ with the photometric data, that is, including $`q`$ within the second parameter set. Mode 2 (Leung & Wilson lw (1977)) which stands for detached systems was used at the beginning, but after a few runs of the DC program the computations clearly evolved towards a semidetached configuration with the secondary component filling its critical Roche lobe. Thus, mode 5 of the WD code was employed until the final solutions were found. In this mode, the parameter $`\mathrm{\Omega }_2`$ cannot be adjusted and it is set equal to the critical filling-lobe value. The solutions never converged to contact configurations. Since some systematic differences were obtained between the observed and modelled light curves (due to the O’Connell effect), we tried to improve the binary modelling by including one or two hot spots on the surface of the stars. The spots were placed at the equator of the stars (colatitude $`=90\mathrm{°}`$) while the other spot parameters where adjusted with the WD code. The range of the adjusted temperature for the secondary star, $`T_2`$, depends on the adopted value for $`T_1`$. To estimate the errors of the other parameters we considered the differences between the values that arised from different solutions (i.e., using the “spectroscopic” value of $`q`$, 0.55, or leaving it to increase until 0.64, to allow a better fit to the photometric data). The adopted solution, with $`q=0.61`$, is a compromise between the optimal fit of the light curve and the tolerable deviation from the spectroscopic data. Table 4, lists the model parameters. Fig. 2 and 3 show the modelled light and radial velocity curves, respectively, derived from the one spot solution and plotted along with the observations, and their corresponding residuals (O-C). The final value of the orbital inclination is incidentally high, very close to $`90\mathrm{°}`$. This fact should produce a total eclipse at the secondary minimum, but, as the dimensions of both components are quite similar, this is not a noticeable feature of the observed light curve. Absolute values of the masses, dimensions and bolometric magnitudes of HV 2543 were computed from the one spot solution. They are listed in Table 5 where $`R_1`$ and $`R_2`$ correspond to mean values of the derived polar, back and side radii of the stars. The final configuration of HV 2543 shows two components of similar dimensions but different masses and temperatures, with the secondary less massive star filling its Roche lobe. ### 3.3 $`BV`$ Field photometry We combined the $`B`$ frames acquired during 1997 to make a master $`B`$ image, on which we performed a profile fitting photometry. We used this photometry together with that performed on the $`V`$ master frame to generate a colour-magnitude diagram, which is shown in Fig. 4. In this figure, the presence of an OB association covering the whole frame ($`9\mathrm{}`$ of diameter) is evident. Ostrov et al. (yoetal (1999)) found evidence of a stellar association surrounding the massive binary system Sk–67$°$105. The relatively small angular distance between the two stellar groups ($`9.5\mathrm{}`$) suggests that they are probably related. We have obtained spectroscopic data for five stars in the neighbourhood of Sk-67$`\mathrm{°}`$105. For these stars, we have estimated the reddening by comparison between the intrinsic $`(BV)_0`$ colours given by Schmidt-Kaler (s-k (1982)) and FitzGerald (fg (1970)) and the observed colours. We derived $`E(BV)=0.17\pm 0.015`$, while for HV 2543 itself we obtained $`E(BV)=0.20`$. A detailed study of these OB associations will be presented in a forthcoming paper. ## 4 Discussion We note that the radius derived for the O8 component is substantially larger than the one obtained by Niemela & Bassino (virpi (1994)). In fact, the spectroscopic data suggest that the luminosity of the O9 star is larger than that of the O8 star, which caused Niemela & Bassino to refer to the O9 component as “primary”, even though it is the one presenting the largest radial velocity amplitude, and consequently the less massive component of the binary system. Hence, we explored alternative solutions to the light curve that could yield a smaller radius for the O8 component. The values $`R_112.75`$, $`R_215.31`$, $`M_121.91`$ and $`M_219.34`$, in solar units, provide a good fit to the light curve, although they require a value of $`q`$ near 0.9, which is not compatible with the radial velocity data. Hence we have discarded this solution, which on the other hand, implies a rather low value of the distance modulus, namely 18.18, assuming $`A_V=0.62`$ (see below) and adopting the temperature scale of Chlebowski & Garmany (chg (1991)). This value would result even smaller if a lower temperature scale is adopted. On the other hand, the temperature difference between the two components of HV 2543 resulting from the light curve analysis is larger than that suggested by the corresponding spectral types. This fact might be due to the difference between the spectral features of the unperturbed back sides and the heated inner sides of the stars. Also should we have in mind that the mass transfer history of the system might account for significant departures of the He abundances relative to those regarded as normal for LMC members, and this effect could influence some spectral lines of the secondary component. To solve these puzzles, high dispersion spectroscopy would be desirable. The semimajor axis and star dimensions of this system are alike those determined by Pritchard et al. (prit (1998)) for HV 2241, but the masses are somewhat smaller. We presume that these systems have experienced case A mass transfer, being now in the slow stage of mass exchange. In these cases, the mass gainer should be indistinguishable of a normal star, excepting that it would stand on an isochrone corresponding to a shorter age (see Vanbeveren et al. vbv (1998)). From the stellar models of Schaerer et al. (sch (1993)) (for single stars) we found that a $`25\mathrm{M}_{\mathrm{}}`$ star takes $`6.7`$ Myr until its radius grows to $`15\mathrm{R}_{\mathrm{}}`$, but such star would have an effective temperature of only 30000 K. The radius and effective temperature derived for the O8 star are consistent whit those of a single star of some 3.7 Myrs and $`40\mathrm{M}_{\mathrm{}}`$, completely out of the range of masses compatible with the radial velocity data. We note that a similar problem arises from the analysis of AB Crucis (Lorenz et al. lor (1994)). It is clear that there are a series of phenomena that we do not fully understand, and consequently, our analysis can not be considered definitive. The errors given for the derived parameters must be considered with caution. True errors are not easy to estimate analytically, since they depend on the importance of phenomena that are not properly accounted for, such as wind shocks, uncertainty of the adopted temperature scale, radiation pressure effects, etc. In fact, the sizes and positions of the spots used to model empirically the O’Connell effect are rather arbitrary, and an equally satisfactory solution could be found with other parameters, but these details do not affect meaningly the derived star dimensions. From our photometry we determine $`(BV)=0.11\pm 0.015`$ for HV 2543. This value is somewhat redder than $`0.18`$, obtained by Isserstedt (isser (1979)). However, given that Isserstedt does not detect the variability of HV 2543, we think that both measures are still in reasonable agreement. If we assume a $`(BV)_0`$ of -0.31 (Schmidt-Kaler s-k (1982)), $`R=3.1`$ (Koornneef koo (1982)) then it results an $`A_V=0.62`$ for HV 2543. This value does not depend on which temperature scale we adopt, since for the range of temperatures of the O-type stars the $`(BV)`$ colours are degenerated. Estimating the bolometric corrections according to Massey & Hunter (mh (1998)), we derive a distance modulus of $`(mM)_0=18.31\pm 0.2`$ to $`18.40\pm 0.2`$, depending on the adopted temperature scale. The error accounts for the uncertainties in the estimates of $`R`$ and the bolometric corrections. This distance modulus must be considered with caution, since this system has experienced strong mass transfer and exhibits the above mentioned anomalies. ## 5 Conclusions We have presented a CCD $`V`$ light curve of the eclipsing binary HV 2543 (Sk-67$`\mathrm{°}`$117) in the LMC. The light curve appearance is almost symmetric, with a slight ($`0.02V`$ magnitudes) O’Connell effect. From $`BV`$ photometry of the surrounding field, we found that HV 2543 probably belongs to an OB association not previously identified. We have analysed this light curve and the published radial velocity data using the Wilson-Devinney code, finding that HV 2543 is a semidetached system with the less massive and less luminous member filling its Roche lobe. The primary minimum occurs when the O8 component is behind the O9. From our analysis, we derived fundamental parameters for this system, obtaining $`M_1=25.63\pm 0.7`$, $`M_2=15.63\pm 1.0`$, $`R_1=15.54\pm 0.4`$ and $`R_2=13.99\pm 0.4`$ in solar units. Some discrepancies between the present results and those derived from the previous spectroscopic analysis (Niemela & Bassino virpi (1994)), should be thoroughly addressed in future works. We found that this system has experienced significant mass exchange, being the present most massive star the originally less massive. ###### Acknowledgements. The authors acknowledge use at CASLEO of the CCD and data acquisition system supported under U.S. National Science Foundation grant AST-90-15827 to R. M. Rich. The focal reducer in use at CASLEO was kindly provided by Dr. M. Shara. This research has made use of the Astronomical Data Center catalogs. We are indebted to the staff of CASLEO for valuable help during the observing runs. We are grateful to V.S. Niemela for her vigorous comments made during the preparation of this paper. We would also like to thank to the referee for the suggestions that allowed to improve the presentation of the paper.
warning/0003/astro-ph0003422.html	ar5iv	text	# A MICROLENSING SEARCH FOR COLD MOLECULAR CLOUDS IN VIRGO ## 1 Rationale The number of microlensing surveys underway is growing rapidly. Most of these surveys aim to detect compact objects in the dark halo of our own Galaxy or in the nearby M31 halo by monitoring stars in the LMC, SMC, or M31 itself. These projects have proved very fruitful detecting and placing limits on the dark matter in the Galaxy halo. The latest results indicate that less than $`20\%`$ of the Galaxy halo is in the form of MACHOs (see refs 1 and 2). There are, however, a number of shortcomings of the current surveys. Firstly, one cannot measure the distance to the MACHO and there is currently some controversy over whether the events seen towards the LMC are due to Galactic dark matter at all. Secondly, the objects have to be extremely compact to produce a microlensing effect at the distance of the LMC. Hence, current surveys provide no constraint on the hypothesis that the dark matter is composed of diffuse cold molecular hydrogen clouds (see e.g. refs 3 and 4). Thirdly, source blending greatly complicates the calculation of survey efficiencies. We have begun a project which addresses these issues and provides completely new information on the nature of dark matter in galaxy clusters. We are monitoring 645 quasar candidates behind the Virgo galaxy cluster in order to detect dark objects in the cluster in the mass range $`5\times 10^7`$ $`M_{}`$$`5\times 10^3`$ $`M_{}`$. The upper mass limit is set by the length of the monitoring program and the lower limit by the frequency of the monitoring. Quasar source size can also have an effect on the lower mass limit. If the quasar’s angular size is larger than the angular diameter of the Einstein ring of the lens, the magnification drops significantly and the event would not be detected. If all the mass in the Virgo cluster were in dark objects within the above mass range then we can expect to see significant numbers of microlensing events during our monitoring program (see ref 5 for details). Several features of our monitoring program deserve to be emphasised. Firstly, if we see a microlensing event we have a reasonably accurate measure of the distance to the MACHO, allowing us to calculate its mass. Secondly, we can detect objects down to surface mass densities of a few tens of $`\mathrm{g}\mathrm{cm}^2`$ as these would act as point mass lenses when placed in Virgo. Such objects, for example cold molecular gas clouds, are accessible to no other current microlensing survey. Thirdly, because we are monitoring $`600`$ objects over a 30 sq. deg. field, we do not have blending of source images and the associated problems that blending causes. One potential worry for a survey such as ours is the fact that quasars are intrinsically variable. If quasar variability is large on the time scales over which we are monitoring we would have to set our microlensing magnification threshold so high that the expected number of events would become very low. To date, little systematic quasar monitoring has been carried out on time scales matching those of our project, thus we have been careful to examine the feasibility of our project. As reported in this article, our pilot project proves that intrinsic quasar variability is not a serious problem on these time scales. A by-product of this microlensing survey is the information provided on the variability of quasars over these timescales, which will lead to constraints on models of quasar fuelling mechanisms. ## 2 Observations and data analysis During the period 1999 February - June we obtained 28 R-band Schmidt plates of the central 30 sq. deg. of the Virgo cluster. The data are complete to an R-band magnitude of $`20`$. Using previous plates of the same area in the U and B bands we identified $`645`$ candidate quasars according to their colour. Of these we expect $`\mathrm{}<10\%`$ to be contaminant stars. Figure 1 shows $`2`$ typical quasar light curves extracted from this data set. It can be seen qualitatively that the quasar candidates vary little during the monitoring period. We have quantified this by comparing the level of variability in the quasar candidates with the variability of $`3000`$ stars in our field. The candidate quasars are statistically no more variable than the stars. This confirms the feasibility of using quasars for microlensing surveys for monitoring periods of a few months. Our simulations show that with this low level of variability we can set a low detection threshold, and identify microlensing events down to a peak magnification of 0.46 mag. We have run a matched filter analysis on each of the $`645`$ candidate-quasar light curves. The filter used was a MACHO light curve of peak amplification $`0.6`$ mag. We used a range of values for the event time-scale and the time of peak magnification. For each quasar light curve the event with the highest $`\mathrm{\Delta }\chi ^2`$ for a match to a MACHO event was retained. To this best event for each quasar, we performed a four parameter MACHO light curve fit. To define selection criteria to pick out likely lensing events, we ran the same set of procedures on the light curves of the 3000 control stars in the field – where we do not expect any real microlensing events. Figure 2 shows the distribution of points in the plane of $`\mathrm{\Delta }\chi ^2`$ for the four parameter MACHO fit against the fitted impact parameter $`b_{min}`$. Quasar candidates are shown by crosses and stars as stars. We define the following criteria for an event to be considered as a serious microlensing candidate: $``$ $`\mathrm{\Delta }\chi ^2>100`$ $``$ $`\text{Impact parameter}b_{min}<0.8`$ These criteria leave two events in the candidate region to be checked. Neither of these proved to be real events: one is a known highly variable BL Lac object and the other a probable stellar RR Lyrae contaminant. From this null detection we can now put constraints on the nature of dark matter in Virgo. ## 3 Detection Efficiency and Predicted Number of Events To calculate the detection efficiency of the survey we have generated artificial microlensing events and added them to the 3000 stellar light curves. For each event time-scale we have generated $`10,000`$ artificial light curves with impact parameters distributed uniformly between $`0`$ and $`1`$ Einstein ring radius and added them to the stellar light curves, placed at random dates. We then applied our matched filter analysis, the fitting routine, and the selection criteria to these artificial light curves and measured how many of the artificial events were recovered. Our detection efficiency is shown as a function of event time-scale on the left hand side of Figure 3. Using the detection efficiency we can calculate the number of microlensing events we would expect to see in our experiment if the dark matter in Virgo were in the form of dark objects of a particular mass. We assume that the Virgo cluster can be modelled as an isothermal sphere with a one-dimensional velocity dispersion of $`673\mathrm{km}\mathrm{s}^1`$. Our expected number of events is shown on the right hand side of Figure 3 and peaks at $`5`$, if all the dark mass in Virgo were in the form of $`10^5`$$`M_{}`$objects. ## 4 Conclusions and future work From our null detection, and the fact that we would have expected $`5`$ events for lenses of mass $`1\times 10^5`$ $`M_{}`$, we can use Poisson statistics to put constraints on the mass in Virgo in the form of dark objects of this mass. The result of this pilot project is that less than $`1/2`$ the mass in Virgo can be in the form of $`10^5`$ $`M_{}`$objects at $`90\%`$ confidence. The most important thing we have learnt from the pilot project is that a quasar monitoring microlensing project is technically feasible and we hope, over the next 2 years, to acquire of order $`100`$ more plates of Virgo. Our simulations indicate that this will allow us to improve our constraints on the nature of dark matter in Virgo by a factor $`510`$. A second season of monitoring Virgo started in February 2000. We would also like to extend the experiment by monitoring other clusters - especially the Perseus cluster in the North as this is the most massive nearby cluster. Alternatively, another strategy would be to monitor several more distant clusters (thereby avoiding the need for a very wide field of view). In the long term we would also like to monitor a control field of quasars with no intervening massive cluster in order to precisely quantify quasar variability. ## Acknowledgements We are immensely grateful to the Schmidt observing team: Fred Watson, Malcolm Hartley, Russell Cannon, Paul Cass, Ken Russell and Delphine Roussiel for taking all the photographic plates of Virgo used in this analysis. ## References
warning/0003/cond-mat0003311.html	ar5iv	text	# Coulomb gap in a model with finite charge transfer energy. ## I INTRODUCTION Doping of solids might lead to drastic qualitative changes in their properties. The metal-insulator transition (MIT) is a spectacular manifestation of this. The understanding of the driving forces of the MIT is a long-standing problem. In the early seventies, the prediction was made that on the dielectric side of the MIT the long-range Coulomb interactions deplete the density of one-electron excitations (DOE) $`g(\epsilon )`$ near the Fermi energy $`\mu `$. Further, analytical calculations of $`g(\epsilon )`$ with Coulomb correlation taken into consideration have been performed on the metallic side of the MIT. Altshuler and Aronov showed that for the metallic case $`g(\epsilon )`$ in three dimensions has a cusp-like dependence $`g(\epsilon )\|\epsilon \mu \|^{1/2}`$ near $`\mu `$. This was later confirmed in electron tunneling experiments for amorphous alloys and granular metals . On the insulating side of the MIT charge transport occurs via inelastic electron tunneling hopping between states localized on the impurity sites with one-electron energies close to $`\mu `$. Mott demonstrated that at low temperatures electrons seek accessible energy states by hopping distances beyond the localization length, leading to a hopping conductivity $`\sigma (T)\mathrm{exp}(T_0/T)^\nu `$ with $`T_0`$ being a characteristic temperature depending on localization length and with the hopping exponent $`\nu =1/4`$ for the non-interacting case in three dimensions. Efros and Shklovskii (ES) argued that the ground state of a system with long-range Coulomb interactions is stable with respect to one-particle excitations only if $`g(\epsilon )`$ in the vicinity of $`\mu `$ has the symmetric shape $$g(\epsilon )\|\epsilon \mu \|^{D1}$$ (1) with the universal exponent $`D1`$ depending only on the dimensionality $`D`$ of the system. In particular, ES predicted that in $`D=3`$ $`g(\epsilon )=\frac{3}{\pi }\left(\frac{\chi }{e^2}\right)^3(\epsilon \mu )^2`$, where $`\chi `$ is the dielectric constant and $`e`$ is the electron charge. Because $`g(\epsilon )`$ vanishes only at $`\epsilon =\mu `$, this is called a “soft” Coulomb correlation gap with a width $`\mathrm{\Delta }\epsilon e^3(N_0/\chi ^3)^{1/2}`$, where $`N_0`$ is the DOE far away from $`\mu `$. The power law (1) gives a hopping exponent $`\nu =D/(D+3)`$ at low temperatures, so for three-dimensional system with long-range Coulomb interactions $`\nu =1/2`$. The intriguing hypothesis about universality of (1) has stimulated further theoretical research, both analytical and numerical . To establish the hypothesis (1) Efros used the ground-state stability conditions for localized electrons (LES) with respect to charge transfer $$\epsilon _j\epsilon _i\frac{e^2}{\chi r_{ij}}>0,$$ (2) where $`\epsilon _i`$ and $`\epsilon _j`$ are the one-particle energies of a neutral donor on a site $`i`$ and of a charged donor on a site $`j`$, respectively, and $`r_{ij}`$ is the distance between the sites $`i`$ and $`j`$. The conditions (2) were used to heuristically derive a non-linear integral equation for $`g(\epsilon )`$ and then assymptotic analysis of this equation leads to the power law (1). LES have been studied using the so-called classical donor-acceptor ($`d`$-$`a`$) model (see, e. g. Ref.). Within this model, the system considered is modeled by a continuous sample with randomly distributed $`k\times N`$ ($`k1`$) acceptor and $`N`$ donor sites. Each acceptor site is negatively charged whereas out of $`N`$ only $`k\times N`$ donors have a positive charge which leads to a large number of configurations of charged donors. Moreover, each of these configurations must obey not only conditions (2) but also more complicated conditions related to many-particles excitations (e. g., charge transfer involving four, six, etc. sites). Efros conjecture about the universality implies that $`g(\epsilon )`$ does not depend on peculiarities of the particular model and, as a consequence, further theoretical studies of LES were confined to a lattice $`d`$-$`a`$ model proposed in Ref.. In this model, $`N`$ donors are localized on all the sites of a D-dimensional lattice and the negative charge from $`k\times N`$ acceptors is uniformly smeared over the lattice sites so that each site $`i`$ has a charge $`e(n_ik)`$, where $`n_i=1`$ if a donor on the site $`i`$ is ionized and $`n_i=0`$ if a donor is neutral. Disorder in this model is ensured by introducing randomly distributed one-site potentials. Monte Carlo simulations on very large specimens of the lattice $`d`$-$`a`$ model, however, have given rise to doubts about the universality of the $`g(\epsilon )`$ behavior. Another hint about possible non-universal behavior of $`g(\epsilon )`$ has come from the intriguing and still not completely unfolded problem whether the so called spin-glass phase does exist in the classical $`d`$-$`a`$ model (see, e. g. Ref.). Grannan and Yu studied the classical three-dimensional $`d`$-$`a`$ model with $`k=0.5`$ but with the total acceptor charge uniformly distributed over donor sites as in the lattice $`d`$-$`a`$ model. In this case, the classical $`d`$-$`a`$ model is equivalent to a model of Ising spins, localized on randomly distributed sites, with pairwise Coulomb interactions, a model in which a transition into the spin-glass state was found to occur at non-zero temperature. It was then concluded that such a transition should exist in all $`d`$-$`a`$ models (with and without smearing of negative charge, defined on a lattice or on a continuous sample) as well because of the Efros universality hypothesis. Voita and Schreber , however, have shown that the spin glass transition does not exist in the lattice $`da`$ model. Besides, in recent work by one of us it was unequivocally demonstrated that the ground state of the classical $`d`$-$`a`$ model and that of the model studied in Ref. are qualitatively different. An analysis of histograms $`[Q_{\alpha \beta }]`$ of the so called overlaps $`Q_{\alpha \beta }=\frac{1}{N}_i\delta (n_i^\alpha ,n_i^\beta )`$ (here $`\alpha `$ and $`\beta `$ refer to different pseudo-ground states (PGS) obtained by direct descents) has revealed that, indeed, for the model studied in Ref. $`[Q_{\alpha \beta }]`$ has a symmetric Gaussian shape with the maximum at $`Q_{\alpha \beta }=0`$ and with the dispersion $`Q_{\alpha \beta }^2N^1`$. This means that a large number of microscopically different PGS’s does exist in the model and according to Parisi’s theory this implies the existence of a spin-glass state at low temperatures. Further Monte Carlo simulations at finite temperatures revealed the typical finite-size scaling of the spin-glass susceptibility. In the classical $`d`$-$`a`$ model, however, $`[Q_{\alpha \beta }]`$ has its maxima at $`Q_{\alpha \beta }=1`$ which means that all PGS’s generated are the same from microscopical point of view. The absence of microscopically different PGS’s in the classical $`d`$-$`a`$ model was explained by the pinning of all PGS’s on the electric field created by the discretely distributed acceptor charges. Therefore, it is highly desirable to study the properties of not only the classical $`d`$-$`a`$ model, but of its various modifications as well. In the present work we consider a modified classical $`d`$-$`a`$ model (MCDAM) in which acceptors can be neutral, so the energy $`\mathrm{\Delta }`$ of the charge transfer from a donor to an acceptor ($`d^0+a^0d^++a^{}`$, where $`d^0`$, $`d^+`$, $`a^0`$, $`a^{}`$ stand for a neutral donor, a charged donor, a neutral acceptor, a charged acceptor, respectively) has to be finite. The classical $`d`$-$`a`$ model might be then viewed as the limit of the MCDAM as $`\mathrm{\Delta }\mathrm{}`$. We have investigated the shape of the Coulomb gap (i. e. $`g(\epsilon )`$ both for donor and acceptors in the vicinity of the Fermi level) in two- and three-dimensional MCDAMs at T=0 and found that the behavior of $`g(\epsilon )`$ is in strong contradiction to the Efros conjecture about the universality of $`g(\epsilon )`$. The rest of the paper is organized as follows. In Section II we introduce the MCDAM and arrive at some rigorous results which follow from a symmetry of the MCDAM with respect to the exchange of donor and acceptor sites. Further, the algorithm of energy minimization for the MCDAM including a discussion about inherent finite size effects is presented in Section III. Section IV is devoted to a description of the main results obtained. In Section V we discuss possible causes of universality violation in the MCDAM, analyze experimental data available in the literature and predict possible experimental situations in which the non-universal behavior of $`g(\epsilon )`$ might be observed. And finally, a summary is presented in Section VI. ## II BACKGROUND ### A Model We consider a $`D`$-dimensional system of volume $`L^D`$, in which an equal number of acceptor and donor sites $`N`$ are allocated according to the Poisson distribution with a density $`n=N\times L^D`$. It is convinient to choose energy unit $`E_0`$ as an energy of the Coulomb interaction between a pair of acceptors, say, localized on the average distance $`n^{1/D}`$, $`E_0=e^2n^{1/D}/\chi `$. In typical bulk semicondustors $`n10^{18}`$ cm<sup>-3</sup> and $`\chi 10`$, so $`E_00.02`$ eV. Hereafter all expressions will be written in dimensionless units $`n^{1/D}`$ for length and $`E_0`$ for energies. A microscopic state of a particular spatial arrangement of the donor and acceptor sites (henceforth referred to as the sample R) is determined by a set of occupation numbers $`(n_a,n_d)\{n_a(i),n_d(k),i=1,2,\mathrm{},N`$, $`k=1,2,\mathrm{},N\}`$ determined in the following way. For the acceptors, $`n_a(i)=1`$ if an acceptor on an acceptor site $`i`$ has captured an electron and $`n_a(i)=0`$ if an acceptor is neutral. For the donors, $`n_d(k)=1`$ if a donor on a donor site $`k`$ is neutral and $`n_d(k)=0`$ if a donor has given an electron away. We investigate the LES from the dielectric side of the MIT, so $`\alpha _B<1`$ ($`\alpha _B`$ is the localization length of the electron on donor). The energy of the sample, assuming that all the interactions are of Coulomb origin, then is $`E(n_a,n_d)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}{\displaystyle \frac{n_a(i)n_a(j)}{r_{ij}^{aa}}}+`$ (3) $`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{kl}{}}{\displaystyle \frac{(1n_d(k))(1n_d(l))}{r_{kl}^{dd}}}`$ (4) $``$ $`{\displaystyle \underset{i,k}{}}{\displaystyle \frac{(1n_d(k))n_a(i)}{r_{ik}^{ad}}}\mathrm{\Delta }{\displaystyle \underset{i}{}}n_a(i),`$ (5) where indices $`i,j`$ and $`k,l`$ number acceptor and donor sites, respectively, $`r_{ij}^{aa}`$, $`r_{kl}^{dd}`$ and $`r_{ik}^{ad}`$ are the distances between the acceptors on the sites $`i`$ and $`j`$, between the donors on the sites $`k`$ and $`l`$, and between the acceptor on the site $`i`$ and the donor on the site $`k`$, correspondingly, and $`\mathrm{\Delta }`$ is the energy of charge transfer between acceptor and donors. When charge transfer occurs in the system, the energy of the sample changes by $`\delta E(n_a,n_d)`$ $`=`$ $`{\displaystyle \underset{i}{}}\epsilon _a(i)\delta n_a(i)+{\displaystyle \underset{k}{}}\epsilon _d(k)\delta n_d(k)+`$ (6) $`+`$ $`{\displaystyle \underset{i,k}{}}{\displaystyle \frac{\delta n_a(i)\delta n_d(k)}{r_{ik}^{ad}}}+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}{\displaystyle \frac{\delta n_a(i)\delta n_a(j)}{r_{ij}^{aa}}}+`$ (7) $`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{kl}{}}{\displaystyle \frac{\delta n_d(k)\delta n_d(l)}{r_{kl}^{dd}}},`$ (8) where $`\epsilon _a(i)`$ is the one-electron excitation (OEE) energy for the acceptors $$\epsilon _a(i)\frac{\delta E(n_a,n_d)}{\delta n_a(i)}=\underset{ji}{}\frac{n_a(j)}{r_{ij}^{aa}}\underset{k}{}\frac{1n_d(k)}{r_{ik}^{ad}}\mathrm{\Delta },$$ (9) $`\epsilon _d(i)`$ is the corresponding OEE energy for the donors and $`\delta n_a(i)`$ ($`\delta n_d(k)`$) denotes the change of the occupation number on the acceptor (donor) site $`i`$ ($`k`$). If a microscopic state ($`n_a^0`$,$`n_d^0`$) of the sample is the ground-state of this sample then for any excitation the relation $$\delta E(n_a^0,n_d^0)0$$ (10) holds. The specific appearance of the conditions (10) depends on what excitations are allowed in the model system considered. In the present paper we investigate the simplest case when only pairs of sites are involved in the charge transfer which, in turn, is allowed to occur in four different ways: (i) via electron hops between a pair of the acceptors $`\{n_a(i)=1,n_a(j)=0\}\{n_a(i)=0,n_a(j)=1\}`$; (ii) via electron hops between a pair of donors $`\{n_d(k)=1,n_d(l)=0\}\{n_d(k)=0,n_d(l)=1\}`$; (iii) via ionization process $`\{n_a(i)=0,n_d(k)=1\}\{n_a(i)=1,n_d(k)=0\}`$ and (iv) via recombination process $`\{n_a(i)=1,n_d(k)=0\}\{n_a(i)=0,n_d(k)=1\}`$. For each of those processes there is an unique set of $`\{\delta n_a(i),\delta n_a(j),\delta n_d(k),\delta n_d(l)\}`$. For instance, for the acceptor-acceptor hops $$\delta n_a(i)=1,\delta n_a(j)=1,\delta n_d(k)=0,\delta n_d(l)=0.$$ (11) Substituting (11) into (8) one obtains the ground-state stability relation with respect to the charge transfer between the pair of acceptors on the sites $`i`$ and $`j`$ $$\epsilon _a^0(j)\epsilon _a^1(i)\frac{1}{r_{ij}^{aa}}0,$$ (12) where $`\epsilon _a^{1(0)}(i)`$ denotes $`\epsilon _a(i)`$ if $`n(i)=1(0)`$. The stability conditions with respect to the other three manners of the charge transfer are obtainable in the similar manner. The relation (12) implies that $`\epsilon _a`$’s for the neutral acceptors are, in general, larger than $`\epsilon _a`$’s for the charged acceptors. Furthermore, the pair of neutral and charged acceptors might be located on any distance and therefore in the thermodynamic limit the chemical potential for the acceptors (i. e. an energy level which separates the energies of the neutral and charged acceptors) is determined as $$\mu _a=\text{min}\{\epsilon _a^0(i)\}=\text{max}\{\epsilon _a^1(i)\}.$$ (13) Alike, there exist the chemical potential $`\mu _d`$ for the donors as well. Moreover, the stability relations with respect to the ionization and recombination lead to $$\mu _a=\mu _d=\mu .$$ (14) Despite the finite size of samples we investigated, the relation (14) with the chemical potentials calculated from (13) is valid within the limits of accuracy of our calculations (see Sect. III). A macroscopic state of the sample R is characterized by degree of acceptor ionization $$C_a(𝐑)=\frac{1}{N}\underset{i}{}n_a(i),$$ (15) by the DOE for acceptors $$g_a(\epsilon _a,𝐑)=\frac{1}{N}\underset{i}{}\delta (\epsilon \epsilon _a(i))$$ (16) and by the corresponding DOE $`g_d(\epsilon _d,𝐑)`$ for the donors. Note, that for the finite samples (especially for the relative small systems we were able to investigate) $`C_a(𝐑)`$, $`g_a(\epsilon _a,𝐑)`$ and $`g_d(\epsilon _d,𝐑)`$ depend essentially on the particular implementation R of the spatial distributions of the donor and acceptor sites (if a sample would be big enough all quantities would be self-averaging). Therefore, in order to obtain reliable results, one has to work with the quantities $`C_aC_a(𝐑)`$, $`g_a(\epsilon )g_a(\epsilon _a,𝐑)`$ and $`g_d(\epsilon )g_d(\epsilon _d,𝐑)`$, where $`\mathrm{}`$ denotes the average over a number of R’s. Note, that the values $`g_{a(d)}(\epsilon _{a(d)},𝐑)d\epsilon `$ obtained for independent R’s are scattered according to the Gaussian distribution with the mean $`g_{a(d)}(\epsilon )d\epsilon `$ and the standard deviation $`\sqrt{g_{a(d)}(\epsilon )d\epsilon }`$. In the region of the Coulomb gap $`g_{a(d)}(\epsilon )d\epsilon 10^4`$ and dispersion is several orders of magnitude larger than the mean. Therefore, in order to reduce the statistical noise in the final $`g_{a(d)}(\epsilon )`$ dependences an average is needed over a sufficient large amount of independent samples (we performed calculations with up to 10<sup>4</sup> samples). ### B Acceptor-donor symmetry Let us rewrite the energy (5) in terms of the OEE energies (9) $`E(n_a,n_d)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}\epsilon _a(i)n_a(i)`$ (17) $``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k}{}}\epsilon _d(k)(1n_d(k)){\displaystyle \frac{\mathrm{\Delta }}{2}}{\displaystyle \underset{i}{}}n_a(i).`$ (18) The system investigated is electrically neutral, i. e. for any sample $$\underset{i}{}n_a(i)=\underset{k}{}(1n_d(k)).$$ (19) Then, the energies of the microscopic states ($`n_a,n_d`$) and ($`n_a^{},n_d^{}`$) for a sample R and its “mirror” reflection R (when the donor and acceptor sites exchange places keeping the spatial arrangement of sites unchanged), are equal under the following conditions $$\epsilon _a(i)+\epsilon _d^{}(i)=\epsilon _d(k)+\epsilon _a^{}(k)=\mathrm{\Delta }$$ (20) and $$n_a^{}(i)=(1n_d(i))n_d^{}(k)=(1n_a(k)).$$ (21) The stability relations (12) for the ground-state ($`n_a^0`$,$`n_d^0`$) of the sample R transform into stability relations for the ground-state ($`n_a^0`$,$`n_d^0`$) of the sample R through the relations (20,21) as well. Since averaging over samples includes all possible pairs R and R, it follows from the symmetry relations (20) and (21) along with the definition (16) that $`g_a(\epsilon )`$ can be mapped to $`g_d(\epsilon )`$ using the relation $$g_d(\epsilon )=g_a(\epsilon \mathrm{\Delta })$$ (22) The symmetry of the model imposes also a relation between the Fermi energy $`\mu `$ (13,14) and the parameter $`\mathrm{\Delta }`$ of the model. Expressing $`n_{a[d]}(i[k])`$ in terms of the Heaviside’s step functions $`n_{a[d]}=\theta (\mu \epsilon _{a[d]}(i[k]))`$, the quantity $`C_a`$ can be written in the form $$C_a=_{\mathrm{}}^\mu g_a(\epsilon )𝑑\epsilon =_\mu ^{\mathrm{}}g_d(\epsilon )𝑑\epsilon $$ (23) The symmetry relation (22) transforms (23) into an integral relation $$_{\mu \mathrm{\Delta }}^{\mathrm{}}g_d(\epsilon )𝑑\epsilon =_\mu ^{\mathrm{}}g_d(\epsilon )𝑑\epsilon ,$$ (24) which has a meaning only if $$\mu =\frac{\mathrm{\Delta }}{2}.$$ (25) Thus, the Fermi energy of our model system in the thermodynamic limit is a fundamental quantity depending only on the energy of charge transfer from an acceptor to a donor. ## III METHOD ### A Algorithm of energy minimization We start from a random allocation of $`N`$ donor and $`N`$ acceptor sites in the continuous $`D`$-dimensional system (generate a sample R) with the density $`n=1`$, so that the system has a linear size $`L=N^{1/D}`$ and then charge randomly chosen $`C_a\times N`$ both donors and acceptors (usually we take $`C_a=0.7`$), i. e. generate an initial microscopic state (IMS) ($`n_a`$,$`n_d`$) of the sample R. Further, we search for such microscopic state ($`n_a^0`$,$`n_d^0`$) which obeys the stability conditions (12) with respect to the four mechanisms of the charge transfer allowed in our model. We used an algorithm which is an extension of the algorithm proposed in Ref. to the case $`\mathrm{\Delta }\mathrm{}`$. The algorithm consists of the three main steps. In order to save computer time, first, we look for pairs $`a^0a^{}`$ ($`d^0d^+`$) for which the “crude” stability relation $`\mathrm{\Delta }\epsilon \epsilon _{a(d)}^0\epsilon _{a(d)}^1>0`$ is violated. Then, the energy of the system is decreased by transferring an electron between such pair of sites for which $`\mathrm{\Delta }\epsilon `$ has its minimal non-positive value. This process is repeated until a state is reached, in which $`\mathrm{\Delta }\epsilon >0`$ for all possible $`a^0a^{}`$ and $`d^0d^+`$ pairs (step I). In the similar manner, we further minimize the energy of the system with respect to the “true” stability relations (12) for the charge transfer between the $`a^0a^{}`$ and $`d^0d^+`$ pairs (step II). And, finally, in the step III we diminish the energy of the system with respect to the stability relations for ionization and recombination processes. Since ionization and recombination processes change the degree $`C_a`$ of the system ionization, each time after one of these processes takes place during calculations, we go back to the step II. Repeating the steps II and III, we finally arrive at a microscopic state ($`n_a^0`$,$`n_d^0`$) for which all four stability conditions are fulfilled. We name the procedure ($`n_a`$,$`n_d`$)$``$ ($`n_a^0`$,$`n_d^0`$) via above steps I,II and III as “a single descent”. It should be noted, however, that the state ($`n_a^0`$,$`n_d^0`$) is not necessarily the ground state of the sample R since for the ground state, in general, not only the simplest relations (12) with only pairs of sites included, but the more complicated relations involving quadruplets, sextets, etc. of sites have to be fulfilled. Therefore, the state ($`n_a^0`$,$`n_d^0`$) (after Ref. ) hereafter will be referred to as the pseudo-ground state (PGS) of the sample R. Then, two questions naturally arise: How close the PGS and the ground state of the given sample are and how this may influence the output of our calculations? In order to answer the first question, we calculate and analyze the histograms $``$ for the so-called overlaps $$Q_{\alpha \beta }=\frac{1}{N}\underset{i}{}\delta (n_a^\alpha ,n_a^\beta ),$$ (26) where indices $`\alpha `$ and $`\beta `$ refer to PGS’s which are obtained by means of the single descent on the same sample but with different IMS ($`n_a`$,$`n_d`$). If two PGS’s are identical then $`Q_{\alpha \beta }=1`$. We calculated for the $`D=2`$ system with $`N=500`$ at $`\mathrm{\Delta }=0`$ the mean $`Q(𝐑)=Q_{\alpha \beta }_{\alpha \beta }`$ for the sequence of 100 PGS’s generated by single descents from the different IMS of the same sample R. We further acquire $`Q(𝐑)`$ for 100 different samples and obtain that the mean $`\overline{Q}Q(𝐑)_𝐑=0.96`$. It means that in PGS generated by the single descent only 20 acceptors out of 500 are, in average, in the “wrong” states compared to those in the true ground state of the sample. In order to evaluate how the “erroneousness” of PGS influences the outcome of our calculations we perform an analysis of ground states obtained by means of the so called multirank descents. Descent of rank $`m`$ comprises of a consequence of the single descents on the same sample with different IMS when calculations are stopped after the lowest observed PGS energy repeats $`m`$ times. We calculate $`\overline{Q}`$ (all other parameters were the same as described in the previous paragraph, where actually the case $`m=0`$ was explored) for descents with different ranks $`m=5,10,15`$ and found that, for instance, for $`m=15`$ (which implies drastic increase in the computation time) $`\overline{Q}=0.990`$. $`g_a(\epsilon )`$ and $`g_d(\epsilon )`$ obtained from the PGS’s generated by means of the single descents and by means of descents with $`m=10`$, say, do not differ within the limits of statistical errors. So, we conclude, that reliable results can be obtained by means of single descents already, thereby saving a lot of computer time and resources. ### B Finite-size effects Due to constraints in computer resources, the largest samples, we were able to deal with, comprise up to $`N=2000`$ donor and $`N=2000`$ acceptor sites ($`L45`$ for $`D=2`$ and $`L12`$ for $`D=3`$). Such relative small sizes of the samples investigated might influence the outcome of calculations. Detailed analysis of finite size effects on the results obtained will be presented in Section IV and here we want to make two remarks about inherent finite size effects in the model system considered. First, as follows from (12), the energies $`\epsilon _a^0`$ for the neutral acceptors and $`\epsilon _a^1`$ for the charged ones in finite samples at $`T=0`$ cannot be further away than $`(L\times \sqrt{D})^1`$. This implies that $`g(\epsilon _a)=0`$ within the $`\epsilon _a`$ interval $$\|\epsilon _a\mu \|<(2L\times \sqrt{D})^1$$ (27) Of course, the same holds for donors as well. The relation (27) gives the estimation how close to $`\mu `$ data on the energy spectrum are, in principle, obtainable from the calculations on finite samples. Secondly, as follows from (9) the energies $`\epsilon _a`$ and $`\epsilon _d`$ for the finite samples are sensitive to the location of the donor and acceptor sites. Therefore, the Fermi energy $`\mu `$ for finite samples does differ, in general, from sample to sample. A straightforward averaging of $`g(\epsilon )`$ over different samples might thus lead to a distortion of the $`g(\epsilon )`$ shape especially in the region where the Coulomb gap is observed. In order to avoid this undesired effect, we used a trick first proposed in Ref.. During accumulation of the results for $`g(\epsilon )`$ we added together $`g(\epsilon )`$ for the same values of $`\epsilon \mu (𝐑)`$ rather than for the same values of $`\epsilon `$. Here $`\mu (𝐑)`$ denotes the Fermi energy for a finite sample R calculated as $$\mu (𝐑)=\frac{1}{2}\left(\text{min}\{\epsilon _a^0(i)\}+\text{max}\{\epsilon _a^1(i)\}\right),$$ (28) Such way of doing $`g(\epsilon )`$ average entirely excludes the influence of the fluctuations of the Fermi energy in the finite samples on the shape of the Coulomb gap. Finally, we remark that all the data presented below were obtained for the open boundary condition. In order to ensure that results obtained are not determined by the type of the boundary conditions used in calculations, we performed calculations of the two-dimensional MCDAM at $`\mathrm{\Delta }=0`$ with different $`N`$ and found that periodic boundary conditions only effectively reduce the linear size of a sample, leaving the qualitative shape of the parameters calculated unchanged. ## IV RESULTS According to the symmetry relation (22) $`g_a(\epsilon )`$ and $`g_d(\epsilon )`$ can be easily mapped to each other for any values of $`\epsilon `$ and hence all the results presented below concern the acceptor sites only. One can expect that the width of the Coulomb gap $`\mathrm{\Delta }\epsilon `$ and the energy scale in our model $`E_0=e^2n^{1/D}/\chi `$ are of the same order of magnitude. Fig.1 shows $`g_a(\epsilon \mu )`$ in the vicinity of the Fermi energy $`\mu `$ obtained for the two-dimensional samples with $`N=1000`$ and various values of $`\mathrm{\Delta }`$. As it is seen, $`g_a(\epsilon \mu )`$ depends considerably on $`\mathrm{\Delta }`$ except for a narrow window $`\|\epsilon \mu \|0.05`$, where all data merge into some “universal” curve symmetric with respect to $`\mu `$, the curve which can be anticipated to obey the Efros universality hypothesis (1). However, a double-logarithmic plot of the “universal” $`g_a(\epsilon \mu )`$ (insert in the Fig.1), reveals that the behavior of $`g_a(\epsilon \mu )`$ in the “universality” region is not even a power law. The width of this “universality” region is comparable to the width of the region where $`g_a(\epsilon \mu )=0`$ due to the finite size effects (for the data presented in Fig.1 relation (27) gives $`\|\epsilon \mu \|<0.011`$), so it is plausible to suggest that the “universal” behavior of $`g_a(\epsilon \mu )`$ is governed by the finite-size effects. This is clearly demonstrated in Fig.2 where $`g_a(\epsilon \mu )`$ are shown for several sizes of the samples investigated. The $`\epsilon `$ window where finite size effects are severe, shrinks considerably with increasing $`N`$ for all values of $`\mathrm{\Delta }`$ we investigated. For instance, $`g_a(\epsilon \mu )`$ for $`N=500`$ and $`N=1000`$ at $`\mathrm{\Delta }=0`$ (see Fig.2a,c) merge when $`\|\epsilon \mu \|0.2`$ while corresponding curves for $`N=1000`$ and $`N=1500`$ are indistinguishable already at $`\|\epsilon \mu \|0.1`$. The statistical noise observed for the curves in Fig.2 is quite small even close to $`\mu `$ and hence, the influence of insufficient large statistics on the results obtained is excluded. Note, that the “universal” behavior of $`g(\epsilon )`$ in the vicinity of $`\mu `$ obtained for the classical $`da`$ model (see Fig.3 in Ref.) is most likely due to the finite size effects as well. In the region $`\|\epsilon \mu \|0.2`$, where the curves for all $`N`$ collapse into a single curve (and where we believe the thermodynamic limit is reached), the behavior of $`g_a(\epsilon \mu )`$ is described by a power law $`g_a(\epsilon \mu )\|\epsilon \mu \|^\gamma `$. The deviation from the power-law observed far away from $`\mu `$ ($`\|\epsilon \mu \|0.7`$) is due to the boundaries of the Coulomb gap which, as was mentioned above, are $`1`$ in units of $`E_0`$. One can see from a comparison of the data shown in Fig.2 for different $`\mathrm{\Delta }`$, that the exponent $`\gamma `$ depends considerably on $`\mathrm{\Delta }`$. Furthermore, values of $`\gamma `$ in the region $`\epsilon \mu >0`$ and those in the region $`\epsilon \mu <0`$ differ as well with this difference increasing with increasing $`\mathrm{\Delta }`$. The data for $`\gamma `$ obtained for the two-dimensional MCDAM are summarized in Fig.3 where a significant deviation of $`\gamma `$ from the value $`D1`$ predicted by the hypothesis (1) is observed at all values of $`\mathrm{\Delta }`$ investigated except for the case $`\mathrm{\Delta }=0`$ when $`\gamma 1`$ within the limits of statistical accuracy. Note, that the deviation of $`\gamma `$ from its predicted value grows monotonically with increasing $`\mathrm{\Delta }`$. At $`\mathrm{\Delta }=10`$ where the features of the MCDAM are expected to be nearly the same as those of the classical $`da`$ model with all the acceptors being ionized (indeed, the degree of the acceptor ionization $`C_a0.9`$ for the two-dimensional MCDAM at $`\mathrm{\Delta }=10`$, see Fig.6 below) the deviation from the Efros exponent is very large. The main results for $`g_a(\epsilon \mu )`$ obtained for the three-dimensional MCDAM are summarized in Figs. 4 and 5. It is seen, that the behavior of $`g_a(\epsilon \mu )`$ in three dimensions does not differ qualitatively from the behavior of $`g_a(\epsilon \mu )`$ in two dimensions. Some quantitative differences observed arise from the fact that at given $`N`$ (the parameter which determines the amount of computer memory needed for the calculations) the linear size of a two-dimensional sample with a given density of sites is larger than that of a three-dimensional sample with the same density of sites and thereby, the finite size effects for three-dimensional samples with given $`N`$ are more pronounced compared to those for the two-dimensional samples with the same $`N`$. For example, the lower boundary of the region where $`g_a(\epsilon \mu )`$ can be described by the power law $`\|\epsilon \mu \|^\gamma `$ shifts towards larger $`\|\epsilon \mu \|0.4`$ values (see inserts in Fig.4). Remarkably, the exponent $`\gamma `$ does not reach the value $`D1`$ predicted by the universality hypothesis (1) even at $`\mathrm{\Delta }=0`$ (Fig.5). Unlike $`g_a(\epsilon \mu )`$ in the vicinity of the Coulomb gap, the density of ionized acceptors $`C_a`$ (15) describes the state of the entire sample and therefore reaches the thermodynamic limit much faster than $`g_a(\epsilon \mu )`$. This allows us to obtain quite accurate results for $`C_a`$ from data on a relatively small amount of samples with $`N=500`$ only. Fig.6 shows the variations of $`C_a`$ with $`\mathrm{\Delta }`$ both for two and three dimensions. In three dimensions almost all acceptors become ionized ($`C_a1`$) rather soon while for two dimensions even for the largest $`\mathrm{\Delta }`$ investigated around 10 % of the acceptors remain neutral. So, one can say, that the three-dimensional MCDAM at $`\mathrm{\Delta }7`$ reduces already to the classical $`da`$ model. It is known that the classical $`da`$ model exhibits in three dimensions the, so called, Coulomb fluctuational catastrophe . For calculations on finite samples it implies that statistical fluctuations of $`\mu (𝐑)`$ grow dramatically with increasing $`\mathrm{\Delta }`$ which is the case in our calculations (see Table I). Therefore, in order to reduce the statistical noise in three dimensions, the average of $`g_a(\epsilon \mu )`$ over a much larger (compared to $`D=2`$) number of samples is needed. Note, that $`\mu (𝐑)`$ in both two and three dimensions are scattered according to the Gaussian distribution with the mean $`\overline{\mu }`$ obeying the relation (25). ## V DISCUSSION The behavior of $`g_a(\epsilon \mu )`$ calculated within the region of the Coulomb gap for the model (5) is in strong contradiction to the universality hypothesis (1). Despite the fact that $`g_a(\epsilon \mu )`$ is indeed described by the power law $`\|\epsilon \mu \|^\gamma `$ in a wide range of $`\epsilon `$ inside the region of the Coulomb gap, the exponent $`\gamma `$ is considerably smaller than that predicted by the hypothesis (1) both for the two- and three-dimensional cases. Moreover, the exponent $`\gamma `$ depends significantly on $`\mathrm{\Delta }`$ and is different for the cases $`\epsilon >\mu `$ and $`\epsilon <\mu `$. It is believed that information about $`g(\epsilon )`$ might be directly obtained from tunneling and photoemission experiments and recent experiments on boron-doped silicon crystals have shown that the density of one-electron excitations at higher energies obeys a power-law with an exponent slightly less than $`0.5`$ which is in good agreement with our results for $`D=3`$ and $`\mathrm{\Delta }8`$. However, the non-metallic samples show around the Fermi energy a nearly quadratic Coulomb gap, so the question arises whether our results could be related to the intermediate asymptotic behavior observed? Here we want to make three remarks concerning this question: First, the power law $`g_a(\epsilon \mu )\|\epsilon \mu \|^\gamma `$ is valid above a value $`\epsilon _0(N)`$ below which the finite size effects take over (Figs.2 and 4). It seems from our results, that $`\epsilon _0(N)\mu `$ when $`N\mathrm{}`$. In two dimensions we were able to obtain size-independent results down to $`\epsilon _00.1`$, i. e. for $`90\%`$ of the whole Coulomb gap, the halfwidth of which is $`1`$ in units of $`E_0`$. Secondly, as follows from the ground-state stability relations (2), the distance $`r_{ij}`$ between a neutral donor, with an energy, say, $`\epsilon _i^1[\epsilon ,0]`$ ($`\epsilon `$ here is the halfwidth of a narrow band around $`\mu =0`$) and a charged donor with an energy $`\epsilon _j^0[0,\epsilon ]`$ should be not less that $`\frac{1}{2\epsilon }`$. I. e., sites with energies $`\epsilon _i^1[\epsilon ,0]`$ cannot be inside a D-dimensional sphere of radius $`R_{sp}=\frac{1}{2\epsilon }`$ and with the center in a site with the energy $`\epsilon _j^0[0,\epsilon ]`$. Assuming that all such spheres do not intersect, the total volume occupied by the spheres is $$V_{sp}=N\times S(D)\left(\frac{1}{2\epsilon }\right)^D_0^\epsilon g(\epsilon ^{})𝑑\epsilon ^{}$$ (29) where $`S(D)`$ is the volume of a D-dimensional sphere with the radius equal to unity. Since $`V_{sp}`$ cannot exceed the total volume $`V`$ of a sample ($`V=N`$ at $`n=1`$) we arrive at the inequality $$_0^\epsilon g(\epsilon ^{})𝑑\epsilon ^{}\frac{(2\epsilon )^D}{S(D)},$$ (30) which is valid for all $`\epsilon `$ if $$g(\epsilon )\frac{D\times 2^D}{S(D)}\|\epsilon \|^{D1}$$ (31) The universality hypothesis (1) then is a limit case of (31). The density of sites with energies $`\epsilon _i^1[\epsilon ,0]`$ indeed decreases when $`\epsilon 0`$, so the assumption (29) for the spheres with finite radii seems to be plausible. However, simultaneously $`R_{sp}\mathrm{}`$ and consequently the plausibility of the assumption (29) and thereby of the hypothesis (1) becomes questionable. And finally, the universality hypothesis (1) can be also obtained as the asymptotic behavior of a non-linear integral equation for $`g(\epsilon )`$ as $`\epsilon 0`$, the equation which, in turn, is heuristically obtained from the stability condition (2). The derivation of this integral equation (given, for example, in Ref.) is based on the implicit assumption that the sites with charged donors are randomly distributed in space according to the Poisson statistics. However, it was unequivocally demonstrated in computer studies of the Coulomb gap that charged donor sites with energies close to $`\mu `$ tend to form clusters (Ref., Fig. 6). We conclude that $`g_a(\epsilon \mu )`$ in the region of the Coulomb gap in model (5) has a power law behavior for all energies down to $`\mu `$ and that the universality hypothesis of Efros (1) is questionable. Note, that our results are in contradiction not only to the universality hypothesis (1), but to the inequality (31) as well. Up to now, all exponents found are in good agreement with this inequality. E. g. in Ref. specimens of 40 000 and 125 000 sites for two- and three-dimensional samples were investigated in the Efros’ lattice model and the power law $`g_a(\epsilon \mu )\|\epsilon \mu \|^\gamma `$ was found with $`\gamma =1.2\pm 0.1`$ and $`\gamma =2.6\pm 0.2`$ for two and three dimensions, respectively. The main conclusion of our results and those of Ref. is that Efros’ lattice model can not be used to as a reliable approximation to the classical $`da`$ model with Poisson impurity distribution. Energy levels of donor (acceptor) impurities are usually close to the bottom (top) of the conduction (valence) band. Since in the most common semiconductors the energy gap $`E_g10^4`$ K and $`E_020`$ K, $`\mathrm{\Delta }1`$ and one may ask what physical relevance does the model (5) with a finite $`\mathrm{\Delta }10`$ have, except for being a pure academic exercise? However, in the case of deep impurities the energy levels for some donor–acceptor pairs are extremely close to each other not excluding even the case $`\mathrm{\Delta }=0`$. Table II shows some donor–acceptor pairs with $`\mathrm{\Delta }10`$ in the most common semiconductors. The solubilities of these impurities are rather low, thereby reducing the temperature at which the Coulomb gap with features described by the model (5) can be observed. Fortunately, these temperatures are high enough ($`10÷20`$ K) for modern experimental techniques and hence experimental observation of the Coulomb gap in the semiconductors with deep impurities is possible to accomplish. ## VI SUMMARY We have studied a model of impurities in semiconductors with infinite-range Coulomb interactions between donors, between acceptors and between donors and acceptors. A new parameter introduced in the model is the finite energy $`\mathrm{\Delta }`$ of charge transfer between donors and acceptors, a parameter which enables processes of ionization of neutral impurities and of recombination of charged impurities. In the particular case of equal amounts of donor and acceptor impurities, we derived rigorous relations for the symmetry of the model with respect to exchange of donor and acceptor sites. We also extended the previously known algorithm to find the ground state including the stability relations with respect to ionization and recombination processes and performed computer studies of the model proposed at zero temperature on a number of two- and three-dimensional samples with randomly distributed $`N`$ donors and $`N`$ acceptors. We explored the energy region around the Fermi energy $`\mu `$ where the Coulomb gap in the density of one-electron excitations $`g(\epsilon )`$ is observed. The analysis of the calculated histograms $`g(\epsilon )`$ revealed that the behavior of $`g(\epsilon )`$ obtained from the simulations on finite samples in the immediate neighborhood of $`\mu `$ is determined solely by the finite size effects. In the region where finite size effects become negligible $`g(\epsilon )`$ is described by a power law with an exponent considerably depending on the parameter $`\mathrm{\Delta }`$ and on the sign of $`\epsilon \mu `$. Our findings challenge the Efros universality hypothesis. Moreover, our results are in contradiction to the main inequality (31) of which Efros’ universality hypothesis is a particular case. We have reexamined the heuristic derivation of the Efros hypothesis and shown that some implicit assumptions which lead to universality are questionable. From the analysis of experimental data on admixtures in semiconductors we put forward possible experimental situations where one could observe the Coulomb gap with the features being the same as those of the model with a finite $`\mathrm{\Delta }`$. ## Acknowledgements This research was supported by The Swedish Natural Science Council and by The Swedish Royal Academy of Sciences.
warning/0003/hep-th0003061.html	ar5iv	text	# Black holes on the brane ## I Introduction Recent developments in string theory have shown that if matter fields are localized on a 3-brane in $`1+3+d`$ dimensions, while gravity can propagate in the $`d`$ extra dimensions, then the extra dimensions can be large (see, e.g., ). The extra dimensions need not even be compact, as in the 5-dimensional warped space models of Randall and Sundrum . (See also for earlier work.) In particular, they showed that it is possible to localize gravity on a 3-brane when there is one infinite extra dimension. If matter on a 3-brane collapses under gravity, without rotating, to form a black hole, then the metric on the brane-world should be close to the Schwarzschild metric at astrophysical scales in order to preserve the observationally tested predictions of general relativity. Collapse to a black hole in the Randall-Sundrum brane-world scenario was studied by Chamblin et al. (see also ). They gave a ‘black string’ solution which intersects the brane in a Schwarzschild solution. Here we give an exact localized black hole solution, which remarkably has the mathematical form of the Reissner-Nördstrom solution, but without electric charge being present. Instead the Reissner-Nördstrom-type correction to the Schwarzschild potential can be thought of as a ‘tidal charge’, arising from the projection onto the brane of free gravitational field effects in the bulk. These effects are transmitted via the bulk Weyl tensor (see below). The Schwarzschild potential $`\mathrm{\Phi }=M/(M_\mathrm{p}^2r)`$, where $`M_\mathrm{p}`$ is the effective Planck mass on the brane, is modified to $$\mathrm{\Phi }=\frac{M}{M_\mathrm{p}^2r}+\frac{Q}{2r^2},$$ (1) where the constant $`Q`$ is a ‘tidal charge’ parameter, which may be positive or negative. A geometric approach to the Randall-Sundrum scenario has been developed by Shiromizu et al. (see also ), and proves to be a useful starting point for formulating the problem and seeing clear lines of approach. The field equations in the bulk are (modifying the notation of ) $$\stackrel{~}{G}_{AB}=\stackrel{~}{\kappa }^2\left[\stackrel{~}{\mathrm{\Lambda }}\stackrel{~}{g}_{AB}+\delta (\chi )\left(\lambda g_{AB}+T_{AB}\right)\right],$$ (2) where the tildes denote bulk quantities. The fundamental 5-dimensional Planck mass $`\stackrel{~}{M}_\mathrm{p}`$ enters via $`\stackrel{~}{\kappa }^2=8\pi /\stackrel{~}{M}_\mathrm{p}^3`$. The brane tension is $`\lambda `$, and $`\stackrel{~}{\mathrm{\Lambda }}`$ is the bulk cosmological constant. The brane is located at $`\chi =0`$ (so that $`x^4=\chi `$ is a natural choice for the fifth dimension coordinate), and $`g_{AB}=\stackrel{~}{g}_{AB}n_An_B`$ is the induced metric on the brane, with $`n_A`$ the spacelike unit normal to the brane. The brane energy-momentum tensor is $`T_{AB}`$, and $`T_{AB}n^B=0`$. The brane is a fixed point of the $`Z_2`$ symmetry. ## II Field equations on the brane The field equations induced on the brane arise from Eq. (2), the Gauss-Codazzi equations and the matching conditions with $`Z_2`$-symmetry, and they may be written as a modification of the standard Einstein equations, with the new terms carrying bulk effects onto the brane : $$G_{\mu \nu }=\mathrm{\Lambda }g_{\mu \nu }+\kappa ^2T_{\mu \nu }+\stackrel{~}{\kappa }^4S_{\mu \nu }_{\mu \nu },$$ (3) where $`\kappa ^2=8\pi /M_\mathrm{p}^2`$. The energy scales are related to each other and to the cosmological constants via $$M_\mathrm{p}=\sqrt{\frac{3}{4\pi }}\left(\frac{\stackrel{~}{M}_\mathrm{p}^2}{\sqrt{\lambda }}\right)\stackrel{~}{M}_\mathrm{p},\mathrm{\Lambda }=\frac{4\pi }{\stackrel{~}{M}_\mathrm{p}^3}\left[\stackrel{~}{\mathrm{\Lambda }}+\left(\frac{4\pi }{3\stackrel{~}{M}_\mathrm{p}^3}\right)\lambda ^2\right].$$ (4) Typically, the fundamental Planck scale is much lower than the effective scale in the brane-world: $`\stackrel{~}{M}_\mathrm{p}M_\mathrm{p}`$. Local bulk effects on the matter are transmitted via the ‘squared energy-momentum’ tensor $`S_{\mu \nu }`$, but since we will consider vacuum solutions, the precise form of $`S_{\mu \nu }`$ (see ) will not be needed. In vacuum, $`T_{\mu \nu }=0=S_{\mu \nu }`$, and we also choose the bulk cosmological constant to satisfy $`\stackrel{~}{\mathrm{\Lambda }}=4\pi \lambda ^2/3\stackrel{~}{M}_\mathrm{p}^3`$, so that $`\mathrm{\Lambda }=0`$ by Eq. (4). Then Eq. (3) reduces to $$R_{\mu \nu }=_{\mu \nu },R_\mu {}_{}{}^{\mu }=0=_\mu {}_{}{}^{\mu },$$ (5) where $`_{\mu \nu }`$ is the limit on the brane of the projected bulk Weyl tensor : $$_{AB}=\stackrel{~}{C}_{ACBD}n^Cn^D.$$ (6) The Weyl symmetries ensure that this is symmetric and tracefree ($`_{[AB]}=0=_A^A`$) and has no orthogonal components ($`_{AB}n^B=0`$, so that $`_{AB}_{\mu \nu }\delta _A{}_{}{}^{\mu }\delta _{B}^{}^\nu `$ as $`\chi 0`$). It carries the influence of nonlocal gravitational degrees of freedom in the bulk onto the brane, including the tidal (or Coulomb) and transverse traceless (gravitational wave) aspects of the free gravitational field. (See for a fuller discussion of $`_{\mu \nu }`$ in the general case.) On the brane, in the vacuum case, this tensor satisfies the divergence constraint $$^\mu _{\mu \nu }=0,$$ (7) where $`_\mu `$ is the brane covariant derivative. In view of the Bianchi identities on the brane, this is an integrability condition for the field equation $`R_{\mu \nu }=_{\mu \nu }`$. For static solutions, Eqs. (5) and (7) form a closed system of equations on the brane . A vacuum solution outside a mass localized on the brane must satisfy equations (5) and (7). This leads to a prescription for mapping 4-dimensional general relativity solutions to brane-world solutions in 5-dimensional gravity: a stationary general relativity solution with tracefree energy-momentum tensor gives rise to a vacuum brane-world solution in 5-dimensional gravity. The 4-dimensional general relativity energy-momentum tensor $`T_{\mu \nu }`$ (where $`T_\mu {}_{}{}^{\mu }=0`$) is formally identified with the bulk Weyl term on the brane via the correspondence $`\kappa ^2T_{\mu \nu }_{\mu \nu }.`$ The general relativity conservation equations $`^\nu T_{\mu \nu }=0`$ correspond to the constraint equation (7) on the brane. In particular, Einstein-Maxwell solutions in general relativity will lead to vacuum brane-world solutions. This is the observation that led us to the Reissner-Nördstrom-type solution. ## III Solutions with tidal charge Algebraic symmetry properties imply that in general we can decompose $`_{\mu \nu }`$ irreducibly with respect to a chosen 4-velocity field $`u^\mu `$ as $$_{\mu \nu }=\left(\frac{\stackrel{~}{\kappa }}{\kappa }\right)^4\left[𝒰\left(u_\mu u_\nu +\frac{1}{3}h_{\mu \nu }\right)+𝒫_{\mu \nu }+2𝒬_{(\mu }u_{\nu )}\right],$$ (8) where $`h_{\mu \nu }=g_{\mu \nu }+u_\mu u_\nu `$ projects orthogonal to $`u^\mu `$. Here $`𝒰=\left({\displaystyle \frac{\kappa }{\stackrel{~}{\kappa }}}\right)^4_{\mu \nu }u^\mu u^\nu `$ is an effective energy density on the brane arising from the free gravitational field in the bulk—but note that this energy density need not be positive. Indeed, as we argue below, $`𝒰<0`$ is the natural case. The effective anisotropic stress from the free gravitational field in the bulk is the spatially tracefree and symmetric part, i.e., $`𝒫_{\mu \nu }=\left({\displaystyle \frac{\kappa }{\stackrel{~}{\kappa }}}\right)^4\left[h_{(\mu }{}_{}{}^{\alpha }h_{\nu )}^{}{}_{}{}^{\beta }\frac{1}{3}h_{\mu \nu }h^{\alpha \beta }\right]_{\alpha \beta },`$ where round brackets denote symmetrization. The effective energy flux from the free gravitational field in the bulk is $`𝒬_\mu =\left({\displaystyle \frac{\kappa }{\stackrel{~}{\kappa }}}\right)^4h_\mu {}_{}{}^{\alpha }_{\alpha \beta }^{}u^\beta .`$ In a static vacuum, with $`u^\mu `$ along the static Killing direction, we have $`𝒬_\mu =0`$, and the effective conservation equation (7) reduces to the single spatial equation $$\frac{1}{3}\mathrm{D}_\mu 𝒰+\frac{4}{3}𝒰A_\mu +\mathrm{D}^\nu 𝒫_{\mu \nu }+A^\nu 𝒫_{\mu \nu }=0,$$ (9) where $`\mathrm{D}_\mu `$ is the projection (orthogonal to $`u^\mu `$) of the brane covariant derivative , and $`A_\mu =u^\nu _\nu u_\mu `$ is the 4-acceleration. Static spherical symmetry means that $`A_\mu =A(r)r_\mu ,𝒫_{\mu \nu }=𝒫(r)\left[r_\mu r_\nu \frac{1}{3}h_{\mu \nu }\right],`$ for some functions $`A(r)`$ and $`𝒫(r)`$, where $`r`$ is the areal distance and $`r_\mu `$ is a unit radial vector. The Reissner-Nördstrom-type solution in Eq. (1) corresponds to the solution $$𝒰=\left(\frac{\kappa }{\stackrel{~}{\kappa }}\right)^4\frac{Q}{r^4}=\frac{1}{2}𝒫$$ (10) of Eq. (9). We can verify that the solution in Eqs. (1) and (10) satisfies the field equations (5), using natural coordinates, for which the metric on the brane is $$ds^2=A(r)dt^2+B(r)dr^2+r^2(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2).$$ (11) Then it may be verified that $`A=B^1=1+{\displaystyle \frac{\alpha }{r}}+{\displaystyle \frac{\beta }{r^2}},`$ (12) $`_t{}_{}{}^{t}=_r{}_{}{}^{r}=_\theta {}_{}{}^{\theta }=_\phi {}_{}{}^{\phi }={\displaystyle \frac{\beta }{r^4}},`$ (13) where $`\alpha `$ and $`\beta `$ are constants. Equations (12) and (13) satisfy all the field equations (5), and hence also the divergence equations (7). By Eqs. (8) and (10), we see that $`\beta =Q`$, and the far-field Newtonian limit \[see Eq. (1)\] shows that $`\alpha =2M/M_\mathrm{p}^2`$. In summary, we have shown that an exact black hole solution of the effective field equations on the brane is given by the induced metric $$g_{tt}=(g_{rr})^1=1\left(\frac{2M}{M_\mathrm{p}^2}\right)\frac{1}{r}+\left(\frac{q}{\stackrel{~}{M}_\mathrm{p}^2}\right)\frac{1}{r^2},$$ (14) where $`q=Q\stackrel{~}{M}_\mathrm{p}^2`$ is a dimensionless tidal charge parameter. The projected Weyl tensor, transmitting the tidal charge stresses from the bulk to the brane, is $$_{\mu \nu }=\left(\frac{q}{\stackrel{~}{M}_\mathrm{p}^2}\right)\frac{1}{r^4}\left[u_\mu u_\nu 2r_\mu r_\nu +h_{\mu \nu }\right].$$ (15) ## IV Properties of the black hole The 4-dimensional horizon structure of the brane-world black hole depends on the sign of $`q`$. For $`q0`$, there is a direct analogy to the Reissner-Nördstrom solution, with two horizons: $$r_\pm =\frac{M}{M_\mathrm{p}^2}\left[1\pm \sqrt{1q\frac{M_\mathrm{p}^4}{M^2\stackrel{~}{M}_\mathrm{p}^2}}\right].$$ (16) As in general relativity, both horizons lie inside the Schwarzschild horizon: $`0r_{}r_+r_\mathrm{s}=2M/M_\mathrm{p}^2`$, and there is an upper limit on $`q`$: $$0qq_{\mathrm{max}}=\left(\frac{\stackrel{~}{M}_\mathrm{p}}{M_\mathrm{p}}\right)\left(\frac{M}{M_\mathrm{p}}\right)^2.$$ (17) The intriguing new possibility that $`q<0`$, which is impossible in the general relativity Reissner-Nördstrom case, leads to only one horizon, lying outside the Schwarzschild horizon: $$r_+=\frac{M}{M_\mathrm{p}^2}\left[1+\sqrt{1q\frac{M_\mathrm{p}^4}{M^2\stackrel{~}{M}_\mathrm{p}^2}}\right]>r_\mathrm{s}.$$ (18) In the $`q<0`$ case, the (single) horizon has a greater area than its Schwarzschild counterpart, so that bulk effects act to increase the entropy and decrease the temperature of the black hole. In general relativity, the electric field in the Reissner-Nördstrom solutions acts to weaken the gravitational field, and the same is true for the brane-world black hole with $`q>0`$. This can be seen clearly from Eq. (1). By contrast, the $`q<0`$ case corresponds to the opposite effect, i.e., bulk effects tend to strengthen the gravitational field. By Eq. (10), we see that in this case, the effective energy density $`𝒰`$ on the brane contributed by the free gravitational field in the bulk is negative. This is in accord with the (Newtonian) notion that the gravitational field of an isolated mass has negative energy density. Furthermore, it agrees with the perturbative analysis by Sasaki et al. and the nonperturbative analysis of Maartens . The gravitational field generated by a source on the brane tends to squeeze test matter in the 5th direction, thus acting as an attractive field with a negative energy contribution. The tidal acceleration measured by static observers along the 5th direction $`n^A`$ is $`\stackrel{~}{R}_{ABCD}u^An^Bu^Cn^D=\left({\displaystyle \frac{\stackrel{~}{\kappa }}{\kappa }}\right)^4𝒰+\frac{1}{6}\stackrel{~}{\kappa }^2\stackrel{~}{\mathrm{\Lambda }},`$ where the right hand side follows from Eqs. (2) and (6). The negative bulk cosmological constant contributes to acceleration towards the brane, reflecting its confining role on the gravitational field. In order for $`𝒰`$ to reinforce confinement, it must be negative. In other words, negative tidal charge $`q<0`$ is the physically more natural case. (See also for further discussion of negative energy density from bulk effects.) Furthermore, $`q<0`$ ensures that the singularity is spacelike, as in the Schwarzschild solution, whereas $`q>0`$ leads to a timelike singularity, which amounts to a qualitative change in the nature of the general relativistic Schwarzschild solution. It is widely assumed that astrophysical black holes cannot exhibit macroscopic electric charge due to the presence of neutralizing plasma in their vicinity. Discussions of astrophysical black hole phenomena are therefore commonly restricted to the Kerr family of black holes. The solution presented here raises the possibility that an effective Reissner-Nördstrom metric emerges by geometric considerations and is hence not constrained by the above argument. The tidal charge $`q`$ affects the geodesics and the gravitational potential, so that indirect limits may be placed on it by observations. Current observational limits on $`\|q\|`$ are rather weak, since the correction term in Eq. (1) dies off rapidly with increasing $`r`$, and astrophysical measurements (lensing and perihelion precession) probe mostly (weak-field) solar scales. These measurements require the correction term in the gravitational potential to be much less than the Schwarzschild term, so that $$\|q\|2\left(\frac{\stackrel{~}{M}_\mathrm{p}}{M_\mathrm{p}}\right)^2M_{}R_{}.$$ (19) This still allows for large values of $`\|q\|`$, which would modify the spacetime geometry of a nonrotating black hole in the strong-gravity regime , with implications for example for the last stable circular orbit for compact binaries. Although the strong-gravity regime is not currently directly accessible to observations, indirect limits could emerge from the way in which tidal charge modifies general relativity strong-gravity effects. This deserves further investigation. Further indirect limits could arise from the effect of tidal charge on primordial black holes, which also merits further analysis. An intriguing possibility is that strong tidal effects in the early universe could lead to black hole formation even in the absence of gravitational collapse of matter. If such matter-free tidal collapse is possible, then the endstate is a metric with $`M=0`$ and $`q<0`$: $$g_{tt}=(g_{rr})^1=1+\left(\frac{q}{\stackrel{~}{M}_\mathrm{p}^2}\right)\frac{1}{r^2},$$ (20) and the horizon on the brane is given by $$r_\mathrm{h}=\frac{\sqrt{q}}{\stackrel{~}{M}_\mathrm{p}}.$$ (21) The tidal charge parameter arises from purely scalar (Coulomb-like) effects of the free gravitational field in the bulk, given the static spherical symmetry. In the original Randall-Sundrum model, as well as in perturbative studies of it (see, e.g., ), the bulk is assumed to be exactly anti-de Sitter in the absence of any source on the brane. The geometric approach of , which we have adopted here, makes no assumptions about the bulk metric, other than that it satisfies the 5-dimensional Einstein equations with cosmological constant. Thus the bulk need not be conformally flat in the absence of sources on the brane. This means that in general, the tidal charge parameter $`q`$ will be determined by both the brane source, i.e., the mass $`M`$, and any Coulomb part of the bulk Weyl tensor that survives when $`M`$ is set to zero. ## V Conclusion We have not investigated fully the effect of the brane-world black hole on the bulk geometry, and in particular the nature of the off-brane horizon structure. This has been done for solutions which reduce to the Schwarzschild black hole on the brane . In these solutions, the bulk metric is $`d\stackrel{~}{s}^2=\left({\displaystyle \frac{6}{\stackrel{~}{\kappa }^2\stackrel{~}{\mathrm{\Lambda }}}}\right){\displaystyle \frac{1}{z^2}}\left[dz^2+g_{\mu \nu }(x^\alpha )dx^\mu dx^\nu \right],`$ where $`g_{\mu \nu }`$ is the Schwarzschild metric. We have adopted a different approach: instead of starting from an induced Schwarzschild metric, we have solved the effective field equations for the induced metric on the brane (which form a closed system, since the metric is static), and found a generalization of the Schwarzschild solution. It turns out that if $`ds^2`$ is given by our solution, as in Eqs. (11)–(14), then $`d\stackrel{~}{s}^2=f(z)\left[dz^2+ds^2\right]`$ cannot satisfy the bulk field equations for any $`f(z)`$ if $`q0`$. Finding an exact form of the bulk metric that is consistent with our exact induced metric on the brane is more difficult than the case where the induced metric is Schwarzschild. Perturbative studies, which start from an exactly anti-de Sitter background, show that the first weak-field correction of the Newtonian potential on the brane is proportional to $`1/r^3`$ (see ). Our solution has by contrast a $`1/r^2`$ correction, so that it is incompatible with the long-distance limit on the brane. (Such a correction can also arise in thick-brane models .) However, in the short-distance limit on the brane, the lowest order correction to the potential is proportional to $`1/r^2`$ . This term will dominate the $`1/r`$ term, which reflects the fact that gravity becomes effectively 5-dimensional at high energies. Thus our solution should describe well the strong-gravity regime on the brane. In the short-distance limit, the perturbative analysis shows that $`q={\displaystyle \frac{M}{\stackrel{~}{M}_\mathrm{p}}},`$ so that the tidal charge is negative, as we argued above, and is determined by the black hole mass (which is to be expected if the background bulk Weyl tensor vanishes). In order to pursue our nonperturbative analysis, we need to look at the off-brane equations for the curvature (see for the general form of these equations). The induced field equations on the brane given in Eq. (5) are supplemented by off-brane equations which determine $`_𝐧_{AB},_𝐧_{ABC},_𝐧R_{ABCD},`$ where $`_𝐧`$ is the Lie derivative along $`n^A`$, $`R_{ABCD}`$ is the 4-dimensional Riemann tensor, and $`_{ABC}=g_A{}_{}{}^{D}g_{B}^{}{}_{}{}^{E}\stackrel{~}{C}_{DECF}^{}n^F,`$ with $`_{\mu \nu \sigma }=0`$ on the brane. These equations together with Eq. (5) form a closed system . In conclusion, we have shown how a Reissner-Nördstrom-type metric satisfies the effective field equations on the brane in Randall-Sundrum-type gravity, which form a closed system because of staticity. There is no electric charge present, but instead a tidal charge, arising from the imprint of the free gravitational field in the bulk on the brane. The tidal charge correction to the Schwarzschild potential is negative, and the solution describes the strong-gravity regime on the brane. Negative $`q`$ leads to an horizon on the brane that is outside the Schwarzschild horizon, corresponding to lower temperature and greater entropy. Current observations place only weak limits on the tidal charge, and in principle significant tidal modifications could arise in the strong-field regime or in the early-universe case of primordial black holes. Further investigation is in progress to probe these modifications and any indirect limits that they may impose on the tidal charge, as well as to find the off-brane behaviour of the horizon. Acknowledgements: VR was supported by a Royal Society grant while at Portsmouth, and thanks the Relativity and Cosmology Group for hospitality. RM thanks IUCAA, Pune for hospitality during a visit, which was partially supported by the Royal Society. We thank Bruce Bassett, Marco Bruni, Roberto Casadio, Roberto Emparan, David Langlois, Jose Senovilla, Carlo Ungarelli and David Wands for useful discussions and comments, and especially Tetsuya Shiromizu.
warning/0003/astro-ph0003336.html	ar5iv	text	# Galactic Evolution of Beryllium and Oxygen ## 1. Introduction First attempts to measure beryllium abundances in metal-poor stars by Molaro & Beckman (1984) and Molaro, Beckman & Castelli (1984) demonstrated that stars in the early Galaxy formed with much lower Be abundances than in the present epoch. First detection of Be in metal-poor stars was achieved by Rebolo et al. (1988) and further studies by S. Ryan, G. Gilmore, A. Boesgaard, P. Molaro, R. J. García López and their respective collaborators revealed a clear linear correlation with iron. Accelerated protons and $`\alpha `$-particles in cosmic rays interact with ambient CNO in ISM and create Be. According to the standard Galactic Cosmic Ray (GCR) theory, these interactions in the general ISM should have given a quadratic relation between Be and O. Alternatively, spallation of cosmic ray CNO nuclei accelerated out of freshly processed material could account for the primary character of the observed early galactic evolution of Be. Another production site is the collective acceleration by SN shocks of ejecta-enriched matter in the interiors of superbubbles. In these two cases, the evolution of Be should reflect the production of CNO from massive stars. Oxygen is mostly produced by Type II SNe while iron is produced in both, Type II and in Type Ia SNe. The fact that Type Ia SNe have longer lifetime progenitors has been commonly used to argue that oxygen must be overabundant in very old stars. Observational evidence for high \[O/Fe\] ratios in many metal-poor stars has been reported over the last two decades. Based on the study of \[O i\] lines at 6300 and 6363 Å in evolved stars (though the second line at 6363 Å is not visible in very metal-poor stars and the analysis is based only on one line), several authors have found that \[O/Fe\]$`=0.30.4`$ dex at \[Fe/H\]$`<1`$ and is constant until \[Fe/H\]$`3`$ (e.g. Barbuy 1988 and Kraft et al. 1992). In contrast with this result, oxygen abundances derived in unevolved stars using the O i IR triplet at 7774 Å (Abia & Rebolo 1989; Tomkin et al. 1992; King & Boesgaard 1995; and Cavallo, Pilachowski, & Rebolo 1997) point towards linearly increasing \[O/Fe\] values with decreasing \[Fe/H\] and reaching a ratio $`1`$ for stars with \[Fe/H\]$`3`$. This may suggest a higher production of oxygen during the early Galaxy. We discuss in this paper the comparison of these abundances with those derived from OH lines located in the near-UV part of the spectra of metal-poor stars, and their relation with beryllium abundances consistently derived from the same spectra. ## 2. Observations and Analyses The observations were carried out in different runs using the UES ($`R=\lambda /\mathrm{\Delta }\lambda 50000`$) of the 4.2-m WHT at the Observatorio del Roque de los Muchachos (La Palma), and the UCLES ($`R60000`$) of the 3.9-m AAT. The spectral region observed spanned typically from 3080 to 3300 Å, where several OH lines and the Be ii doublet at 3131 Å are located. Details of the abundance analyses can be found in García López, Severino, & Gomez (1995) and Israelian, García López, & Rebolo (1998) for beryllium and oxygen, respectively. The stellar parameters play an important role in the abundance determinations from near-UV lines. Effective temperatures ($`T_{\mathrm{eff}}`$) for our stars were estimated using the Alonso et al. (1996) calibrations versus $`VK`$ and $`by`$ colors, which were derived by applying the infrared flux method (IRFM), and cover a wide range of spectral types and metal content. These temperatures were used to compute synthetic spectra around the Be ii doublet and slightly modified, within the error bars provided by the calibrations, until obtaining a good reproduction of this region in the observed spectra. Metallicities were adopted from literature values obtained from high resolution spectra. Adopted gravities, derived using the accurate parallaxes measured by Hipparcos (ESA 1997), are larger by 0.28 dex in average than the values adopted by Israelian et al. (1998). This implies a mean small reduction of 0.09 dex in the oxygen abundances inferred from the OH lines with respect to that work, which does not affect significantly their original results. Beryllium abundances reported here were obtained from the Be ii resonance doublet located at 3130.421 and 3131.065 Å. The first line, which is also the strongest one, is severely blended with atomic and molecular lines of other species, and the abundance determination usually relies only on the other line, more isolated and weaker. The abundances used in this work are those presented in García López (1999). ## 3. Oxygen Israelian et al. (1998) presented new oxygen abundances derived from near-UV OH lines (which form in the same layers of the atmosphere as \[O i\]) for 24 metal-poor stars. They have concluded that the \[O/Fe\] ratio of metal-poor stars increases from 0.6 to 1 between \[Fe/H\]=$``$1.5 and $`3`$, with a slope of $`0.31\pm 0.11`$. Contrary to the previously accepted picture (see e.g. Bessell, Sutherland, & Ruan 1991, who used older model atmospheres with a coarser treatment of the opacities in the UV), these new oxygen abundances derived from low-excitation OH lines, agreed well with those derived from high-excitation lines of the O i IR triplet at 7774 Å. The comparison with oxygen abundances derived using O i data from Tomkin et al. (1992) showed a mean difference of $`0.00\pm 0.11`$ dex for the stars in common. Boesgaard et al. (1999a) made a similar analysis of several metal-poor stars using a different set of OH lines. They found a very good agreement with the results obtained by Israelian et al. (1998), and basically the same dependence of \[O/Fe\] versus metallicity. This is clearly seen in the upper panel of Figure 1. The UV “missing opacity” problem discussed by Balachandran & Bell (1998), which could affect both oxygen and beryllium abundance determinations from these lines, has been studied recently by Allende Prieto & Lambert (2000). These authors have found a good agreement between $`T_{\mathrm{eff}}`$s obtained from the IRFM and from the near-UV continuum for stars with $`4000T_{\mathrm{eff}}6000`$ K when accurate Hipparcos gravities are used. This also agrees with our good reproduction of the near-UV spectral region using the IRFM temperatures. This result indicates that the model atmospheres used provide an adequate description of the near-UV continuum forming region. In any case, even if a not well understood opacity problem would exist as described by Balachandran & Bell, it would have a minor effect on the OH results since most of the stars in the samples of Israelian et al. and Boesgaard et al. are hotter than the Sun and very metal-poor. The corrections to oxygen abundances for individual stars would be lower than 0.15 dex, not changing significantly the \[O/Fe\] vs. \[Fe/H\] trend. A new non-LTE analysis of the O i IR triplet for a sample of 38 metal-poor stars performed by Mishenina et al. (2000) gives consistent results with those of Abia & Rebolo (1989), Tomkin et al. (1992) and Kiselman (1993), and indicates that the mean value of the non-LTE correction in unevolved metal-poor stars is typically 0.1-0.2 dex. These authors confirmed the \[O/Fe\] vs \[Fe/H\] trend discussed by Israelian et al. (1998) and Boesgaard et al. (1999a) from the OH lines, without finding any trend of oxygen abundances with $`T_{\mathrm{eff}}`$ or $`\mathrm{log}g`$. It is also worthwhile to mention that the O i IR triplet is not affected by 3D effects, convection and small-scale inhomogeneities in the stellar atmosphere (Asplund et al. 1999). In addition, oxygen abundances derived form this triplet are not significantly affected by chromospheric activity either. The central panel of Fig. 1 shows a compilation of oxygen abundances derived using the IR triplet. The larger scatter observed as compared with the measurements based on OH lines can be associated with the different scales of stellar parameters ($`T_{\mathrm{eff}}`$, gravities, and metallicities) adopted by the authors of each set of stars, and to the fact that some measurements have not been corrected for non-LTE effects. Very recently, Carretta, Gratton, & Sneden (2000) performed an independent analysis of 32 metal-poor stars hotter than 4600 K using the IR triplet, and provide LTE and non-LTE oxygen abundances which are significantly lower than those found in previous works. A preliminary attempt to understand the reasons for this discrepancy can be done by looking in detail into their most metal-poor star (BD +3740, \[Fe/H\]$`=2.66`$) where a surprisingly low oxygen abundance \[O/Fe\]$`=0.38`$ is claimed. A recent study of stellar parameters based on the non-LTE analysis of iron lines (Thévenin & Idiart 1999), gives a lower effective temperature (by 140 K) and a higher gravity (by 0.3 dex) than the values adopted by Carretta et al. for this star. Using these latter parameters we obtain an LTE oxygen abundance 0.4 dex higher than Carretta et al. (i.e. \[O/Fe\]$`{}_{\mathrm{LTE}}{}^{}=1.05`$), and for a star with these parameters the non-LTE correction to the oxygen abundance is of the order of 0.05 dex (Mishenina et al. 2000), much lower than the 0.25 dex value used by Carretta et al. We therefore arrive at a value \[O/Fe\]$`{}_{\mathrm{NLTE}}{}^{}1.0`$, in good agreement with the OH determination by Boesgaard et al. (1999a). Corrections to the stellar parameters as inferred from the non-LTE analysis of Fe lines clearly have an impact on the oxygen abundances which we will address in a forthcoming paper. Israelian et al. (1998) found four dwarfs in their sample for which oxygen abundances derived using \[O i\] were in good agreement with those derived from OH when Hipparcos gravities are used. Several oxygen measurements for unevolved stars based on the \[O i\] 6300 Å line are compiled in the lower panel of Fig. 1. This figure shows a similar trend than that observed for the abundances recently derived from forbidden lines by Carretta et al. (2000). The presence of a linear trend of \[O/Fe\] versus metallicity in Fig. 1 strongly depends on the only two measurements available at \[Fe/H\]$`2`$. These two measurements have been reported by Fulbright & Kraft (1999) for the subgiants BD +371458 and BD +233130, which were also considered by Israelian et al. (1998) and Boesgaard et al. (1999a; only BD +371458 in this case). The analysis carried out by Fulbright & Kraft is based on gravities derived from LTE iron ionization balance of these subgiants where it is well known that non-LTE effects are strong (Thévenin & Idiart 1999; see also Idiart & Thévenin, this conference). Allende Prieto et al. (1999) have shown that gravities derived using this technique in metal-poor stars do not agree with the gravities inferred from accurate Hipparcos parallaxes. They find that gravities are systematically underestimated when derived from ionization balances and that upward corrections of $`0.5`$ dex can be required at metallicities similar to those of our stars, in good agreement with Thévenin & Idiart. We remark here that any underestimation of gravities will also strongly underestimate the abundances inferred from the forbidden line. For the two stars under discussion our Hipparcos based gravities are 0.45 and 1.05 dex (for BD +371458 and BD +233130, respectively) higher than derived by Fulbright & Kraft, and would imply corrections in the oxygen abundances similar to those indicated in Fig. 1 (a detailed analysis would imply also the correction for the assumed metallicities). Our conclusion is that the uncertainties in the gravities of these subgiants allow the abundances inferred from the forbidden line to be consistent with those estimated from the OH lines or the triplet. Actually, consistency with the other oxygen indicators is achieved for the high gravities inferred from Hipparcos when consistent analyses are made, and this could be taken as an indication that the high gravities are indeed the correct ones. ## 4. Beryllium and oxygen The dependence of $`\mathrm{log}`$(Be/H) on \[Fe/H\] and on \[O/H\] (using the abundances derived from the OH UV lines) is essentially linear (García López 1999; Boesgaard et al. 1999b), but with different slopes: $`1.1`$ and $`1.5`$, respectively. No evidence of a primordial plateau of Be down to $`\mathrm{log}`$(Be/H)=$``$13.5 is found. Figure 2 shows the increase of the \[Be/O\] ratio with increasing metallicity and a slope of $`0.4`$. This relation provides an observational constraint to the Galactic Cosmic Ray theories. Three types of GCR models exist at present which try to explain their observed evolution. These are 1) a pure primary GCR from superbubbles (Ramaty, this conference), 2) a hybrid model based on GCR and superbubble accelerated particles (Cassé, this conference), which could be accomplished by a pure superbubble model (Parizot & Drury, this conference), and 3) standard GCR (Olive, this conference). Apparently all these models can be adopted for both, variable and flat \[O/Fe\]. However, models presented by R. Ramaty and K. Olive show more consistency when variable \[O/Fe\] is adopted. Chemical evolution models of the early Galaxy where stellar lifetimes are taken into account and assuming that Type Ia SN appear at a Galactic age of 30 million years can also explain the evolution of oxygen delineated in Fig. 1. (Chiappini et al. 1999.). The evolution of oxygen proposed in this paper also helps to understand the evolution of <sup>6</sup>Li versus \[Fe/H\] and the <sup>6</sup>Li/Be ratio at low metallicities in the framework of standard Galactic Cosmic Ray Nucleosynthesis (Fields & Olive 1999). In addition, Ramaty et al. (1999) have proposed that a delay between the effective deposition times into the ISM of Fe and O (only a fraction of which condensed in oxide grains) can explain a linear trend of \[O/Fe\]. It has been suggested (Vangioni-Flam & Cassé, this conference) to use magnesium as metallicity indicator instead of oxygen. However, given the existence of unevolved halo stars with negative \[Mg/Fe\] ratios (McWilliam 1997; Carney et al. 1997), this approach may not lead to better results. For example, the subdwarf BD+3740 has \[O/Fe\]$`1`$ (see previous Section) while its \[Mg/Fe\]$`=0.28`$ (Fuhrmann et al. 1995). Yield of Mg depends on the extent of mixing (Argast et al. 2000) and its primordial abundance can be changed due to the operation of the MgAl cycle. ## References Abia, C., & Rebolo, R. 1989, ApJ, 347, 186 Allende Prieto, C., García López, R. J., Lambert, D. 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warning/0003/astro-ph0003409.html	ar5iv	text	# 1 Introduction ## 1 Introduction Although photoionization and recombination are inverse processes as they occur in nature, they are usually treated in independent theoretical frameworks. This basic inconsistency, directly related to ionization balance in radiatively ionized media, and consequent inaccuracies, propagate through to the photoionization models employed in astrophysics. A further division, largely artificial, is made in theoretical methods used to compute electron-ion recombination rates. Two sets of data are usually calculated: (i) ’radiative recombination’ (RR), calculated using background, or non-resonant, photoionization cross sections, and (ii) ’di-electronic recombination’ (DR) representing the contribution of autoionizing resonances, first shown to be important by Burgess (1964). That this procedure is not only theoretically unsatisfactory, but also impractical in most cases, is seen from both theoretical calculations and experimental measurements of photoionization and recombination cross sections. The simple reason is that the resonances are inseparable from the background. The cross sections contain, in general, extensive and interacting Rydberg series of resonances; the non-resonant and resonant contributions are not accurately separable (except, possibly, for few-electron, highly charged ions). The large number of photoionization cross sections computed under the Opacity Project exhibit these features in detail, particularly for many electron systems (The Opacity Project Team 1995, 1996 - compiled publications and data) In addition, the cross sections for photoionization and recombination of excited states, particularly metastable levels, may contain even more complicated resonances than the ground state (Luo et al 1990). Experimentally, of course, the measurements always yield a combined ’(RR + DR)’ cross section (albeit in limited energy ranges usually accessible in experimental devices). Therefore a theoretical method that accounts for both the resonant and the non-resonant recombination in a unified manner is desirable, and has been developed (e.g. Nahar and Pradhan 1994, Zhang et al 1999), based on the close coupling (CC) approximation using the R-matrix method (Burke & Seaton 1984, Berrington et al 1987, Hummer et al 1993) as used in the Opacity Project and the Iron Project (hereafter OP and IP). Photoionization cross sections may be computed essentially for all bound states, level of excitation ($`n,\mathrm{},SL\pi ,SLJ\pi )`$, energy range, and with resolution of resonances. In principle, the cross section for the inverse photo-recombination process is given by detailed balance. However, since recombination takes place to an infinite number of bound states of (e + ion) system, it becomes impractical (and as it turns out, unnecessary) to do so for the very higly excited levels above a certain n-value (chosen to be 10 in practice). For recombination into levels with n$`>`$10, the non-resonant contribution, relative to the resonant contribution per unit energy is negligible owing to the density of resonances as n $`\mathrm{}`$. In that range we employ a precise theoretical treatment of DR based on multi-channel quantum defect theory and the CC approximation (Bell & Seaton 1985, Nahar & Pradhan 1994) to compute the recombination cross section. Among the problems that manfiest themselves in the CC photoionization/recombination calculations are: the accuracy and convergence of the eigenfunction expansion for the ion, relativistic fine structure effects, the contribution from non-resonant recombination into high-n levels as $`E0;n\mathrm{}`$, and resolution of narrow resonances with increasing n and/or $`\mathrm{}`$, and radiation damping thereof. Experimental work is of importance in ascertaining the accuracy of theoretical calculations and the magnitude of various associated effects, since most of the photoionization/recombination data can only be calculated theoretically. In recent years there has also been considerable advance in the measurements of (e + ion) recombination cross sections on ion storage rings (e.g. Kilgus et al, 1990, 1993, Wolf et al. 1991), and photoionization cross sections using accelerator based photon light sources (R. Phaneuf et al., private communication). We compare the CC calculations for both atomic processes with the latest experimental data. ## 2 Theory The CC approximation takes account of the important coupling between the energtically accessible states of the ion in the (e + ion) system. The target ion is represented by an $`N`$-electron system, and the total wavefunction expansion, $`\mathrm{\Psi }(E)`$, of the ($`N`$+1) electron-ion system of symmetry $`SL\pi `$ or $`J\pi `$ may be represented in terms of the target eigenfunctions as: $$\mathrm{\Psi }(E)=A\underset{i}{}\chi _i\theta _i+\underset{j}{}c_j\mathrm{\Phi }_j,$$ (1) where $`\chi _i`$ is the target wavefunction in a specific state $`S_iL_i\pi _i`$ or $`J_i\pi _i`$ and $`\theta _i`$ is the wavefunction for the ($`N`$+1)-th electron in a channel labeled as $`S_iL_i(J_i)\pi _ik_i^2\mathrm{}_i(SL\pi \mathrm{or}J\pi ))`$; $`k_i^2`$ being its incident kinetic energy. $`\mathrm{\Phi }_j`$’s are the correlation functions of the ($`N`$+1)-electron system. Bound and continuum wavefunctions for the (e + ion) system are obtained on solving the CC equations at any total energy E $`<`$ 0 and E $`>`$ 0 respectively. The coupling between the eigenfunctions of the energetically inaccessible states of the ion (’closed channels’), and the accessible states (’open channels’), gives rise to resonance phenomena, manifested as infinite Rydberg series of resonances converging on to the excited states of the target ion. With the bound and the continuum (free) states of the (e + ion) system, atomic cross sections may be calculated for electron impact excitation (EIE), photoionization, and recombination — the free and the bound-free processes. Radiative transition (bound-bound) probabilities may also be obtained. The R-matrix method, and its relativistic extension the Breit-Pauli R-matrix method (BPRM), are the most efficient means of solving the CC equations, enabling in particular the resolution of the resonances in the cross sections at a large number of energies. ## 3 Photoionization Photoionization calculations are carried out for all levels ($`n,\mathrm{},SL\pi ,SLJ\pi `$, with $`n10,\mathrm{}=n1`$. Typically this means several hundred bound levels of each atom or ion. The cross section for each level is delineated at about thousand energies, or more, to map out the resonance structure in detail. In recent studies, photoionization cross sections of all ions of C, N, and O were computed (Nahar and Pradhan 1997, Nahar 1998) with more extended eigenfunction expansions, resolution of resonances, and number of levels, than the earlier OP data. Overall, new photoionization data for over 40 atoms and ions, with improvements over the OP data (e.g. currently in TOPbase, Cunto et al. 1993), has now been calculated for: low ionization stages of iron: Fe I, II, III, IV, and V (references in Bautista and Pradhan 1998), Ni II (Bautista 1999), the C-sequence ions (Nahar & Pradhan 1991,1992), the Si-sequence ions (Nahar & Pradhan 1993). Unified (e + ion) photo-recombination cross sections and rates (total and level-specific) have also been obtained, as discussed below. ### 3.1 Comparison with experiments The recent ion-photon merged beam experiment by Kjeldsen et al. (1999) on the photoionization cross sections of the ground state of C II shows an extremely rich and detailed resonance structures (Fig. 1). There is excellent agreement between the theory and experiment, both in terms of magnitude and details of the background and resonances. However, the theoretical calculations were in LS coupling, neglecting fine structure, that clearly manifests itself in the additional peaks seen in the experimental cross sections (new relativistic calculations are in progress). New photoionization experiments have been carried out for positive atomic ions at the Advance Light Source (ALS) in Berkeley, where a photon light source is used on an accelerator that produces the ion beams. These extremely high resolution measurements provide an unprecedented check on the details of the theoretical cross sections, particularly resonance structures and fine structure effects. Fig. 2 compares the O II cross section from theory (Nahar 1998), and experiments at the ALS done by the Reno group headed by R. Phaneuf. The experimental cross sections include not only the photoionization of the ground state $`2s^22p^3(^4S^o)`$ but also the metastable excited states $`2s^22p^3(^2D^o,^2P^o)`$. $$h\nu +OII(2s^22p^3{}_{}{}^{4}S_{}^{o},^2D^o,^2P^o)e+OIII(2s^22p^2{}_{}{}^{3}P,^1D,^1S)$$ The complicated features arise from several series of resonances converging on to the excited states of the residual ion O III. There is very good agreement between the CC calculations and experiment, verifying the often expressed, but not heretofore established, claim of about 10% accuracy of the theoretical cross sections. The situation may be more complicated, and the uncertainties larger, for more complex atomic systems. These results also show that metastable states may need to be included in atomic photoionization models of astrophysical sources. ## 4 Unified method for (e + ion) recombination Photoionization calculations described above are for total photoionization from a given level into all excited levels of the residual ion. However, for photo-recombination the calculations must be repeated to obtain the cross section into the ground state of the ion alone. Detailed balance then applies precisely (Milne relation) as $$\sigma _{\mathrm{RC}}(ϵ)=\frac{\alpha ^2}{4}\frac{g_i}{g_j}\frac{(ϵ+I)^2}{ϵ}\sigma _{\mathrm{PI}},$$ (2) where $`\sigma _{\mathrm{RC}}(i_o)`$ is the photo-recombination cross section, $`\sigma _{\mathrm{PI}}`$ is the photoionization cross section into the ground state $`i_o`$, $`\alpha `$ is the fine structure constant, $`ϵ`$ is the photoelectron energy, and $`I`$ is the ionization potential in Rydberg atomic units. Recombination can take place into the ground or any of the excited recombined (e+ion) states. The contributions of these bound states to the total $`\sigma _{\mathrm{RC}}`$ are obtained by summing over the contributions from individual cross sections. $`\sigma _{\mathrm{RC}}`$ thus obtained from $`\sigma _{\mathrm{PI}}`$, including the autoionizing resonances, corresponds to the total (DR+RR) unified recombination cross section in an ab initio manner. Recombination into the high-$`n`$ states, i.e. $`n_{\mathrm{max}}<n\mathrm{}`$, is computed assuming DR to dominate over the non-resonant background contribution. The CC approximation can then be used to calculate DR collision strengths $`\mathrm{\Omega }_{\mathrm{DR}}`$, as an extension of the theory of DR by Bell and Seaton (1985). These two main parts of the unified recombination calculations, and other parts, are described in detail in Zhang et al (1999). For very highly charged ions, such as the H- and He-like ions with large radiative decay rates for core transitions, radiation damping effects can be significant. As in other CC calculations for excitation and photoionization, resonances are resolved at a suitably fine mesh to enable perturbative radiative damping, and to ensure that the neglected resonances do not significantly affect the computed rates. Relativistic fine structure is considered in the BPRM calculations for highly charged ions. ### 4.1 Comparison with experiments Although experimental results are available for relatively few ions in limited energy ranges, and mostly for simple atomic systems such as the H-like and He-like ions, they are useful for the calibration of theoretical cross sections. Zhang et al (1999) have compared in detail the BPRM cross sections with experimental data from ion storage rings for $`e+CVCIV,e+CVICV,e+OVIIIOVII`$, with close agreement in the entire range of measurements for both the background (non-resonant) cross sections and resonances. The reported experimental data is primarily in the region of low-energy resonances that dominate recombination (mainly DR) with H- and He-like ions. The recombination rate coefficients, $`\alpha _R`$, obtained using the cross sections calculated by Zhang et al. agree closely with those of Savin (1999) who used the experimental cross sections to obtain ’experimentally derived DR rates’. However, these rates do not include contributions from much of the low energy non-resonant RR and very high energy regions. The total unified $`\alpha _R(T)`$ which include all possible contributions is, therefore, somewhat higher than that obtained from limited energy range. In Fig. 3, the solid curve corresponds to the total unified $`\alpha _R`$. The dotted curve and the dot-long-dash curves are the rates using cross sections from Zhang et al (1999) and Savin (1999) respectively, in the limited energy range in experiments (the two curves almost merge). The short-and-long dash curve is the total $`\alpha _R`$ in LS coupling (Nahar & Pradhan 1997) which, at high temperatues, is higher than the new BPRM rates including fine structure and radiation damping (solid curve). The dot-dash curve is the DR rate by Badnell et al (1990), which is lower than the others. the dashed and the long-dashed curves are RR rates by Aldrovandi & Pequignot (1973), and Verner and Ferland (1996); the latter agrees with the present rates at lower temperatures. ### 4.2 Photoionization/recombination of Fe XXV Fe XXV is one of the most important ions in X-ray spectroscopy (see the review article by Pradhan in this volume). We have completed the calculations for: (i) photoionization cross sections for fine structure levels up to n = 10, including those for ionization into the ground level, and (ii) total and level-specific unified recombination cross sections and rate coefficients. The most commonly observed lines of Fe XXV correspond to the K-$`\alpha `$ transitions between the $`n`$= 1 and 2 levels: (1) the ’z’ line $`1s2s(^3S_1)1s^2(^1S_0)`$, (2) the ’y’ line, $`1s2p(^3P_1^o)1s^2(^1S_0)`$, (3) the ’x’ line, $`1s2p(^3P_2^o)1s^2(^1S_0)`$, and (4) the ’w’ line $`1s2p(^1P_1^o)1s^2(^1S_0)`$. Recombination rate coefficients into these levels are given in Fig. 4; these vary smoothly with temperature, except a slight ”shoulder” at high temperature due to DR. ## 5 Ionization equilibrium The new photoionization/recombination data should enable more accurate calculations for photoionization equilibrium $$_{\nu _0}^{\mathrm{}}\frac{4\pi J_\nu }{h\nu }N(X^z)\sigma _{PI}(\nu ,X^z)𝑑\nu =\underset{j}{}N_eN(X^{z+1})\alpha _R(X_j^z;T),$$ (3) and, coronal equilibrium $$C_I(T,X^z)N_eN(X^z)=\underset{j}{}N_eN(X^{z+1})\alpha _R(X_j^z;T),$$ (4) where $`\alpha _R(X_j^z;T)`$ is the total electron-ion recombination rate coefficient of the recombined ion of charge $`z`$, $`X_j^z`$, to state j at electron temperature T, $`C_I`$ is the rate coefficient for electron impact ionization, and $`\sigma _{PI}`$ is the photoionization cross section evaluated at photon frequency $`\nu `$ and convoluted with the isotropic radiation density J<sub>ν</sub> of the source; N<sub>e</sub>, $`N(X^{z+1})`$, and $`N(X^z)`$ are the densities for the free electrons, and the recombining and recombined ions respectively. Coronal ionization fractions for C,N,O using the unified recombination rates are also computed (Nahar and Pradhan 1997; Nahar 1999). ## 6 Conclusion We carry out ab initio large scale close coupling R-matrix calculations for (i) photoionization cross sections, and (ii) electron-ion recombination rate coefficients. The predicted theoretical features in $`\sigma _{PI}`$ are being observed in the recent sophisticated experiments. The unified method for (e+ion) recombination has been benchmarked with available experimental measurements. Our study of unified electron-ion recombination rates exhibit a general pattern with temperature. Although generally applicable to all systems, the close coupling unified method is especially suitable for the strong coupling cases where the broad and overlapping resonances dominate the near-threshold region in the electron-ion recombination process, and other methods may not be accurate. Total and state-specific unified recombination rate coefficients and photoionization cross sections are available for about 40 atoms and ions: Carbon: C I, C II, C III, C IV, C V, C VI Nitrogen: N I, N II, N II, N IV, N V, N VI, N VI Oxygen: O I, O II, O III, O IV, O V, O VI, O VII, O VII C-like: F IV, Ne V, Na VI, Mg VII, Al VIII, Si IX, S XI Si and S: Si I, Si II, S II, S III, Ar V, Ca VII Iron: Fe I, Fe II, Fe III, Fe IV, Fe V, Fe XIII, Fe XXV These photoionization/recombination datasets are self-consistent, and should yield more accurate astrophysical photoionization models. 4. Acknowledgements This work is supported partially by the NSF and NASA. REFERENCES Aldrovandi, S.M.V. & Pequignot, D. 1973, Astron. Astrophys. 25, 137 Badnell, N.R., Pindzola, M.S., & Griffin, D.C. 1990, Phys. Rev. A 41, 2422 Bautista, M.A. 1999, Astron. Astrophys. Suppl. 137, 529 Bautista, M.A. & Pradhan A.K. 1988, Astrophys. J. 492, 650 Bell, R.H. and Seaton, M.J. 1985, J. Phys. 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A 58, 4593 Nahar, S.N. 1999, ApJS 120, 131 Nahar, S.N., & Pradhan, A.K. 1991, Phys. Rev. A 44, 2935 Nahar, S.N., & Pradhan, A.K. 1992, Phys. Rev. A 45, 7887 Nahar, S.N., & Pradhan, A.K. 1993, J. Phys. B 26, 1109 Nahar, S.N., & Pradhan, A.K. 1994, Phys. Rev. A 49, 1816 Nahar, S.N., & Pradhan, A.K. 1997, ApJS 111, 339 Savin, D.W. 1999, Astrophys. J. 523, 855 The Opacity Project 1 & 2, compiled by the Opacity Project team (Institute of Physices, London, UK, 1995,1996) Wolf A, Berger J, Bock M, Habs D, Hochadel B, Kilgus G, Neureither G, Schramm U, Schwalm D, Szmola E, Müller A, Waner M, and Schuch R 1991 Z. Phys. D Suppl. 21, 569 Verner, D.A. & Ferland G. 1996, Astrophys. J. Suppl. 103, 467 Zhang, H.L., Nahar, S.N., & Pradhan, A.K. 1999, J. Phys. B 32, 1459
warning/0003/math-ph0003030.html	ar5iv	text	# 1 Introduction ## 1 Introduction Nonlinear dynamics, which is a relatively new field of study, is a very important frontier for probing natural phenomena. Active research efforts that focus on the mathematics and physics of nonlinear dynamical systems have emerged worldwide in various fields, including fluid dynamics, plasma physics, astrophysics, and even string theory . Among its unique features is the ability to describe a variety of patterns and particle-like traveling solutions . Other notable features of the theory include its description of solitons and breather modes, features in quantum optics , molecular and solid state physics phenomena , and solitons in nuclear and particle physics . In contrast with linear theories which exhibit smooth regular motion, nonlinear models require nonlinear partial differential equations (NPDE) and show strong couplings between different mechanisms and parts of the system. Also, the nonlinear interactions involve multiple scales and are related with self-similar patterns or fractals. The NPDE solutions of physical interest are mostly localized and demonstrate good stability in time and through scattering with each other. Their shapes are related to the velocity, thus making the nonlinear patterns distinct from linear waves. In the asymptotic domain these solutions consist of isolated traveling pulses that are free of interactions. Close to the scattering domain, the nonlinear solutions obey nonlinear superposition principles. The main challenge in any nonlinear analysis is the construction of localized or finitely supported analytical solutions for the NPDE of interest. This challenge includes issues regarding the inexistence of a superposition principle for such solutions. Recent examples show that the traditional nonlinear tools (inverse scattering, group symmetry, functional transforms) are not always applicable . On the other hand, from an experimental point of view one knows that patterns which are observed in nature – either stationary, growing, or propagating – generally have finite space-time extension and a multi-scale structure. Since soliton, even when they are localized, have an infinite extent, one needs other appropriate structures, and eventually self-similar bases. In this paper, wavelet-inspired approaches for localized solutions of NPDE are explored. We propose a new similarity formalism for the qualitative analysis, and clasiffication, of soliton solutions of nonlinear equations. This method provides relations between the characteristics of such solutions (amplitude, width and velocity) without the need of solving the corresponding NPDE. The method uses the multi-resolution analysis where traditional tools like the Fourier integrals or linear harmonic analysis are inadequate for describing the system. Wavelets are functions that have a space-dependent scale which renders them an invaluable tool for analyzing multi-scale phenomena. Wavelets have been used in signal processing, in problems involving singular potentials, pattern recognition, image compression, turbulence and even in radar and acoustic problems . Moreover, the introduction of wavelet analysis in a study of NPDE is very natural because they can acommodate everything from strong variations, even singularities, to a smooth behavior. The second purpose of the paper is to show an example of construction of a nonlinear basis for a NPDE with nonlinear dispersion. There are many physical reasons favoring wavelets in the construction of nonlinear bases. For example, the breakup process of fluid drops has been shown to be self-similar, and in particular, the singularities (necks) look identical at any scale . Other nonlinear oscillations of liquid drops, shells, bubbles, or even neutron stars involve such types of behavior . We introduce in section 2 a qualitative similarity analysis for NPDE that yields relations between the amplitude, width and velocity of their traveling solutions, with many examples and predictions. In section 3, we introduce a nonlinear wavelet-like frame, associated with localized analytic solutions of a modified KdV equation, with nonlinear dispresion. ## 2 Qualitative similarity analysis Finding analytic solutions for the nonlinear equations which describe physical phenomena is rather an exception than the rule. The traditional nonlinear tools like the inverse scattering theory or functional transforms are not always applicable . Also, when phenomena of interest have many space (or time) scales, hundred times less or smaller than the dimension (or time scale) of the whole system, numerical methods may fail. This is the case of sharp propagating perturbations developing discontinuities, or shock waves. A simple option is the expansion of solutions in a basis of appropriate chosen linear modes. The most common choice are the Fourier series which have the advantage of orthogonality, but can hardly discriminate local behavior of phenomena. Moreover, information about the order of magnitude of the Fourier coefficients of a signal or wave $`u(x,t)`$ is not sufficient for making conclusions about the size or scale of $`u`$. In these special situations, the use of bases of multiresolution analysis and wavelets has become popular . We introduce a qualitative analysis for the localized traveling solutions belonging to any type of NPDE, in terms of Morlet continuous wavelet approach. In our case of interest, traveling localized perturbations, the most important information are provided by the scale of the pulse (or the width) denoted $`L`$, the amplitude $`A`$, and the group velocity $`V`$. The qualitative analysis in this section provides simple relations between these three parameters, without actually solving the equation. It also gives an estimation about the specific scale of the solutions. The procedure consistsin the substitution of all the following terms in the NPDE, according to the rule $$u_t\pm Vu_x,u\pm A,u_x\pm A/L,u_{xx}\pm A/L^2\mathrm{},$$ (1) and so forth for higher order of derivatives. Consequently, the NPDE is mapped into an algebraic equation in $`A,L`$ and $`V`$. In Table 1 we present several examples of application of this substitution to some well-known and widely used NPDE in physics. The validity of the method follows from the expansion of the soliton-like solution $`u(x)`$ in Morlet wavelets $$\mathrm{\Psi }_\alpha (x)=\pi ^{1/4}e^{i\alpha x\frac{x^2}{2}},$$ (2) where $`\alpha `$ describes the scale of this mother wavelet. The support of any Morlet wavelet is mainly confined in the $`(1,1)`$ interval. We have the discrete Morlet wavelet expansion of $`u`$ $$u(x)=\underset{j}{}\underset{k}{}C_{j,k}\mathrm{\Psi }_\alpha (2^jxk),$$ in terms of integer translations and dyadic dilations of the mother wavelet. In order to reduce the number of scales needed, the range of the summation should be choosen with an eye to the underlying physics. We use in the following the asymptotic formula describing the pointwise behavior of the Morlet wavelet series around of a point $`x_0`$ of interest . For a chosen $`x_0`$ and scale $`j`$, there is only one $`k`$ and $`\|ϵ\|1`$ such that the support of the corresponding $`\mathrm{\Psi }_{\alpha ,j,k}`$ contains this point, $`k=2^jx_0+ϵ`$. We can express the solution and its derivatives in a neighborhood of this point $$u(x_0)\mathrm{\Psi }(ϵ)\underset{j}{}C_{j,2^jx_0+ϵ}\underset{j}{}u_j(x_0),$$ $$\frac{du}{dx}(x_0)i\mathrm{\Psi }(ϵ)\underset{j}{}2^j\alpha C_{j,2^jx_0+ϵ}=i\alpha \underset{j}{}2^ju_j(x_0),$$ (3) for $`\alpha `$ chosen enough large compared to $`ϵ`$. Since the coefficient $`1/\alpha 2^j`$ represents the scale for each $`\mathrm{\Psi }_{\alpha ,j,k}`$ Morlet wavelet, we can define it as a characteristic half-width $`L_j`$. And we finaly have in $`x_0`$, from eqs.(3) $$\frac{d^nu}{dx^n}(x_0)\underset{j}{}\frac{u_j(x_0)}{L_j^n},$$ (4) where $`u_j(x_0)=\mathrm{\Psi }_\alpha (ϵ)C_{j,2^jx_0+ϵ}\mathrm{\Psi }_\alpha (0)C_{j,2^jx_0}`$. Eq.(4) is the many scales generalization of the simpler formula in eq(1). With eq.(4) in hand we can investigate the structure of hypothetic soliton solutions, by choosing $`x_0`$ in the neighborhood of the maximum value of the solution, $`u(x_0,0)=A`$. Around this maximum, such solutions can be described very well by a unique scale $`L`$, and hence the solution and its derivatives can be approximated with the corresponding dominant term, by the substitutions in eqs.(1). In Table 1 we present a series of examples of NPDE, identified in the first column by the name and the form of the equation. In the second column, we write the corresponding traveling localized solution, if an analytical form is known. Such solutions provide special relations between $`L,A`$ and $`V`$, which are given in the third column. In the last column we introduce for comparison the results of this Morlet qualitative analysis, that is the relations between the three parameters, provided by eqs.(1). The usefulness of the approach may be checked, by a quick comparison between the fourth and the fith columns. While the results in the fourth column are possible only when one knows the analytical solutions, the results presented in the last column, obtained by the similarity approach, result directly from the NPDE, without actually solving it. The first line in the Table presents a linear case, for comparative purposes. The above Morlet wavelet approximation provides a correct expression for the dispersion relation ($`V=c`$ $``$ $`k^2=\omega ^2/c^2`$) with no constraint on either the amplitude $`A`$ or on the width $`L`$. The case of the KdV equation is described in the second row of the Table 1. The method gives a general expression for $`L=L(A,V)`$. For $`L`$ to be related to $`A`$ only, from the fourth column it results that the velocity $`V`$ must be proportional to $`A`$. In this case we obtain exactly the well-known relation (column three) among the parameters in the exact solution. A prediction of the method is if we allow $`V`$ to depend on a power of $`A`$. This means solutions with a higher nonlinear coupling between the shape and kinematics. A side effect would be a lower limit for $`A`$. Smaller solitons than this limit can move vith velocity proportional to the amplitude only. The same result is obtained for the MKdV equation (third row), except that in this case $`V`$ needs to be proportional to the square of $`A`$ in order to have $`L`$ a simple function of $`A`$ only. This prediction is again identical with that in the exact solution (third column). Moreover, the same relations remain valid even for the new exotic solutions of the MKdV equation of compacton type, $$u(x,t)=\frac{\sqrt{32}k\mathrm{cos}[k(x4k^2t)]^2}{3(1\frac{2}{3}\mathrm{cos}[k(x4k^2t)]^2)},$$ which has $`L=5\pi /6k`$ or $`\pi /6k`$, that is $`L1/A`$ like in the Table 1. Next example (4-th row) is provided by a generalised KdV equation, in which the dispersion term is quadratic $$\eta _t+(\eta ^2)_x+(\eta ^2)_{xxx}=0.$$ (5) Eq.(5), known as K(2,2) from the two quadratic terms, admits compact supported traveling solutions, named compactons . The compactons are powers of trigonometric functions defined on a half-period, and zero otherwise. In general, they have the form $`Acos^ad(xct)`$, and different from solitons, their width is independent of the amplitude. This is the fact that provides a connection with wavelet bases. They are characterize by a unique scale, and it is this feature that makes it possible to introduce a nonlinear basis starting from this “mother” function. For eq.(5) the compacton solution is given by $`\eta _c(xVt)={\displaystyle \frac{4V}{3}}cos^2\left[{\displaystyle \frac{xVt}{4}}\right],`$ (6) if $`\|xVt\|<2\pi `$ and zero otherwise. Here the velocity is a function of the amplitude. Notice that the width $`L=4`$ of the wave is independent of the amplitude. The quadratic dispersion term is characteristic for the nonlinear coupling in a chain. The general compacton solution for eq.(5) is actually a ”dilated” version of eq.(6). That is, a combination of the first rising half of the squared cos in eq.(6), followed by a flat domain of arbitrary length ($`\lambda `$), and finaly followed by the second, descending part of eq.(6). Actually, this combination is just a kink compacton joined smoothly with an antikink one $$\eta _{kak}(xVt;\lambda )=\{\begin{array}{cc}0\mathrm{}\hfill & \\ \frac{4V}{3}cos^2\left[\frac{xVt}{4}\right],2\pi xVt0\hfill & \\ \frac{4V}{3},0xVt\lambda \hfill & \\ \frac{4V}{3}cos^2\left[\frac{xVt\lambda }{4}\right],\lambda xVt\lambda +2\pi \hfill & \\ 0\mathrm{}\hfill & \end{array}$$ (7) In Fig. 1 we present a compacton, eq.(6), a kink-antikink pair (KAK) described by eq.(7), both with the same amplitude and velocity. Although the second derivative of this generalized compacton is discontinuous at its edges, the KAK, eq.(7), is still a solution of eq.(5) because the third derivative acts on $`u^2`$, which is a function of class $`C_3`$. Finally, we can construct solutions by placing a compacton on the top of a KAK, like in the third solution in Fig. 1. Such a solution exists only for a short interval of time, since the two structures have different velocities. The solution is given by $$\eta (x,t)=\eta _{kak}(xVt;\lambda )+\left(\eta _c(xV^{}t2\pi )+\frac{4V}{3}\right)\chi (\frac{xV^{}t2\pi }{2\pi }),$$ (8) for $`0<t<(\lambda 4\pi )/(V^{}V)`$. Here $`\chi (x)`$ is the support function, equal with 1 for $`\|x\|1`$ and 0 in the rest, and $`V^{}=3\text{max}\{\eta _c\}/4+2V`$. For the K(2,2) compacton, eq.(6), the exact relations between the parameters are $`A=4V/3`$ and $`L=4`$, . The relation provided by the similarity method, in the last column of the forth row, predicts the existence of the compacton. That is, for a linear dependence between the amplitude and the speed, the half-width is constant and does not depend on $`A`$ or $`V`$. This fact ($`LL_0=const.`$) is a typical feature of K(2,2) compactons. Moreover, in literature there was found numerically that for any compact supported initial data, widder than $`L_0`$, the solution decomposes in time into a series of $`L_0`$ compactons, Fig. 2. For narrower initial data the numeric solution blows up. There is no exact or analytic explanation of this effect, so far. The similarity method can give a hint in this situation, too, by using the graphic of the relation $`L=L(V,A)`$ provided by this qualitative method. In Fig. 3 $`L`$ is ploted versus $`V`$, for several values of $`A`$ (larger values of $`A`$ translate the curves to the right). The half-width of a stable compacton was chosen $`L_0=0.707`$. Above this value, Fig. 3a, wider compact pulses produce an intersection for each curve (each $`A`$) with the axis $`L_0`$ providing series of compactons of different heights, like in the numerical experiments, . Below this $`L_0`$ line, all the curves approache infinite amplitude, providing instability of narrower shapes. Another good example of prediction of the method is exemplified in the case of a general convection-nonlinear dispersion equations, denoted K(n,m) $$\eta _t+(\eta ^n)_x+(\eta ^m)_{xxx}=0.$$ (9) Compacton solution for any $`nm`$ are not known in general, except some particular cases. In this case we find a general relation among the parametrs, for any $`n,m`$, shown in the $`5^{th}`$ and $`6^{th}`$ rows. These general relations $`L(A,V)`$ approache the known relations for the exact solutions, in the particular cases like $`n=m`$ ($`5^{th}`$ row), $`n=m=2`$ ($`4^{th}`$ row), $`n=m=3`$. And $`n=3,m=2`$; $`n=2,m=3`$ in the $`6^{th}`$ row. These results can be used to predict the behavior of solutions for all values of $`n,m`$. Similar analysis can be done in the case of sine-Gordon equation, if we ask that velocity be proportional with $`L^2`$ ($`7^{th}`$ row). In this case we obtaining a transcendental equation in $`A`$, which is just the case of the sine-Gordon soliton. In the $`8^{th}`$ row, we present the cubic nonlinear Schrödinger equation (NLS) which has a soliton solution, too . This equation arises, for example, in nonlinear optics or in the polaron model in solid state physics, . In the general case of a NLS of order $`n`$ ($`9^{th}`$ row), when the general analytical solution is unknown, the method predicts a special $`L=L(A,V)`$ dependence shown in the fourth column and in Fig. 4. Contrary to third order NLS, where the dependence of $`L`$ with $`A`$ is monotonous for $`V=\pm A`$ ($`n=3`$ in Fig. 4), at higher order, the $`L(A)`$ function has a discontinuity in the first derivative. This wigle of the function (Fig. 4 for $`n=4`$) yields at a critical width, producing bifurcations in the solutions and scales. As a consequence, initial data close to this width can split into doublet (or even triplet for higher order NLS) solutions, with different amplitudes. Such phenomena have been put into evidence in several numerical experiments for quintic nonlinear equations . In the following, we present another example of applications of this qualitative approach, related to a new type of behavior of nonlinear systems. Traditional solitons move with constant speed on a rectilinear path (except for the roton, which has a circular trajectory with constant angular velocity). The speed is usually equal with the amplitude scaled with a constant. Higher solitons travel faster and there are no solitons at rest (zero speed asks for zero amplitude). They can travel in both directions with oposite signs for the amplitude. The situation is different in the case of compactons, which allow also stationary solutions. When linear and nonlinear disspersion occur simultaneously, like in the so called K(2,1,2) equation $$u_t+(u^2)_x+(u)_{xxx}+ϵ(u^2)_{xxx}=0,$$ where $`ϵ`$ is a control parameter, the similarity approach yields a dependence of the form $$L=\sqrt{(\pm A+ϵ)/(V\pm A)},$$ which still provides a constant width if $`V=\pm A+2ϵ`$. In this case the speed is proportional with the amplitude, but can change its sign even at non-zero amplitude. Solutions with larger amplitude than a critical one ($`A_{crit}=2ϵ`$) move to the right, solutions having the critical amplitude are at rest, and solutions smaller than the critical amplitude move to the left. This behavior was explored in . However, such a switching of the speed is not necessarily a feature of the nonlinear dispersion. A compacton of amplitude $`A`$ on the top of a KAK solution of amplitude $`\delta `$ $$u(x,t)=Acos^2\left(\frac{xVt}{4}\right)+\delta ,$$ (10) is still a solution of the K(2,2) equation, $`u_t+(u^2)_x+(u^2)_{xxx}=0`$, with the velocity given by $`V=\frac{3}{4}\left(2\delta +A\right)`$. For $`A=2\delta `$ the compacton becomes a stationary anticompacton, embedded in the moving, supporting KAK. Such an example is presented in Fig. 5 for a slow-scale time-dependent amplitude compacton. The induced oscillations in the amplitude transform into oscillations in the velocity. While not the topic of this paper, such a dynamic system has been analysed and it will be published soon. The key to such a conversion of oscillations is the coupling between the traditional nonlinear picture (convection-dispersion-diffusion) and the typical Schrödinger terms. A last application of this method, occurs if the KdV equation has an additional term depending on the square of the curvature $$u_t+uu_x+u_{xxx}+ϵ(u_{xx}^2)_x=0.$$ (11) This is the case for extremely sharp surfaces (surface waves in solids or granular materials) when the hydrodynamic surface pressure cannot be linearized in curvature. Such a new term yields a new type of localized solution fulfilling the relations $$L=\sqrt{\frac{4ϵA}{\pm \sqrt{18ϵA(A\pm V)}1}}.$$ If we look for a constant half-width solution (compacton of $`1/L=\alpha `$) we need a dependence of velocity of the form $`V=(1+\alpha ^2ϵ/8)A+1/8ϵA+\alpha /4`$. There are many new effects in this situation. The non-monoton dependence of the speed on $`A`$ introduces again bifurcations of a unique pulse in dublets and triplets. Also, there is a upper bound for the amplitude at some critical values of the width. Pulses narrower than this critical width drop to zero. Such bumps can exist in pairs of identical amplitude at different widths. They may be related with the recent observed ”oscillations” in granular materials, . As the examples presented in Table 1 proved the above method provides a reliable criterium for finding compact suported solutions. The reason this simple prescription works in so many cases follows from the advantages of wavelet analysis on localized solutions. We stress that this method has little to do with the traditional similarity (dimensional) analysis . In the latter case one obtains relations among powers of $`A,L`$ and $`V`$, not relations with numeric coefficients like those found in our method. ## 3 Compacton kink-antikink pairs and the multiresolution frame A common feature of all NPDE and of the finite differences equations is the existence of compact supported solutions. Compactons and discrete wavelets are typical examples. An interesting general conclusion can be obtained if we look at a one-dimensional model described by the most general NPDE dynamical equation $$_tu=𝒪(x,_x)u,$$ (12) where $`𝒪`$ is a nonlinear differential operator. By taking into account only traveling solutions, this NPDE reduces to a NODE in the coordinate $`\xi =xVt`$ for an arbitrary velocity $`V`$. If $`u(\xi )`$ is a compact solution it results that it is not unique for given initial compact data. If one chooses zero initial value for the solution and its derivatives up to the requested, in a certain point $`\xi _0`$ of the $`\xi `$ axis, these conditions can be fulfilled by any linear combination of disjoint translated versions of one particular solution, placed everywhere on the axis except $`\xi _0`$. Consequently, for such initial data, the solution is not unique. This result shows that the compact supported property of the initial data and the solution, implies its non-uniqueness. Since we can transform the NODE into a nonlinear differential system of order one $$\frac{d\stackrel{}{U}}{dx}=\stackrel{}{F}(\xi ,\stackrel{}{U}),\stackrel{}{U}=(u,_xu,\mathrm{}),$$ (13) we can apply the fundamental theorem of existence and uniqueness to solutions of eq.(13), for given initial data $`\stackrel{}{U}(\xi _0)=\stackrel{}{U}_0`$. If the function $`\stackrel{}{F}`$ in eq.(13) fulfills the Lipschitz condition (its relative variation is bounded) than, for any initial condition, the solution is unique, . Since any linear function is analytic and hence Lipschitz, we conclude that only nonlinear functions $`\stackrel{}{F}`$ allow the existence of compact supported solutions. Thus, a compact soliton implies non-uniqueness in the underlying NPDE, which implies non-Lipschitzian structure of the NPDE and hence the existence of nonlinear terms. In the following we investigate some compact solutions of the K(2,2) equation. The high stability against scattering of the K(2,2) compactons, or compacton generation from compact initial data, suggest they may play the role of a nonlinear local basis. We know form many numerical experiments, , that any positive compact initial data decomposes into finite series of compactons and anticompactons. This suggests that an intrinsic ingredient for a nonlinear basis could be the multiresolution structure of the solutions, similar with the structure of scaling functions in wavelet theory. The compactons given in eqs.(6,7) have constant half-width and hence describes a unique scale, which can cover all the space by integer translations. From the point of view of multi-resolution analysis, the K(2,2) equations acts like a $`L`$-band filter, allowing only a particular scale to emerge for any given set of initial condition. To each scale, from zero to infinity, we can associate a K(2,2) equation with different coefficients. However, the compacton solution is not the unique one with this property. For a given K(2,2) equation, we can thus extend the scale from $`L`$ to any larger scale. These more general compact supported solutions are still $`C_2(𝐑)`$ and are combinations of piece-wise constant and piece-wise $`\mathrm{cos}^2`$ functions. The simplest shape is given by a half-compacton prolonged with a constant level, that is a kink solution. The basis solution is a kink-antikink (KAK) compact supported combination, Fig. 1. Such kink-antikink pairs of different length, can be associated with other compactons, or KAK pairs, one on the top of the other $$\eta _{comp+KAK}(xVt;\lambda )=\{\begin{array}{cc}0\mathrm{}\hfill & \\ \frac{4V}{3}cos^2\left[\frac{xVt}{4}\right],2\pi xVt0\hfill & \\ \frac{4V}{3},0xVt\delta \hfill & \\ \frac{4V}{3}+\frac{4}{3}(V^{}2V)cos^2\left[\frac{xV^{}t}{4}\right],\delta xVt\delta +4\pi \hfill & \\ \frac{4V}{3},\delta +4\pi xVt\lambda \hfill & \\ \frac{4V}{3}cos^2\left[\frac{xVt\lambda }{4}\right],\lambda xVt\lambda +2\pi \hfill & \\ 0\mathrm{}\hfill & \end{array}$$ (14) where $`\delta <\lambda `$ characterizes the initial position (at $`t=0`$) of the top compacton, with respect to the flat part of the KAK solution. The amplitude $`4(V^{}2V)/3`$ of the compacton, and the amplitude $`4V/3`$ of the KAK, are related to their velocities $`V^{}`$ and $`V`$, respectively. The length of the flat part, $`\lambda `$, is arbitrary. A compound solution is not stable in time since the different elements travel with different velocities. The total height of the compacton is $`4(V^{}V)/3`$. Since the higher the amplitude is, the faster the structure travels, the top compacton moves faster than the KAK, and at a certain moment it passes the KAK. Because the area of the solution is conserving, such a compound structure decomposes into compactons and KAK pairs. Similar and even more complicated constructions can be imagined, with indefinite number of compactons and KAK’s, if one just fulfills the $`C_3`$ continuity condition for the square of the total structure. Such structures, defined at the initial moment can interpolate any function, playing a similar role with wavelets or spline bases. It has been also proved that the KAK solutions are stable, by using both a linear stability analysis and Lyapunov stability criteria, . For a given K(2,2) equation, the compacton solution, eq.(6) and in addition the family of KAK solutions, eq.(7) can be organized as a scaling functions system. They act like a low-pass filter in terms of space-time scales and give the opportunity to construct frames of functions from the wavelet model, . For the sake of simplicity we will renormalize the coefficients of the K(2,2) equation such that the support of the simple compacton is one. That is, we take $`\eta _c(x,t)=\eta _{kak}(\pi (xVt),0)`$ on the interval $`\|xVt\|`$ in $`[1/2,1/2]`$. We construct a multiresolution approximation of $`L^2(𝐑)`$, that is an increasing sequence of closed subspaces $`V_j`$, $`j𝐙`$, of $`L^2(𝐑)`$ with the following properties, 1. The $`V_j`$ subspaces are all disjoint and their union is dense in $`L^2(𝐑)`$. 2. For any function $`fL^2(𝐑)`$ and for any integer $`j`$ we have $`f(x)V_j`$ if and only if $`D^1f(x)V_{j1}`$ where $`D^1`$ is an operator that will be defined later. 3. For any function $`fL^2(𝐑)`$ and for any integer $`k`$, we have $`f(x)V_0`$ is equivalent with $`f(xk)V_0`$. 4. There is a function $`g(x)V_0`$ such that the sequence $`g(xk)`$ with $`k𝐙`$ is a Riesz basis of $`V_0`$. In the case of compact solutions of K(2,2) of unit length, we chose for the space $`V_0`$ that which is generated by all translation of $`\eta _c`$ with any integer $`k`$. The subspaces $`V_j`$ for $`j0`$ are generated by all integer translations of the compressed version of this function, namely, by $`\eta _{kak}(2^j\pi (xVt),0)`$. The subspaces $`V_j`$ for $`j0`$ are generated by all integer translations of the KAK solution of length $`\lambda 2^j1`$. For example, $`V_1`$ is generated by $`\eta _{kak}(\pi (x2^jVt),0)`$. The spaces $`V_j,j0`$ are all solutions of K(2,2); the others are not. The function $`g(x)`$ is taken to be $`\eta _{kak}(\pi (xVt),0)`$. It is not difficult to prove that these definitions fulfill restrictions one, three, and four. As for the second criterion, we define the action of the operator $`D^1f(x)=f(2x)`$ if $`f(x)V_j`$ with a $`j`$ positive integer, and $`D^1\eta _{kak}(\pi 2^j(x2^jVt),2^j1)=\eta _{kak}(\pi 2^j(x2^{j+1}Vt),2^{j+1}1)`$ for negative $`j`$. In conclusion, we construct a frame of functions made of contractions of compactons and sequences of KAK solutions. We can write the corresponding two-scale equation which connects the subspaces (the equivalent of eq.(15)), $$\eta _{kak}(\pi (xVt),1)=\eta _{kak}(\pi (xVt),0)+\eta _{kak}(\pi (xVt1),0).$$ (15) We will denote generically by $`\eta _{k,j}`$ the elements of this frame, that is $$\eta _{k,j}(x)=\eta _{kak}(\pi (x2^jVtk),2^j1)\|_{t=0},$$ where $`t=0`$ means that we neglect the time evolution, but the amplitude is still amplified with a factor of $`2^j`$, in virtue of relation $`\eta _{max}=4V/3`$. In the following, we can expand any initial data for the K(2,2) equation in this basis. $$u_0(x)=\underset{k}{}\underset{j}{}C_{k,j}\eta _{k,j}(x).$$ (16) We notice that the following equality holds for $`j^{}`$,$`j`$ $$\eta _{k,j}\eta _{k^{},j^{}}\{\begin{array}{cc}0\hfill & k^{}=k2^{j^{}j},\mathrm{},(k+1)2^{j^{}j}1\hfill \\ =0\hfill & \text{otherwise.}\hfill \end{array}$$ (17) After some rather elaborate algebraic calculations and by using eq.(21), we show that the square of this function (since the equations is nonlinear and of order two) will be given by $$u^2(x)=\underset{k,j}{}\underset{j^{}j}{}\underset{k^{}I}{}C_{k,j}C_{k^{},j^{}}$$ $$\times \left(\underset{i_1=0}{\overset{1}{}}\underset{i_2=0}{\overset{1}{}}\mathrm{}\underset{i_{j^{}j}=0}{\overset{1}{}}\eta _{\sigma (i_1,i_2,\mathrm{},i_{j^{}j}),j^{}}\right)\eta _{k^{},j^{}},$$ (18) where $`I`$ is the range of $`k^{}`$ described in the first line of eq. (21), and $$\sigma (i_1,i_2,\mathrm{},i_{j^{}j})=\underset{l=1}{\overset{j^{}j}{}}i_l2^{j^{}j^{}l+\frac{(j^{}j)(j^{}j+1)l(l+1)}{2}}$$ $$+k2^{(j^{}j)j+\frac{(j^{}j)(j^{}j+1)}{2}}.$$ From eq.(21) we notice that in eq.(22)the unique nonzero terms are those for which $`\sigma (i_1,i_2,\mathrm{},i_{j^{}j})=k^{}`$ with $`k^{}I`$. This result express the following simple fact. The initial data is expanded in different scales and different translations. The translations are mutually orthogonal so they do not give a contribution to the square. When we have to multiply two different scales in the expression of the square, we reduce the wider scale in terms of linear combination of the narrower one by using the two-scale equation, eq.(19). This is what eq.(22) expresses. Out of all the terms in such a product only approximately $`(2^j1)/(2^j^{}1)2^{j^{}j}`$ give non-zero contributions. In other words this number is given by the number of solutions of equation $`\sigma (i_1,i_2,\mathrm{},i_{j^{}j})=k^{}`$, with $`k^{}I`$. This is the primary advantage of treating nonlinear problems with a basis that has a scale criterion. Another advantage is that all the function in the basis are actually contractions or dilations, and translations of only two basic ones. ## 4 Comments and conclusions In the present paper we introduce new applications for wavelets, in the field of the study of localized solutions of nonlinear differential equations. The existence of compactons and discrete wavelets underlines a common feature of NPDE and finite differences equations, that is the existence of compact supported solutions. We propose a new similarity formalism for the qualitative analysis and clasiffication of soliton solutions, without the need of solving the corresponding NPDE. Also, we proved that starting from any unique soliton solution of a NPDE, we can construct a frame of solutions organized under a multiresolution criterium. This approach provides the possibility of constructing nonlinear basis for NPDE. We show that frames of self-similar functions are related with solitons with compact support. In addition, we notice the evidence that compactons fulfil both characteristics of solitons and wavelets, suggesting possible new applications. Such unifying direction between nonlinearity and self-similarity, can bring new applications of wavelets in cluster formation, at any scale, from supernovae through fluid dynamics to atomic and nuclear systems. The similarity approach can be applied with succes to the physics of droplets, bubbles, traveling patterns, fragmentation, fission and inertial fusion. Supported by the U.S. National Science Foundation through a regular grant, No. 9970769, and a Cooperative Agreement, No. EPS-9720652, that includes matching from the Louisiana Board of Regents Support Fund. Figure Captions * Fig. 1 A compacton and a kink-antikink pair solution (KAK) of the equation K(2,2), both having the same amplitude, and hence velocity V. To the right, there is a smaller compacton on the top of KAK. The upper compacton has higher speed, V’. * Fig. 2 A finite series of K(2,2) compactons emerging from initial compact data, with the width larger than the compacton width. * Figs. 3 The half-width $`L`$ versus velocity $`V`$ for the K(2,2) equation, for different amplitudes A. Figure 3a shows widths larger than $`L_{compacton}=3/4`$, and Fig. 3b shows narrower widths, $`L<3/4`$. Amplitude increases from left to right, in the range 0.01-0.85. * Fig. 4 The half-width $`L`$ plot versus amplitude $`V`$, for the third (n=3) and forth (n=4) order NLS equation, in two $`V=\pm A`$ cases. From the figure one notes that the quartic NLS equation yields bifurcations in the solutions. * Fig. 5 The solution of the mixed linear plus nonlinear-dispersion K(2,2) equation, in the case of a solution with a slow oscillating shape.
warning/0003/hep-ph0003065.html	ar5iv	text	# Top quark, heavy fermions and the composite Higgs boson ## Abstract We study the properties of heavy fermions in the vector-like representation of the electro-weak gauge group $`SU(2)_W\times U(1)_Y`$ with Yukawa couplings to the standard model Higgs boson. Applying the renormalization group analysis, we discuss the effects of heavy fermions to the vacuum stability bound and the triviality bound on the mass of the Higgs boson. We also discuss the interesting possibility that the Higgs particle is composed of the top quark and heavy fermions. The bound on the composite Higgs mass is estimated using the method of Bardeen, Hill and Lindner , $`150`$GeV$`m_H`$ $`450`$GeV. Enormous efforts have been made in searching for physics beyond the standard model but up to now a crucial, direct experimental indication is still illusive. One of the most important motivation to study the property of heavy fermions above the energy scale accessible by current accelerators is to look for extra building blocks of nature beyond the three families of the standard model. For this purpose it may be adequate to study fermions in vector-like representations of the electro-weak gauge group with a large bare mass term, rather than the conventional chiral fermions. The main reason for this is from the strong experimental constraints on the S parameter . While experiments favor a negative value of S , a standard chiral doublet of heavy fermions (degenerate in mass) contributes to the S parameter as $`1/6\pi `$. On the contrary, for fermions in the vector-like representation of the electro-weak gauge group, a large bare fermion mass M completely changes the low energy properties of the heavy fermions. As a consequence of the decoupling theorem, heavy fermions’ contribution to the oblique corrections of the standard processes are suppressed by $`\frac{1}{M^2}`$. Especially, their contribution to the S parameter is still positive definite but much smaller in magnitude than the ordinary chiral fermions. Furthermore the heavy fermion contributes to the vacuum expectation value of electroweak symmetry breaking as , $$\delta (f_\pi ^2)=\delta v^2\frac{m^2N_c}{2\pi ^2}(\mathrm{log}\frac{\mathrm{\Lambda }^2}{M^2}),$$ (1) where where $`m`$ is the mass generated by the Yukawa coupling and $`\mathrm{\Lambda }`$ is the cutoff scale of the effective theory. It is interesting to compare the above expression to that of the pion decay constant obtained in the QCD effective action approach , $`f_\pi ^2=\frac{N_c}{4\pi ^2}M_Q^2\mathrm{ln}(\frac{\mathrm{\Lambda }_{QCD}^2}{M_Q^2})`$, where $`M_Q`$ is the constitute quark mass which is similar to $`m`$ in our present discussion. We notice that if in the above Eq. (1) $`mO(v)`$ then several of these heavy fermions would be enough to induce the electroweak symmetry breaking. Therefore if there is a strong attractive forces in the appropriate channel to cause the heavy fermion condensation then they may place the role similar to techniquarks in the technicolor model. This way of dynamical electroweak symmetry breaking, if possible, is remarkable. Contrary to the technicolor model, it avoids the dangerous low energy consequences which may contradict experiments. Also it can be demonstrated that the composite Higgs boson’s mass is proportional to the dynamically generated fermion mass and completely decouples from the bare one, even though the Higgs particle is “composed of” the heavy fermions. This is a consequence of symmetry and be model independent, at least in a system with second order phase transition. Heavy fermions may have many other interesting role in physics beyond the standard model either. For example, they may be responsible for a dynamical generation of light fermion mass matrix ; they appear in the “vector–like extension” of the standard model; they are natural consequences of many grand unification models, and of the super-symmetric preon model . Therefore it is important to investigate the fundamental properties of the heavy vector-like fermions thoroughly. There have been continuous interests in understanding the structure of the standard model at high energies, even up to Planck scale (see for example, and the most recent review which contains many materials, Ref. ). A powerful tool is to use the renormalization group equations to trace the evolution of the coupling constant of the $`\lambda \varphi ^4`$ self-interaction of the Higgs particle. Assuming the standard model remains valid up to certain scale $`\mathrm{\Lambda }`$, an upper bound (the triviality bound, obtained by requiring $`\lambda `$ not to blow up below $`\mathrm{\Lambda }`$) of the Higgs boson mass, $`m_H`$, can be obtained. Meanwhile, requiring the stability of the electro-weak vacuum, we can also obtain a lower bound on $`m_h`$. For the later purpose, in principle one needs to consider the renormalization group improved effective potential and require it be bounded from below. But in practice this turns out to be equivalent to the requirement that the Higgs self-interaction coupling constant $`\lambda `$ does not become negative, below the given scale (see and ref. therein). It is remarkable that for the given experimental value of the top quark mass (here we use $`m_t=174`$GeV), there is an allowed range for the Higgs boson mass, $`130\text{GeV}`$$`m_H`$$`200\text{GeV}`$ , for which the standard model may remain valid up to Planck scale. In this paper we devote to study heavy fermions’ influence to the vacuum stability bound and the triviality bound on the Higgs boson mass. Furthermore, assuming that the Higgs boson is a composite particle, we use the method developed in Ref. to estimate the range of the Higgs boson’s mass<sup></sup><sup></sup>* This paper replaces and is an extension of Ref. .. We find that the top quark also place an important role in the compositeness picture and the composite Higgs boson can be viewed as a mixture of $`\overline{t}t`$ pair and heavy fermion pair. The larger the hierarchy is the more top quark content the composite Higgs boson contains and vise–versa. We start with the following general Lagrangian for heavy fermions, $`=\overline{Q}(i\mathrm{\Delta }/_dM)Q+\overline{U}(i\mathrm{\Delta }/_sM)U+\overline{D}(i\mathrm{\Delta }/_sM)D+g_d\overline{Q}_L\varphi D_R+g_u\overline{Q}_L\stackrel{~}{\varphi }U_R`$ $$+g_d^{}\overline{Q}_R\varphi D_L+g_u^{}\overline{Q}_R\stackrel{~}{\varphi }U_L+h.c..$$ (2) In above $`Q`$ is the $`SU(2)_W`$ doublet and $`U`$ and $`D`$ are singlets with weak hypercharge $`Y_Q`$, $`Y_U`$ and $`Y_D`$, respectively (with the selection rule $`Y_UY_Q=Y_QY_D=Y_\varphi `$). We assume they participate in strong interactions and are in fundamental representations of $`SU(3)_C`$. The subscript $`d`$ ($`s`$) in the covariant derivatives denotes that the corresponding fermion is a $`\text{SU(2)}_W`$ doublet (singlet) and $`\varphi `$ denotes the standard Higgs doublet. We further expect the Yukawa couplings to be of order 1. For simplicity we take all the bare fermion masses to be equal. Also we do not discuss the mixing between heavy fermions and the ordinary fermions here. As is well known, because of the negative sign, fermions turn to destabilize the vacuum. After including heavy fermions the structure of our world changes drastically at high scales, even though vector-like fermions are essentially decoupling below their threshold. At scales much higher than the threshold whether the fermion field is chiral or vector-like does not make any qualitative difference. The only thing matters is the number of independent Yukawa couplings and their strength. The relevant one loop RGEs are listed as below Due to a careless mistake, the Yukawa coupling RGEs given in Ref. contain an error. The top quark effects were not considered correctly., $$16\pi ^2\frac{d\lambda }{dt}=24\lambda ^2+12\lambda A6A^{}(9g_2^2+3g_1^2)\lambda +\frac{9}{8}g_2^4+\frac{3}{4}g_2^2g_1^2+\frac{3}{8}g_1^4,$$ (3) $$16\pi ^2\frac{dg_u}{dt}=\{\frac{3}{2}(g_ug_u^{}g_dg_d^{})+3A8g_s^2\frac{9}{4}g_2^23(Y_Q^2+Y_U^2)g_1^2\}g_u,$$ (4) $$16\pi ^2\frac{dg_d}{dt}=\{\frac{3}{2}(g_dg_d^{}g_ug_u^{})+3A8g_s^2\frac{9}{4}g_2^23(Y_Q^2+Y_D^2)g_1^2\}g_d,$$ (5) $$16\pi ^2\frac{dg_t}{dt}=\{\frac{3}{2}g_t^2+3A8g_s^2\frac{9}{4}g_2^2\frac{17}{12}g_1^2\}g_t,$$ (6) where, $$A=\text{tr}\{g_ug_u^{}+g_dg_d^{}+g_u^{}(g_u^{})^{}+g_d^{}(g_d^{})^{}\}+g_t^2,$$ (7) $$A^{}=\text{tr}\{(g_ug_u^{})^2+(g_dg_d^{})^2+(g_u^{}(g_u^{})^{})^2+(g_d^{}(g_d^{})^{})^2\}+g_t^4,$$ (8) and $$16\pi ^2\frac{dg_s}{dt}=(7+\frac{2}{3}(2N_Q+N_U+N_D)\theta )g_s^3,$$ (9) $$16\pi ^2\frac{dg_2}{dt}=(\frac{19}{6}+2N_Q\theta )g_2^3,$$ (10) $$16\pi ^2\frac{dg_1}{dt}=(\frac{41}{6}+4(2N_QY_Q^2+N_UY_U^2+N_DY_D^2)\theta )g_1^3,$$ (11) where the trace doesn’t sum over color space and $`g_u^{}`$ and $`g_d^{}`$ obey similar equations. In general these Yukawa couplings can be matrices in the flavor space if there are many heavy fermions, and $`g_t`$ is the Yukawa coupling of the top quark ($`g_t=\sqrt{2}m_t/v`$). The symbols $`N_Q`$, $`N_U`$ and $`N_D`$ refer to the number of Q, U and D type of quarks, respectively. We use a simple step function $`\theta =\theta (tlog(M/M_z))`$ to model the heavy fermion threshold effects. All the Yukawa couplings in above renormalization group equations are understood as multiplied by $`\theta `$. Applications using two loop RGEs in the standard model case and beyond was considered in Ref. and it was found that the two loop effects are very small below Planck scale. In the following qualitative discussion, we set $`Y_Q=1/6`$, $`Y_U=2/3`$ and $`Y_D=1/3`$. For simplicity we take $`N_Q=N_U=N_D`$ ($`N`$) and all the Yukawa couplings (after the diagonalization of the coupling matrices) in the initial boundary conditions being identical The ‘up’ and ‘down’ type quarks evolve differently because of different $`U(1)_Y`$ charge, however the isospin splitting is very small for the standard values of the hypercharge.. In fig. 1 we plot the vacuum stability bound and the triviality bound on the Higgs mass as a function of the scale $`\mathrm{\Lambda }`$ for some typical values of the parameters of the heavy fermions. We see that the inclusion of heavy fermions drastically change the Standard model structure at high energies even though they decouple from the low energy world. They tighten the bound on the mass of the Higgs boson as a function of the cutoff scale $`\mathrm{\Lambda }`$. Notice that (in terms of one loop renormalization equations) the upper line (triviality bound) and the lower line (vacuum stability bound) never meet each other. Because the upper line is drawn by requiring $`\lambda `$ not to blow up and the lower line is drawn by requiring $`\lambda 0`$. Between them is the ultra-violet unstable fixed point of $`\lambda `$, so the two lines get close to each other rapidly. We now study the interesting possibility of considering the Higgs particle as a composite object of the heavy vector-like fermions. Applying the above renormalization group analysis to the composite model leads to some interesting results which we present below. We follow the method proposed by Bardeen, Hill and Lindner (BHL) originally developed for the top quark condensate model. The basic idea of the BHL method is the following: Using the collective field method the four–fermi interaction Lagrangian can be rewritten into an effective Higgs–Yukawa interaction Lagrangian at the cutoff scale $`\mathrm{\Lambda }`$. The effective Yukawa interaction Lagrangian is identical to the standard model at the cutoff scale $`\mathrm{\Lambda }`$, but with vanishing wave function renormalization constant of the Higgs field ($`Z_H=0`$) and vanishing Higgs self-coupling ($`\lambda =0`$). Below $`\mathrm{\Lambda }`$ the model is equivalent to the standard model and therefore the coupling constants of the effective theory run according to the standard model renormalization group equations. However the vanishing of $`Z_H`$ at the scale $`\mu =\mathrm{\Lambda }`$ leads to the following boundary conditions of the renormalization group equations: $$g_Y^r\mathrm{},\lambda ^r/(g_Y^r)^40,$$ (12) where $`\lambda ^r`$ and $`g_Y^r`$ are the renormalized Higgs self-coupling and Yukawa coupling, respectively. With the renormalization group equations and boundary conditions, one can predict the mass of the Higgs boson and the fermion mass (or the Yukawa couplings) at the infra-red fixed point. In the present case, of course, the “standard model” often refers to the standard model plus heavy fermions and the “infra-red fixed point” value of $`g_Y`$ refers to its value at the threshold. The minimal top quark condensate model has already been ruled out by experiments. In order to generate the electroweak symmetry breaking scale $`v`$, the top quark mass is required to be at least as large as 218 GeV (corresponding to $`\mathrm{\Lambda }=10^{19}`$ GeV, i.e., Planck scale). The experimental value of the top quark mass indicates that the top quark Yukawa coupling does not diverge up to Planck scale in the standard model and therefore does not meet the compositeness condition of BHL. This can be clearly seen from fig. 2. However, in the present model, since there is no strict experimental constraint on the heavy fermions, the compositeness condition is easily and naturally achievable, that $`g_t`$ blows up below Planck scale with the aid of the heavy fermions. From Eqs. (4), (6) we see that the evolution of the Yukawa couplings are correlated to each other and one ‘blows up’ leads the another to blow up too. When both the top quark and heavy fermions are involved, the situation is more complicated than the simple top condensate model. Running the RGEs down from certain scale, one must take good care of $`g_t`$ to ensure that it reaches the experimental value at the infra-red fixed point. This means that a certain fine-tuning is needed on the initial boundary conditions of the Yukawa coupling RGEs. The composite Higgs boson is now a mixture of $`\overline{t}t`$ pairs and the heavy quark pairs. Fig. 3 and fig. 4 show two typical examples of such a situation. In the situation of fig. 3 the Higgs particle is mainly composed of heavy fermions while in fig. 4 the top quark becomes the major component. Notice that for a given ratio of $`g_Y/g_t`$ in the compositeness boundary condition (for fixed M and N), the composite scale $`\mathrm{\Lambda }`$ is no longer free, rather it is determined by $`g_t^{exp}`$. In fig. 5 we plot the composite Higgs particle’s mass<sup>§</sup><sup>§</sup>§ The Higgs mass in these figures is the renormalized mass at $`\mu =M_Z`$. The renormalized mass is close to the pole mass of the Higgs boson. as a function of the composite scale, $`\mathrm{\Lambda }`$. We chose $`N3`$ to avoid the problem with the non-asymptotic freedom of $`g_s`$. From fig. 5 we see that the allowed range for the Higgs mass is rather narrow against the wide range of the cutoff scale, the bare fermion mass and the number of heavy fermions, except when the heavy fermion bare mass $`M`$ is close to the cutoff $`\mathrm{\Lambda }`$. A lower bound on the Higgs mass can be obtained: $`m_H150`$ GeV. When $`M`$ is getting close to the cutoff scale our results become unstable and are sensitive to the input numerical values of the boundary conditions. In such a situation the scale is not large enough for the couplings to reach the infra-red stable point. It is estimated that the Higgs mass will not exceed 450 GeV, otherwise the whole mechanism become unnatural (in the sense that the Yukawa coupling constant at electroweak scale also becomes substantially larger than 1). In fig. 6 we plot a typical example of the Higgs mass for a given cutoff scale $`\mathrm{\Lambda }_c`$ and $`N`$. We also plot the triviality bound and the vacuum stability bound using the value of the Yukawa coupling constant at the infrared fixed-point, which is determined uniquely by the parameters $`M`$, $`\mathrm{\Lambda }_c`$ and $`N`$ in the compositeness picture, as the initial boundary condition. It is very interesting to notice that $`m_H`$ and $`\mathrm{\Lambda }`$ take the values where the curves of triviality bound and vacuum stability bound (practically) meet each other. This is the unique feature of BHL compositeness picture. The reason behind this is very simple: The infra-red attractive fixed point corresponds to the ultra-violet unstable fixed point. In the sense of Ref. , this picture can be disturbed. However in most cases the infra-red–ultra-violet fixed point structure is influential and rather stable against perturbation. In above we presented an analysis on the properties of heavy fermions in vector-like representations of the standard model gauge group. We pointed out earlier that if they can place the role to break the electro-weak symmetry dynamically the theory has some distinguishable properties: the low energy theory is asymptotically renormalizable and returns to the standard model. From the above RG analysis we realize that the top quark also places an important role in the dynamical symmetry breaking scenario and our model can be viewed as a natural generalization to the top condensate model of BHL. We found that the composite Higgs boson’s mass ranges from 150GeV to 450GeV, and the lighter the Higgs boson is the more top quark content it contains, and vice versa. Our prediction to the mass of the Higgs boson will be testable by LHC and the model will be ruled out if $`m_H`$ is found to be below 150GeV. Acknowledgment: The work of H.Z. is supported in part by the National Natural Science Foundation of China under grant No. 19775005.
warning/0003/hep-ph0003168.html	ar5iv	text	# Charged Particle Ratio Fluctuation as a Signal for QGP ## Abstract In this letter we argue that the event-by-event fluctuations of the ratio of the positively charged and the negatively charged pions provides a signal of quark-gluon plasma. The fact that quarks carry fractional charges is ultimately responsible for this distinct signal. It is of great importance that we have a clear signal of the long-sought quark-gluon plasma (QGP) not only for the experiments at RHIC but also for theoretical reasons. At stake is our fundamental understanding of strong interactions as well as understanding of the state of matter in the very early universe. Proposed signals of this new state of matter abound in literature one of the most studied being the $`J/\psi `$ suppression . In this paper, we propose the event-by-event $`h^+/h^{}`$ fluctuations as a distinct signal of QGP formation. We would also like to stress that this observable is something that can be and already has been calculated on a lattice. The idea is very simple and is reminiscent of the original detection of color in $`e^+e^{}`$ experiment where one measures $`R_{e^+e^{}}{\displaystyle \frac{e^+e^{}\text{Hadrons}}{e^+e^{}\mu ^+\mu ^{}}}=N_c{\displaystyle \underset{q}{}}Q_q^2`$ (1) Here $`Q_q`$ is the charge of each flavor and $`N_c`$ is the number of colors. Note that if the fundamental degrees of freedom were hadrons, $`R_{e^+e^{}}`$ would be very different from this simple counting. We would like to establish that the event-by-event $`h^+/h^{}`$ fluctuations can similarly determine whether the underlying degrees of freedom are quarks and gluons or hadrons. The point is that in the QGP phase, the unit of charge is $`1/3`$ while in the hadronic phase, the unit of charge is 1. The net charge, of course does not depend on such subtleties. However, the fluctuation in the net charge depends on the squares of the charges and hence strongly depend on which phase it originates from. Measuring the charge fluctuation itself, however, is plagued by systematic uncertainties such as volume fluctuations due to the impact parameter variation. In a previous letter, we showed that the multiplicity ratio fluctuation is only sensitive to the density fluctuations and not to the volume fluctuations. The task for us is then to find a suitable ratio whose fluctuation is easy to measure and simply related to the net charge fluctuation. The obvious candidate is the ratio $`F=Q/N_{\mathrm{ch}}`$ where $`Q=N_+N_{}`$ (2) is the net charge and $`N_{\mathrm{ch}}=N_++N_{}`$ (3) is the charge multiplicity. Here $`N_\pm `$ denote the positive and negative multiplicities. Instead of using $`F`$, however, in this paper we propose to use the charge ratio $`R=N_+/N_{}`$. The advantages of using $`R`$ over $`F`$ are that although trivially related, $`R`$ is more fundamental to experiments and the signal is about 4 times amplified in $`R`$ as we show below. To relate $`R`$ with $`F`$, we first rewrite the charge ratio as $`R={\displaystyle \frac{N_+}{N_{}}}={\displaystyle \frac{1+F}{1F}}`$ (4) When $`N_{\mathrm{ch}}Q`$ we can safely say $`\left\|F\right\|1`$. Expanding in terms of $`F`$ yields $`R1+2F+2F^2`$ (5) Defining $`\delta x=xx`$ for any fluctuating quantity $`x`$, it is easy to show $`\delta R^2=R^2R^24\delta F^2`$ (6) where $`\mathrm{}`$ denotes the average over all events. Let us now consider $`\delta F^2`$ more closely. In a previous letter (see also and the upcoming paper ), we showed that a ratio fluctuation can be expressed as $`\delta F^2={\displaystyle \frac{Q^2}{N_{\mathrm{ch}}^2}}\left({\displaystyle \frac{\delta Q}{Q}}{\displaystyle \frac{\delta N_{\mathrm{ch}}}{N_{\mathrm{ch}}}}\right)^2`$ (7) We then showed that when the average ratio is very much different than 1, the fluctuation is driven mainly by the fluctuation in the smaller quantity (for instance $`K/\pi `$ fluctuation is driven by $`K`$ fluctuation). At RHIC we expect $`Q/N_{\mathrm{ch}}5\%`$. Hence the fluctuation in $`F`$ is totally dominated by the fluctuation in $`Q`$ so that $`\delta F^2{\displaystyle \frac{\delta Q^2}{N_{\mathrm{ch}}^2}}`$ (8) If we can detect all charged particles from a heavy ion collision, the net charge $`Q`$ is a fixed quantity and hence will not fluctuate. This implies that $`\delta F^2`$ is very small with a $`4\pi `$ coverage. However, no detector can catch all charged particles. Our study shows that for a realistic detector acceptance, using the grand canonical ensemble is acceptable and that is what we assume here. The corrections to this approximation have been worked out and will be reported in the upcoming paper. (Also see Finite Acceptance Correction at the end of this letter.) Our main observable is then, $`N_{\mathrm{ch}}\delta R^2=4N_{\mathrm{ch}}\delta F^2=4{\displaystyle \frac{\delta Q^2}{N_{\mathrm{ch}}}}`$ (9) to the leading order in the fluctuations and $`1/N_{\mathrm{ch}}`$. So far, we have only considered statistics of the ratio fluctuations. Physics lies in how the charge fluctuation is expressed in terms of the fluctuations in the fundamental degrees of freedom. For simplicity, let us consider a pion gas and a QGP consisting of $`u`$ and $`d`$ quarks and gluons. Our main conclusion does not depend on this simplifying assumption. We will briefly consider the size of the corrections towards the end of the paper. In a pion gas, the fundamental degrees of freedom are of course pions. Hence, $`N_{\mathrm{ch}}=N_{\pi ^+}+N_\pi ^{}`$ and $`\delta Q=\delta N_{\pi ^+}\delta N_\pi ^{}`$ (10) Using thermal distributions and disregarding correlations, we get $`\delta Q^2=\delta N_+^2+\delta N_{}^2=w_\pi N_{\mathrm{ch}}`$ (11) where $`w_\pi \delta N_\pi ^2/N_\pi `$ (12) is slightly bigger than 1 . Hence for a pion gas, $`D_{\mathrm{had}}N_{\mathrm{ch}}\delta R^2\|_{\mathrm{hadron}}4.`$ (13) For a quark-gluon plasma, $`\delta Q=Q_u\delta \left(N_uN_{\overline{u}}\right)+Q_d\delta \left(N_dN_{\overline{d}}\right)`$ (14) where $`Q_q`$ is the charges of the quarks and $`N_q`$ is the number of quarks. The fluctuations $`\delta N_{u,d}^2`$ are measured on lattice and we will shortly get back to the results. For now, let us consider a thermalized gas of non-interacting quarks and gluons to get a physical baseline. Thermal distributions and no correlations yield $`\delta Q^2=Q_u^2w_uN_{u+\overline{u}}+Q_d^2w_dN_{d+\overline{d}}`$ (15) where $`N_{q+\overline{q}}`$ from here on denotes the number of quarks and anti-quarks. The constant $`w_q\delta N_q^2/N_q`$ (16) is slightly smaller than 1 due to the fermionic nature of the quarks. Relating the final charged particle multiplicity $`N_{\mathrm{ch}}`$ to the number of primordial quarks and gluons is not as simple. To make an estimate, we assume that the entropy is conserved and that all the particles involved are massless, in thermal equilibrium and non-interacting. For such particles, the following relation between the entropy density and the particle number density holds: $`\sigma _B=3.6n_B`$ (17) and $`\sigma _F=4.2n_F`$ (18) where the subscript $`B,F`$ signifies the particle types. The total entropy of a quark-gluon gas in a given volume $`V_{qg}`$ is $`S=V_{qg}\sigma _{qg}=3.6N_g+4.2\left(N_{u+\overline{u}}+N_{d+\overline{d}}\right)`$ (19) where $`N_g`$ is the number of gluons inside the volume and $`N_{q+\overline{q}}`$ is the number of quarks and anti-quarks inside the same volume. As the volume expands and cools, eventually the quarks and gluons are converted to pions. Since entropy is conserved, the number of pions coming from these quarks and gluons must be given by $`N_\pi ={\displaystyle \frac{S}{3.6}}=N_g+{\displaystyle \frac{4.2}{3.6}}\left(N_{u+\overline{u}}+N_{d+\overline{d}}\right)`$ (20) The charged multiplicity is $`2/3`$ of $`N_\pi `$ due to isospin symmetry: $`N_{\mathrm{ch}}={\displaystyle \frac{2}{3}}\left(N_g+1.2N_{u+\overline{u}}+1.2N_{d+\overline{d}}\right)`$ (21) Then for massless non-interacting quarks and gluons, $`D_{\mathrm{QGP}}N_{\mathrm{ch}}\delta R^2\|_{\mathrm{QGP}}0.75`$ (22) from Eq.(9) and using Eqs.(15) and (21). The value of $`D_{\mathrm{QGP}}`$ is more than a factor of 5 smaller than the value of $`D_{\mathrm{had}}`$ in Eq.(13)! This is an unmistakable signal of QGP formation from such a simple measurement. We now would like to stress that $`\delta Q^2/N_{\mathrm{ch}}`$ is already calculated on lattice and hence one does not have to rely on the above thermal model calculation. In Ref., Gottlieb et al report their calculation of the quark number susceptibility and the entropy density with 2 flavors of dynamic quarks. These two quantities are directly related to the net charge fluctuation and the charged multiplicity in the following way. From the definition of the charge susceptibility $`\chi _q`$, it is clear that $`\delta Q^2=V_{qg}T\chi _q`$ (23) where $`V_{qg}`$ is the volume of the quark-gluon plasma at the hadronization and $`T`$ is the temperature. Gottlieb et al calculated the quark number density susceptibilities $`T\chi _S=V_{qg}\left(\delta n_u+\delta n_d\right)^2`$ (24) and $`T\chi _{NS}=V_{qg}\left(\delta n_u\delta n_d\right)^2`$ (25) and found that at high temperature both are very close to the non-interacting thermal gas limit $`\chi _S\chi _{NS}2T^2`$ (26) From this result, one can first of all infer that $`u`$ and $`d`$ quark densities are uncorrelated $`\delta N_u\delta N_d0`$ (27) and (28) $`\delta N_u^2\delta N_d^2.`$ These results imply that the charge fluctuation at high temperature follows that of the thermal fermion gas $`\delta Q^2={\displaystyle \frac{4}{9}}\delta N_u^2+{\displaystyle \frac{1}{9}}\delta N_d^2={\displaystyle \frac{5}{9}}V_{qg}T^3`$ (29) or (30) $`T\chi _q={\displaystyle \frac{5}{9}}T^3`$ $`.`$ For the charged multiplicity, we assume that the relations Eqs. (2021) still hold. Equating the entropy of the final pions with that of the primordial quarks and gluons, one obtains $`N_{\mathrm{ch}}={\displaystyle \frac{1}{5.4}}V_{qg}\sigma _{qg}.`$ (31) Ref. reports that the gluon entropy density $`\sigma _g`$ is almost the same as the non-interacting thermal bosons, but the quark entropy density $`\sigma _{u+d}`$ is about one half of that of the non-interacting thermal fermions. Hence, the total entropy from the lattice calculation is $`\sigma _{qg}`$ $`=`$ $`\sigma _g+\sigma _{u+d}`$ (32) $`=`$ $`16\times 3.6\times f_g+24\times 4.2\times \alpha f_q`$ (33) $``$ $`12T^3`$ (34) where $`\alpha 1/2`$ and $`f_{q,g}`$ is the average density per degree of freedom. Using Eqs.(3034), the lattice calculation gives $`D_{\mathrm{lat}}N_{\mathrm{ch}}\delta R^2\|_{\mathrm{lattice}}1`$ (35) which is still 4 times smaller than the pion gas result. This is the main result of this paper. The difference between the pion gas result and the QGP result is distinct enough one should easily see it in the first few days of data collecting at RHIC. We are of course aware that lattice result is not yet exact. We hope that this discussion will actually stimulate more sophisticated lattice calculations of quark number susceptibilities. We also note that this is an opportunity to test what is calculated on lattice in a direct observation. The picture obtained above holds if the following two conditions are met: (i) The detected phase space is a small sub-system of the whole. (ii) The original quarks and gluons stay in the system during or after the hadronization. Both conditions can be met if the rapidity intervals are such that $`y_{\mathrm{total}}y_{\mathrm{accept}}1.`$ (36) Here, $`y_{\mathrm{total}}`$ is the rapidity range allowed by the energetics of the collisions and $`y_{\mathrm{accept}}`$ is the acceptance interval of a given detector. The first of these conditions is needed to ensure that the rest of the system acts as a reservoir and the second condition ensures that the charge diffusion in the rapidity space during and after hadronization is negligible. In real life, of course, Eq. (36) is satisfied in varying degrees. For instance, the STAR at RHIC has $`y_{\mathrm{total}}10`$, $`y_{\mathrm{accept}}3`$ and hence corrections should be taken into account. We would like now to discuss Caveats and corrections due to these and other effects. Hadronization : First, charge conservation fixes the net charge once it is set in the QGP phase. Subsequent hadronization cannot change the net charge. Second, the entropy can only increase during the hadronization. Hence, our estimate of $`\delta Q^2/N_{\mathrm{ch}}`$ from QGP should not only survive the hadronization but also could even become smaller, thus strengthening the signal. Phase Space Cuts : A possible issue of having a rapidity cut is that this implies a box in the momentum space while the argument presented so far dealt with a box in the coordinate space. Within the Bjorken scenario, this is of course not a problem. However this is not a problem even without the Bjorken picture. If anything, it actually helps. Once it is established that the ratio $`h^+/h^{}`$ is independent of the overall volume, the only crucial ingredient in our argument is that the number fluctuations are all Poissonian, namely, $`\delta N^2/N1`$ The quantum statistics makes this ratio deviate from 1. For bosons (pions), restricting to small momenta or rapidity makes this ratio bigger. For fermions (quarks), it makes this ratio smaller. From Eqs.(11), (12), (15) and (16), it is clear that having a momentum cut-off can then only enhance the contrast. Resonance Contributions : As explained in a previous paper, neutral resonances introduce positive correlations between $`N_+`$ and $`N_{}`$ and hence lower the value of $`D_{\mathrm{had}}`$ from 4. In a thermal scenario studied in the same paper, we found that the resonances reduce the fluctuation by about 30 %. Hence for a realistic hadron gas, $`D_{\mathrm{had}}=N_{\mathrm{ch}}\delta R^23.`$ (37) This is still a factor of 3 bigger than the lattice result. Mixtures : If the system is a mixture of a QGP and a hadron gas, the signal should depend on the fractions. To a first approximation, it should be a linear combination of Eqs. (37) and (35) $`N_{\mathrm{ch}}\delta R^2=D_{\mathrm{had}}(1f)+D_{\mathrm{lat}}f`$ where $`f`$ is the QGP fraction. Even if $`f0.5`$, the signal can be still visible. Rapidity Correlations : If gluon interactions dominate the creation of hadrons in high energy hadron-hadron collisions, the unlike charged particles are strongly correlated in rapidity. Such strong correlation can further lower the value of $`D_{\mathrm{had}}`$. If the reduction is big enough to mimic QGP signal, we should certainly see the false signal in even in the nucleon-nucleon collisions. Using the two particle correlation data in $`pp`$ collisions compiled by Whitmore, we performed Monte Carlo simulations for such a scenario. The correction to (37) is found to be small ($`10\%`$). This is due to the experimental fact that even though the correlation is pronounced in the ‘connected’ part of the correlation $`R(y_+,y_{})`$, it is not so pronounced in the full two particle correlation function $`\rho _2(y_+,y_{})=\rho _1(y_+)\rho _1(y_{})[1+R(y_+,y_{})].`$ (38) For details, see . Finite Acceptance Correction : The finite size of the acceptance window introduces a factor of $`(1p)`$ corrections where $`p`$ is the fraction of the total multiplicity inside the acceptance window. This is easy to understand. If the detector sees a 100 % of all charged particles ($`p=1`$), the fluctuations should shrink down to zero due to the global charge conservation. Fortunately, this is a common factor that applies to both Eqs. (37) and (35) so that the ratio stays the same. More details will be reported in . Effects of Rescattering : The partons as well as hadrons are subject to rescattering during the course of a heavy ion collision. This in principle may affect the above charge fluctuations by diffusing the charge in rapidity space. However in the limit $`y_{\mathrm{total}}y_{\mathrm{accept}}1`$ these effects should be very small since they scale as the surface to volume ratio in rapidity space. To estimate the effect in the hadronic phase, we performed a simple Monte Carlo calculation where a Gaussian noise with $`\sigma =0.5`$ are added to each particle’s rapidity originating from a highly correlated source. Doing so increases the value of $`N_{\mathrm{ch}}\delta R^2`$ up to 40 % assuming the STAR acceptance at RHIC. In a slightly different context, Gavin also found that when a strong flow is established, the rescattering effect on the hadron rapidities is relatively small. Strangeness : Adding Kaons to a pion gas will not change the value of $`D_{\mathrm{had}}`$ because their contribution adds exactly the same amount to both numerator (c.f. Eq.(11)) and the denominator. For a quark gluon plasma, the lattice calculation suggests that at high temperature, the strangeness entropy is about 40 % of the $`u+d`$ entropy. Taking this at a face value changes the result (35) by less than 10 %. In conclusion, we showed in this paper that detection of QGP formation is quite possible through the simple measurement of $`h^+/h^{}`$ fluctuation. This measurement should be very feasible for STAR. We also emphasize that this is a direct confirmation of the lattice QCD results. What we considered here is the simplest ratio out of many possible ones that can behave quite differently in the presence of a QGP. For instance, the strangeness anti-strangeness ratio fluctuations can provide us with a valuable handle on the strangeness distributions with or without the formation of a QGP. These and related issues are under active investigation and will be reported elsewhere . S.J. would like to thank M. Bleicher for many discussions. This work was supported by the Director, Office of Energy Research, Office of High Energy and Nuclear Physics, Division of Nuclear Physics, and by the Office of Basic Energy Sciences, Division of Nuclear Sciences, of the U.S. Department of Energy under Contract No. DE-AC03-76SF00098. Note added: After finishing this work, we received a preprint by Asakawa, Heinz and Müller which addresses similar issues.
warning/0003/hep-th0003194.html	ar5iv	text	# Quantization Ambiguity and Supersymmetric Ground State Wave Functions ## 1 Introduction In this paper we study how the quantization ambiguity, which implies that quantization on space with a nontrivial topology such as $`S^1`$ inevitably yields an undetermined parameter into the theory , affects the supersymmetric ground state wave functions of a model of supersymmetric quantum mechanics on $`S^1`$. Quantum mechanics on $`S^1`$ was studied for the first time in the path-integral formalism . We observe that the aforementioned parameter appears as a phase factor in the Feynman kernel due to the nontrivial topology of configuration space and, as a result, there are many distinct propagators labeled by the parameter. One can also consider the same effect in the Lagrangian by adding a total derivative term whose coefficient is given by the parameter. Accordingly, the canonical momentum is shifted by the amount of the parameter, so that it can be interpreted as a constant gauge field. The total derivative term has physical implications at the quantum level for space with nontrivial topology . Quantization on $`S^1`$ is much different from that on one-dimensional Euclidean space. In the language of canonical formalism, the latter case is that the representation of the canonical algebra is uniquely determined up to a unitary equivalent representation. There is essentially one quantum mechanics on the space. In the former case, however, there is an infinite number of inequivalent representations of the fundamental algebra, which is introduced as a generalization of the canonical algebra by Ohnuki and Kitakado in order to formulate quantum mechanics on $`S^1`$. As a result, there exists various quantum mechanics on $`S^1`$. The various quantum mechanics on $`S^1`$ are parametrized by the undetermined parameter. The parameter is interpreted as a constant gauge field . The gauge field can exist and has effects on observables at the quantum level. It is the existence of the gauge field that leads to the various quantum mechanics on $`S^1`$. It may be interesting to study the possible effects of the gauge field on the supersymmetric ground state (zero-energy state) wave functions of supersymmetric quantum mechanics on $`S^1`$. The zero-energy state wave functions can be obtained in closed form because, thanks to the supersymmetry algebra satisfied by the system, the wave functions are obtained by solving simple first-order equations in many cases. Therefore, it may be possible to study the effects as analytically as possible. Supersymmetric quantum mechanics has been studied in great detail and applied to many physics fields . Actually, it provides us with an example of the dynamical supersymmetry breaking by instantons in certain models . In those models, the normalizability of the supersymmetric ground state wave function crucially depends on the leading term in the superpotential, by which we determine whether or not the supersymmetry is broken. The semiclassical instanton approximation has been used to estimate the ground state energy for the system with broken supersymmetry<sup>1</sup><sup>1</sup>1Estimating the ground state energy is a subtle problem. Actually, it is reported in that the instanton calculation is very limited though it gives us an excellent estimation in some cases and the breakdown of symmetry is caused by the interplay of perturbative and nonperturbative effects.. In this paper we will find an another mechanism of supersymmetry breaking. The very existence of the gauge field twists boundary conditions of supersymmetric ground state wave functions. For certain values of the gauge field, the wave functions do not satisfy a required periodic boundary condition and become unphysical though they are normalizable. The supersymmetry breaking does not depend on the structure of the superpotential, unlike the usual supersymmetry breaking discussed in supersymmetric quantum mechanics. Supersymmetry can be broken by the gauge field, that is, the quantization ambiguity. In the next section we shall introduce a model of supersymmetric quantum mechanics on $`S^1`$ after reviewing briefly the quantum mechanics on $`S^1`$ formulated by Ohnuki and Kitakado. And then, we shall discuss how the gauge field affects the supersymmetric ground state wave functions and how it yields supersymmetry breaking. We shall also study an infinite limit of the radius of $`S^1`$. The supersymmetric harmonic oscillator is realized in the limit with the strength of the oscillator being constant. The final section is devoted to conclusions and a discussion, where we shall also discuss the similarities between our mechanism of supersymmetry breaking and that through the boundary conditions of fields in supersymmetric quantum field theory on compactified space. ## 2 Supersymmetric Quantum Mechanics on $`S^1`$ We shall study effects of the quantization ambiguity, which implies that quantization on a space with nontrivial topology yields an undetermined parameter, on supersymmetric ground state wave functions of supersymmetric quantum mechanics on $`S^1`$. Let us consider a system in which there is the fermionic operator $`\widehat{Q}_i`$ that commutes with the Hamiltonian $`\widehat{H}`$ and satisfies the supersymmetry algebra $$[\widehat{Q}_i,\widehat{H}]=0,\{\widehat{Q}_i,\widehat{Q}_j\}=\delta _{ij}\widehat{H},i=1,\mathrm{}N.$$ (1) $`N=2`$ is the simplest case and it is of our interest. Since the Hamiltonian is positive semidefinite, a supersymmetric state $`\widehat{Q}_i\|\mathrm{\Psi }=0`$ is automatically a zero-energy ground state. Conversely, if we have a zero-energy state, it has to be a supersymmetric ground state. Thanks to this property, finding supersymmetric ground states is reduced to solving simple first-order equations instead of solving the second-order equation $`\widehat{H}\|\mathrm{\Psi }=0`$. A key point for our study is that the fermionic operator $`\widehat{Q}_i(i=1,2)`$ should be written in terms of the operators which are appropriate to describe the quantum mechanics on $`S^1`$ as shown in the subsection $`\mathbf{2.2}`$. And we shall study the supersymmetric ground state wave functions of such a system <sup>2</sup><sup>2</sup>2Concerning supersymmetric quantum mechanics and the ground state wave functions, it is known that a supersymmetric quantum mechanical system can be constructed by using the ground state wave functions of a nonsupersymmetric Hamiltonian . In this approach the superpotential can be given by the ground state wave function of the system. And the approach is extended to generally covariant systems such as relativistic particles interacting with external gauge fields and gravitational fields etc. .. ### 2.1 Quantum Mechanics on $`S^1`$ Before we proceed to a model of supersymmetric quantum mechanics on $`S^1`$, it may be important and instructive to review briefly the Ohnuki-Kitakado formulation of quantum mechanics on $`S^1`$ . Those who are familiar with their formulation can skip this subsection and go directly to the subsection $`\mathbf{2.2}`$ where the supersymmetric quantum mechanics on $`S^1`$ is introduced. The discussions below are based on a paper in a part of which the quantum mechanics on $`S^1`$ is summarized clearly. The quantum mechanics on $`S^1`$ is defined by a self-adjoint operator $`\widehat{G}`$ and a unitary operator $`\widehat{W}`$ satisfying the commutation relation $$[\widehat{G},\widehat{W}]=h\text{-}\widehat{W}.$$ (2) The operators $`\widehat{G},\widehat{W}`$, and $`\widehat{W}^{}`$ generate an algebra. Let us construct its representation. We shall start with an eigenvalue equation $$\widehat{G}\|\alpha =h\text{-}\alpha \|\alpha \mathrm{with}\alpha \|\alpha =1,$$ (3) where an eigenvalue $`\alpha `$ is a real number. It is easy to see that $`\widehat{W}(\widehat{W}^{})`$ raises (lowers) the eigenvalues of $`\widehat{G}`$ $$\widehat{G}\widehat{W}\|\alpha =h\text{-}(\alpha +1)\widehat{W}\|\alpha ,\widehat{G}\widehat{W}^{}\|\alpha =h\text{-}(\alpha 1)\widehat{W}^{}\|\alpha .$$ (4) A state vector defined by $$\|n+\alpha \widehat{W}^n\|\alpha ,n=\mathrm{integer},$$ (5) is also an eigenstate of $`\widehat{G}`$: $$\widehat{G}\|n+\alpha =h\text{-}(n+\alpha )\|n+\alpha .$$ (6) For fixed $`\alpha `$, our Hilbert space, denoted by $`_\alpha `$ where the two operators $`\widehat{G},\widehat{W}`$ are defined, is given by completing the vector space of linear combinations of $`\|n+\alpha (n=0,\pm 1,\pm 2,\mathrm{})`$. The set of state vectors forms the orthocomplete system in $`_\alpha `$. Therefore, we have $$m+\alpha \|n+\alpha =\delta _{mn},\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}\|n+\alpha n+\alpha \|=\mathrm{𝟏}_\alpha ,$$ (7) where $`\mathrm{𝟏}_\alpha `$ is an identity operator in $`_\alpha `$. Equation (6) and $`\widehat{W}\|n+\alpha =\|n+1+\alpha `$ define an irreducible representation of the algebra (2) on $`_\alpha `$. The classification of the irreducible representation of the algebra may be done by noting that (i) $`_\alpha `$ and $`_\beta `$ are unitary equivalent Hilbert space if and only if $`\alpha \beta =\mathrm{integer}`$ and (ii) for an arbitrary irreducible representation $``$ of the algebra, there exists a real number $`\alpha `$ such that $``$ is the unitary equivalent of $`_\alpha `$. Thus, the classification is completed; that is, all the inequivalent irreducible representations are given by the Hilbert space $`_\alpha (0\alpha <1)`$. It should be emphasized that the algebra (2) has an infinite number of inequivalent representations characterized by an undetermined parameter $`\alpha `$, as contrary to the usual irreducible representation of the canonical algebra on one-dimensional Euclidean space. So far, we have constructed the $`\widehat{G}`$-diagonal representation. One can also go to the $`\widehat{W}`$-diagonal representation by which we will obtain wave functions in the quantum mechanics on $`S^1`$. For fixed representation space $`_\alpha `$, since $`\widehat{W}`$ is a unitary operator, the eigenvalue equation for it may be written as $$\widehat{W}\|\theta =\mathrm{e}^{i\theta }\|\theta .$$ (8) Its solution is $$\|\theta =\kappa (\theta )\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}\mathrm{e}^{in\theta }\|n+\alpha ,$$ (9) where $`\theta `$ is a real parameter and $`\kappa (\theta )`$ is an arbitrary complex-valued function satisfying $`\left\|\kappa (\theta )\right\|=1`$ and $`\kappa (\theta +2\pi )=\kappa (\theta )`$. It is not difficult to show that $`\|\theta +2\pi n`$ $`=`$ $`\|\theta ,n=\mathrm{integer},`$ (10) $`\theta \|\theta ^{}`$ $`=`$ $`2\pi {\displaystyle \underset{n=\mathrm{}}{\overset{n=+\mathrm{}}{}}}\delta (\theta \theta ^{}+2\pi n),`$ (11) $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\theta }{2\pi }}\|\theta \theta \|`$ $`=`$ $`{\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}\|n+\alpha n+\alpha \|=\mathrm{𝟏}_\alpha ,`$ (12) $`\mathrm{exp}(i\lambda {\displaystyle \frac{\widehat{G}}{h\text{-}}})\|\theta `$ $`=`$ $`\mathrm{e}^{i\lambda \alpha }\kappa (\theta )\kappa ^{}(\theta +\lambda )\|\theta +\lambda ,`$ (13) where $`\mathrm{𝟏}_\alpha `$ is an identity operator in $`_\alpha `$. These correspond to periodicity, orthonormality, completeness, and translation for the eigenstate of $`\widehat{W}`$. Let us note that it may be possible from Eqs. (8) and (13) to identify $`\widehat{G}`$ and $`\widehat{W}`$ with the momentum and the position operators on $`S^1`$, respectively. Now, let $`\|\psi `$ be a state vector and we define a wave function $`\psi (\theta )`$ on $`S^1`$ as follows $$\psi (\theta )\theta \|\psi .$$ (14) Taking the inner product of Eq. (13) with $`\|\psi `$, we obtain $$\theta \|\mathrm{exp}(i\lambda \frac{\widehat{G}}{h\text{-}})\|\psi =\mathrm{e}^{i\lambda \alpha }\kappa ^{}(\theta )\kappa (\theta +\lambda )\theta +\lambda \|\psi ,$$ (15) from which the $`\widehat{W}`$-diagonal representation of $`\widehat{G}`$ is given by $$\theta \|\widehat{G}\|\psi =\left[ih\text{-}\frac{}{\theta }ih\text{-}\kappa ^{}(\theta )\frac{\kappa (\theta )}{\theta }+h\text{-}\alpha \right]\psi (\theta ).$$ (16) We also obtain, from Eq. (8), $$\theta \|\widehat{W}\|\psi =\mathrm{e}^{i\theta }\psi (\theta ).$$ (17) The inner product on $`S^1`$ is expressed in terms of the wave function as $$\chi \|\psi =_0^{2\pi }\frac{d\theta }{2\pi }\chi ^{}(\theta )\psi (\theta ).$$ (18) Thus, the representation of Hilbert space, which is defined by Eqs. (16) and (17), is the space of the square integrable function on $`S^1`$. Let us note that all wave functions have to satisfy the periodic boundary condition $`\psi (\theta +2\pi n)=\psi (\theta )`$, which is a direct consequence of Eq. (10). This periodicity is essential when we study the supersymmetric ground state wave functions of the supersymmetric quantum mechanics on $`S^1`$. Let us next present the physical meaning of the parameter $`\alpha `$. To this end, let us redefine $`\kappa (\theta )`$ by utilizing the arbitrariness of it in such a way that $`\kappa (\theta )=\omega (\theta )\kappa ^{}(\theta )`$, where $`\omega (\theta )`$ has to satisfy $`\left\|\omega (\theta )\right\|=1`$ and $`\omega (\theta +2\pi )=\omega (\theta )`$. It follows that $`\|\theta =\omega (\theta )\|\theta ^{}`$, so that the transformed wave function $`\psi ^{}(\theta )`$ is given by $$\psi ^{}(\theta )=\omega (\theta )\psi (\theta ).$$ (19) According to this redefinition, the $`\widehat{W}`$-diagonal representation for $`\widehat{G}`$ becomes $${}_{}{}^{}\theta \|\widehat{G}\|\psi =[ih\text{-}\frac{}{\theta }+A^{}(\theta )]\psi (\theta ),$$ (20) where we have defined $$A^{}(\theta )A(\theta )+ih\text{-}\omega ^{}(\theta )\frac{\omega (\theta )}{\theta },A(\theta )ih\text{-}\kappa ^{}(\theta )\frac{\kappa (\theta )}{\theta }+h\text{-}\alpha .$$ (21) Equations (19) and (21) stand for the gauge transformation. Therefore, the parameter $`\alpha `$ has the meaning of the gauge field. It is easy to see that the gauge field has the properties (i) $`A(\theta )`$, assumed to be an arbitrary real-valued function satisfying the periodic boundary condition $`A(\theta +2\pi )=A(\theta )`$, can always be made a constant function $`A^{}(\theta )=\alpha `$ by a gauge transformation and (ii) for two constant functions $`A^{}(\theta )=\alpha `$ and $`A^{}(\theta )=\beta `$, these are connected by a unique gauge transformation if and only if $`\beta \alpha `$ is an integer. Thus, we arrive at an important conclusion that all the inequivalent gauge fields are given by $`A_\alpha \alpha (0\alpha <1)`$<sup>3</sup><sup>3</sup>3 Dirac’s approach to the quantization for a constrained system does not yield an infinite number of inequivalent representations; that is, it corresponds to only $`\alpha =0`$. Quantization and embedding $`S^1`$ into higher dimensional space $`𝐑^2`$ is not a “commutative” procedure relating to each other. Let us also note that the Ohnuki-Kitakado formulation is independent of the dynamics, in contrast to Dirac’s approach.. Hereafter, we choose $`\kappa (\theta )=1`$ for simplicity. It is a very special feature of the quantum mechanics on $`S^1`$ that the inequivalent gauge field is restricted to be $`0\alpha <1`$. Another way of looking at it is that if we perform a gauge transformation by $`\psi (\theta )\psi ^{}(\theta )=\mathrm{e}^{in\theta }\psi (\theta ),`$ we see that the gauge fields $`A(\theta )`$ and $`A(\theta )nh\text{-}`$ are equivalent for $`n=`$integer. $`n`$ has to be restricted to be an integer; otherwise, the transformed wave function $`\psi ^{}(\theta )`$ does not satisfy the required periodic boundary condition. Therefore, the inequivalent gauge field is given by $`0A(\theta )<h\text{-}`$ , which means $`0\alpha <1`$. Let us note that the gauge transformation by $`\mathrm{e}^{in\theta }`$ with $`n=`$ noninteger is a singular gauge transformation and is strictly forbidden. Different values of the gauge field give different quantum mechanics on $`S^1`$. It may be helpful to note that the gauge field $`\alpha `$ may correspond to the magnitude of the magnetic flux $`e\mathrm{\Phi }/2\pi h\text{-}c`$ through $`S^1`$ in the Aharanov-Bohm effect. The different magnitude of the flux actually gives different physics. ### 2.2 Supersymmetric Simple Pendulum Now, we are ready to introduce a model of the supersymmetric quantum mechanics on $`S^1`$. According the discussion above, the two operators $`\widehat{G}`$ and $`\widehat{W}`$, which correspond to the momentum and the position of a particle on $`S^1`$, are fundamental. It may be natural to construct a quantum Hamiltonian in terms of these operators. The Hamiltonian in our model is assumed to satisfy the supersymemtry algebra (1), so that the fermionic operators $`\widehat{Q}_i(i=1,2)`$ also have to be given in terms of them. We will discuss the classical counterpart of the quantum hamiltonian constructed in this way later. Let us define the fermionic operator $`\widehat{Q}_i(i=1,2)`$ by $`\widehat{Q}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\widehat{Q}_1+i\widehat{Q}_2\right)`$ (22) $`=`$ $`\left({\displaystyle \frac{1}{\sqrt{2m}R}}\widehat{G}+iV(\widehat{W},\widehat{W}^{})\right)\widehat{\xi }\widehat{q}\widehat{\xi },`$ $`\widehat{\overline{Q}}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\widehat{Q}_1i\widehat{Q}_2\right)`$ (23) $`=`$ $`\left({\displaystyle \frac{1}{\sqrt{2m}R}}\widehat{G}iV(\widehat{W},\widehat{W}^{})\right)\widehat{\overline{\xi }}\widehat{q}^{}\widehat{\overline{\xi }}.`$ $`V(\widehat{W},\widehat{W}^{})`$, which is called the superpotential hereafter, is a hermitian operator in terms of $`\widehat{W}`$ and $`\widehat{W}^{}`$. Here $`m`$ and $`R`$ stand for the mass of a particle and the radius of $`S^1`$, respectively. The fermionic variables $`\widehat{\xi },\widehat{\overline{\xi }}`$ satisfy the algebra $$\{\widehat{\xi },\widehat{\overline{\xi }}\}=1,\widehat{\xi }^2=\widehat{\overline{\xi }}^2=0.$$ (24) Then, the Hamiltonian is given by $`\widehat{H}`$ $`=`$ $`\{\widehat{Q},\widehat{\overline{Q}}\}`$ (25) $`=`$ $`{\displaystyle \frac{1}{2mR^2}}\widehat{G}^2+V^2(\widehat{W},\widehat{W}^{})`$ $``$ $`{\displaystyle \frac{i}{\sqrt{2m}R}}\left(\widehat{G}V(\widehat{W},\widehat{W}^{})V(\widehat{W},\widehat{W}^{})\widehat{G}\right)[\widehat{\xi },\widehat{\overline{\xi }}],`$ where we have used Eq. (24). Here, it may be necessary to discuss the classical counterpart of the quantum Hamiltonian (25). To this end, let us note that the fundamental algebra (2) may be actually inferred by the classical Poisson’s brackets for the angle variable $`\theta `$ and the correponding momentum $`P_\theta `$ in the polar coordinate: $$\{P_\theta ,\mathrm{e}^{i\theta }\}_P=i\mathrm{e}^{i\theta }.$$ (26) If we replace the classical Poisson’s brackets by the commutation relation divided by $`ih\text{-}`$ , we obtain the fundamental algebra (2) by identifying $`\mathrm{e}^{i\theta }`$ and $`P_\theta `$ with $`\widehat{W}`$ and $`\widehat{G}`$, respectively. This is the same identification stated earlier. Therefore, in the classical limit we may replace $`\widehat{W}`$ by $`\mathrm{e}^{i\theta }`$ and $`\widehat{G}`$ by $`P_\theta `$. According to these replacements, we obtain a classical Hamiltonian, ignoring the fermionic variables $`\widehat{\xi },\widehat{\overline{\xi }}`$: $$\widehat{H}H_{cl}=\frac{P_\theta ^2}{2mR^2}+V^2(\mathrm{e}^{i\theta },\mathrm{e}^{i\theta }).$$ (27) If we choose $`V(\mathrm{e}^{i\theta },\mathrm{e}^{i\theta })`$ as $$V(\mathrm{e}^{i\theta },\mathrm{e}^{i\theta })=\sqrt{\frac{mg_NR}{2}}\mathrm{sin}\theta ,$$ (28) the classical Hamiltonian describes a simple pendulum with angle $`2\theta `$. Here $`g_N`$ is the gravitation accelerator constant. On the other hand, given the superpotential (28), the quantum counterpart of it is obtained by $$V(\widehat{W},\widehat{W}^{})=\sqrt{\frac{mg_NR}{2}}\left(\frac{\widehat{W}\widehat{W}^{}}{2i}\right).$$ (29) Having this superpotential, the model of the supersymmetric quantum mechanics on $`S^1`$, that is, the supersymmetric simple pendulum, is given by the Hamiltonian $$\widehat{H}=\frac{1}{2mR^2}\widehat{G}^2+\frac{mg_NR}{2}\mathrm{sin}^2\theta \frac{h\text{-}}{2}\sqrt{\frac{g_N}{R}}\mathrm{cos}\theta [\widehat{\xi },\widehat{\overline{\xi }}]$$ (30) in the $`\widehat{W}`$-diagonal representation. It is understood that $`\widehat{G}=ih\text{-}/\theta +h\text{-}\alpha `$ and $`\alpha `$ is the gauge field discussed in the previous subsection. Let us study the supersymmetric ground state wave functions of the supersymmetric simple pendulum whose hamiltonian is given by Eq. (30). It follows from the algebra (1) that the supersymmetric ground states must be zero-energy states satisfied by $$\widehat{Q}\|\mathrm{\Psi }=0\mathrm{and}\widehat{\overline{Q}}\|\mathrm{\Psi }=0.$$ (31) Let us introduce a matrix representation for the fermionic variables. It is easy to see that the matrix representations given by $$\widehat{\xi }=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right),\widehat{\overline{\xi }}=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)$$ (32) satisfy the algebra (24). Then, it follows that $$[\widehat{\xi },\widehat{\overline{\xi }}]=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\sigma ^3.$$ Using these, the Hamiltonian (30) becomes $`\widehat{H}`$ $`=`$ $`\left({\displaystyle \frac{1}{2mR^2}}\widehat{G}^2+{\displaystyle \frac{mg_NR}{2}}\mathrm{sin}^2\theta \right)\mathrm{𝟏}_{2\times 2}+{\displaystyle \frac{h\text{-}}{2}}\sqrt{{\displaystyle \frac{g_N}{R}}}\sigma ^3\mathrm{cos}\theta `$ (37) $`=`$ $`\left(\begin{array}{cc}\widehat{q}^{}\widehat{q}& 0\\ 0& \widehat{q}\widehat{q}^{}\end{array}\right)\left(\begin{array}{cc}\widehat{H}_+& 0\\ 0& \widehat{H}_{}\end{array}\right).`$ In the matrix representation, the hamiltonian is a $`2\times 2`$ matrix. Since the Hamiltonian (37) commutes with an operator $`\widehat{S}^F\sigma ^3/2`$, the eigenstates of the Hamiltonian is labeled by the eigenvalues of $`\widehat{S}^F`$. Let us call the two states $`\|+`$ and $`\|`$ fermion numbers $`+1/2`$ and $`1/2`$, respectively<sup>4</sup><sup>4</sup>4One may call the two states spin-up and spin-down states. Or by introducing $`\widehat{f}=\frac{1}{2}+\frac{1}{2}[\widehat{\xi },\widehat{\overline{\xi }}]`$ whose eigenvalues are $`0,1`$, one may denote them by fermion numbers $`0,1`$.. The state vector is, now, a two-component vector $$\|\mathrm{\Psi }=\left(\genfrac{}{}{0pt}{}{\|+}{\|}\right).$$ (38) In the $`\widehat{W}`$-diagonal representation, it may be written as $$\mathrm{\Psi }(\theta )=\left(\genfrac{}{}{0pt}{}{\psi _{+\frac{1}{2}}(\theta )}{\psi _{\frac{1}{2}}(\theta )}\right).$$ (39) The Hamiltonian $`\widehat{H}`$ is diagonalized with respect to the fermion number $`\pm 1/2`$. In this matrix representation, Eq. (31) is read as $`\widehat{q}\psi _{+\frac{1}{2}}(\theta )=\left({\displaystyle \frac{1}{\sqrt{2m}R}}(ih\text{-}{\displaystyle \frac{}{\theta }}+h\text{-}\alpha )+i\sqrt{{\displaystyle \frac{mg_NR}{2}}}\mathrm{sin}\theta \right)\psi _{+\frac{1}{2}}(\theta )=0,`$ $`\widehat{q}^{}\psi _{\frac{1}{2}}(\theta )=\left({\displaystyle \frac{1}{\sqrt{2m}R}}(ih\text{-}{\displaystyle \frac{}{\theta }}+h\text{-}\alpha )i\sqrt{{\displaystyle \frac{mg_NR}{2}}}\mathrm{sin}\theta \right)\psi _{\frac{1}{2}}(\theta )=0.`$ (40) Solutions for Eq. (40) are found to be $$\psi _{+\frac{1}{2}}(\theta )=\frac{1}{\sqrt{I_0(2z)}}\mathrm{exp}(i\alpha \theta \frac{z}{h\text{-}}\mathrm{cos}\theta ),\psi _{\frac{1}{2}}(\theta )=\frac{1}{\sqrt{I_0(2z)}}\mathrm{exp}(i\alpha \theta +\frac{z}{h\text{-}}\mathrm{cos}\theta ),$$ (41) where $`I_0(2z)`$ in the normalization factor is the zeroth-order modified Bessel function and we have defined a dimensionless parameter $$\frac{mR^2}{h\text{-}}\sqrt{\frac{g_N}{R}}\frac{mR^2}{h\text{-}}\omega \frac{z}{h\text{-}}.$$ These are normalizable solutions. Thus, one may say that the zero-energy states, that is, supersymmetric ground states, exist in the model and the supersymmetry is unbroken. This is, however, a hasty conclusion. In addition to the normalizability, all the wave functions have to satisfy the periodic boundary condition $`\mathrm{\Psi }(\theta +2\pi )=\mathrm{\Psi }(\theta )`$, that is, $`\psi _{\pm \frac{1}{2}}(\theta +2\pi )=\psi _{\pm \frac{1}{2}}(\theta )`$, which follows from Eq. (10) <sup>5</sup><sup>5</sup>5Let us note that the two-component spinor state (38) does not have a minus sign under $`2\pi `$ rotation in this case because the rotation is done by the usual rotation matrix in two dimensions.. It is easy to see from Eq. (41) that $$\psi _{\pm \frac{1}{2}}(\theta +2\pi )=\mathrm{e}^{i2\pi \alpha }\psi _{\pm \frac{1}{2}}(\theta ).$$ (42) The boundary condition for the zero-energy state wave functions is twisted by the gauge field $`\alpha `$. The zero-energy state wave functions do not satisfy the required periodic boundary condition except for $`\alpha =\mathrm{integer}`$. Since the inequivalent representation is given by $`0\alpha <1`$, they are inconsistent with the periodic boundary condition and become unphysical wave functions for $`0<\alpha <1`$. Therefore, the supersymmetry can be broken due to the gauge field $`\alpha `$. Let us note that the Witten index $`\mathrm{Tr}(1)^{\widehat{f}}=n_B^{E=0}n_F^{E=0}`$ vanishes in our model. It is easy to see that $`n_B^{E=0}=n_F^{E=0}=1`$ for $`\alpha =\mathrm{integer}`$ and that $`n_B^{E=0}=n_F^{E=0}=0`$ for $`\alpha =`$ noninteger. Unlike the usual supersymmetry breaking, in which the leading term of the superpotential determines whether or not supersymmetry is broken, our breaking of supersymmetry does not depend on the structure of the superpotential. It is entirely due to the existence of the gauge field $`\alpha `$, which is an inevitable consequence of the quantization ambiguity when one quantizes the theory on topologically nontrivial space such as $`S^1`$. The gauge field has the effect of twisting the boundary conditions of the zero-energy state wave functions. Among the various supersymmetric quantum mechanics on $`S^1`$ led by the gauge field, it includes theories with broken supersymmetry due to the gauge field. Let us note that in this context there is no mechanism to determine the values of $`\alpha `$ or what values of $`\alpha `$ we should take. Note that for noninteger values of the gauge field $`\alpha `$, the gauge field can not be removed in Eq. (40) by the gauge transformation $$\psi _{\pm \frac{1}{2}}(\theta )\psi _{\pm \frac{1}{2}}^{}(\theta )=\mathrm{e}^{i\alpha \theta }\psi _{\pm \frac{1}{2}}(\theta ).$$ (43) As we stated before, the inequivalent gauge field is given by $`0\alpha <1`$. Any gauge field in this range cannot be connected by a regular gauge transformation with $`\mathrm{e}^{in\theta }(n=\mathrm{integer})`$. Only a singular gauge transformation can do it, but it destroys the required periodic boundary condition for the transformed wave function. The singular gauge transformation is strictly forbidden, so that the gauge field cannot be gauged away. Let us briefly comment on the ground state energy. We have a physical supersymmetric ground state wave function for $`\alpha =\mathrm{integer}`$, so that the ground state energy is exactly zero. On the other hand, for $`0<\alpha <1`$, supersymmetry is broken. The ground state energy is nonzero (positive). Estimating the ground state energy is a subtle problem as studied in . The Hamiltonian (37) cannot be solved analytically. The bosonic potential $`\frac{mg_NR}{2}\mathrm{sin}^2\theta `$ is periodic and the classical vacuum has a periodic structure. One may expect that there is a instantonlike classical solution, which gives a finite Euclidean action, connecting the two vacua with different fermion numbers $`(\pm 1/2)`$. Actually, there exists such a classical solution in our model. It is given by $`\mathrm{cos}\theta _{cl}(\tau )=\pm \mathrm{tanh}(\omega (\tau \tau _0))`$ with the classical Euclidean action being $`2z/h\text{-}`$ . And the fermion zero mode exists in this classical background. Therefore, we expect tunneling to occur between the two vacua. According to the semiclassical argument, the tunneling effect shifts the ground state energy to give an exponentially small amount of energy in the form of $`\mathrm{exp}(2z/h\text{-})\times \mathrm{cos}2\pi \alpha `$<sup>6</sup><sup>6</sup>6This band structure $`\mathrm{cos}2\pi \alpha `$ can be understood from the effective action obtained by the transition amplitude $`K(\theta _f,t;\theta _i,0)=\theta _f\|\mathrm{exp}(i\widehat{H}t/h\text{-})\|\theta _i=_{n=\mathrm{}}^+\mathrm{}_{nwinding}𝒟\theta \mathrm{exp}(iS_{eff}/h\text{-})`$, where the effective action is given by $`S_{eff}=𝑑t\frac{mR^2}{2}(\frac{d\theta }{dt})^2\frac{mg_NR}{2}\mathrm{sin}^2\theta +\frac{1}{2}\sqrt{\frac{g_N}{R}}\mathrm{cos}\theta [\xi ,\xi ^{}]+i\xi ^{}\frac{d\xi }{dt}\alpha \frac{d\theta }{dt}`$. The “topological” term $`\alpha \dot{\theta }`$ is the origin of such band structure.. On the other hand, for very small $`z`$, we can resort to perturbation theory to obtain the energy spectrum of the Hamiltonian. The ground state energy is given by $`E_0\frac{1}{2mR^2}(\alpha ^2+O(z^2))`$, where we have set $`h\text{-}=1`$. The gauge field $`\alpha `$ is a dominant contribution to the ground state energy in this case. It may be interesting to consider the $`R\mathrm{}`$ limit. So far, we have fixed the radius $`R`$ of $`S^1`$. If $`R`$ varies to become large, we expect that the arc of an arbitrary part of $`S^1`$ will approach a straight segment. In the limit of $`R\mathrm{}`$, one-dimensional Euclidean space will be recovered. In order to study the limit, let us define a variable $`xR\theta `$. In terms of this new variable, the Hamiltonian (37) is written as, remembering $`\widehat{G}=ih\text{-}/\theta +h\text{-}\alpha `$ in the $`\widehat{W}`$-diagonal representation, $$\widehat{H}=\left(\frac{h\text{-}^2}{2m}\left(\frac{}{x}+i\frac{\alpha }{R}\right)^2+\frac{mg_NR}{2}\mathrm{sin}^2(\frac{x}{R})\right)\mathrm{𝟏}_{2\times 2}+\frac{h\text{-}}{2}\sqrt{\frac{g_N}{R}}\sigma ^3\mathrm{cos}(\frac{x}{R}).$$ (44) If we take the limit of $`R\mathrm{}`$ naively, it becomes trivial for the Hamiltonian to yield the one for a free particle on $`S^1`$: $$\widehat{H}=\frac{h\text{-}^2}{2m}\frac{^2}{x^2}\mathrm{𝟏}_{2\times 2}.$$ (45) In order to obtain an interacting theory, one has to take the limit, keeping a relation given by $$\frac{m}{h\text{-}}\sqrt{\frac{g_N}{R}}\frac{m\omega }{h\text{-}}=\left(\mathrm{strength}\mathrm{of}\mathrm{oscillator}\right)^2=\mathrm{const}.$$ (46) Then, we obtain $$\widehat{H}=\left(\frac{h\text{-}^2}{2m}\frac{^2}{x^2}+\frac{m\omega ^2}{2}x^2\right)\mathrm{𝟏}_{2\times 2}+\frac{h\text{-}\omega }{2}\sigma ^3+O(\frac{1}{R^2}).$$ (47) This is the well-known Hamiltonian for the supersymmetric harmonic oscillator with angular frequency $`\omega =\sqrt{g_N/R}`$. Likewise, by taking the same limit, the fermionic operators $`\widehat{Q},\widehat{\overline{Q}}`$ become $`\widehat{Q}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2m}}}\left(\widehat{p}+iW(\widehat{x})\right)\widehat{\xi }+O({\displaystyle \frac{1}{R^2}})\widehat{Q}_{susy}+O({\displaystyle \frac{1}{R^2}}),`$ $`\widehat{\overline{Q}}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2m}}}\left(\widehat{p}iW(\widehat{x})\right)\widehat{\overline{\xi }}+O({\displaystyle \frac{1}{R^2}})\widehat{\overline{Q}}_{susy}+O({\displaystyle \frac{1}{R^2}}),`$ (48) where we have defined $`W(\widehat{x})m\omega \widehat{x}`$ and $`\widehat{p}ih\text{-}/x`$. It is easy to check that the Hamiltonian (47) satisfies the supersymmetry algebra (1) with the supercharges (48) if we use the canonical commutation relation $`[\widehat{p},\widehat{x}]=ih\text{-}`$ and Eq. (32). The usual supersymmetry, by which the supersymmetry transformations between boson $`(\widehat{x})`$ and fermion $`(\widehat{\xi },\widehat{\overline{\xi }})`$ are generated, is realized in the limit of $`R\mathrm{}`$ with Eq. (46). The supersymmetric ground state wave functions for the Hamiltonian (47) are obtained by solving the first-order equation $`\widehat{Q}_{susy}\|\mathrm{\Psi }=0`$ and $`\widehat{\overline{Q}}_{susy}\|\mathrm{\Psi }=0`$. Using the same matrix representation as before, their solutions are $$\psi _{+\frac{1}{2}}^{h.o}(x)\mathrm{exp}(+\frac{m\omega }{2h\text{-}}x^2),\psi _{\frac{1}{2}}^{h.o}(x)=\left(\frac{m\omega }{\pi h\text{-}}\right)^{1/4}\mathrm{exp}(\frac{m\omega }{2h\text{-}}x^2).$$ (49) There are two candidates for the supersymmetric ground state (zero-energy state) wave functions. The one is physical and its wave function is given by $`\psi _{\frac{1}{2}}^{h.o}(x)`$. The other one $`\psi _{+\frac{1}{2}}^{h.o}(x)`$ is unphysical because of its non-normalizability. Therefore, we have one supersymmetric ground (zero-energy) state. This is consistent with the exact energy spectrum of the Hamiltonian (47). As easily seen from the Hamiltonian, there exists one supersymmetric ground state. In fact, the solutions (49) can be obtained by taking the the limit of $`R\mathrm{}`$ with the relation (46) in Eq. (41). By noting $`I_0(z)\mathrm{e}^z/\sqrt{2\pi z}`$ for large $`z`$, we obtain $$\stackrel{~}{\psi }_{+\frac{1}{2}}\mathrm{exp}(+\frac{m\omega }{2h\text{-}}x^2),\stackrel{~}{\psi }_{\frac{1}{2}}(x)=\left(\frac{m\omega }{\pi h\text{-}}\right)^{1/4}\mathrm{exp}(\frac{m\omega }{2h\text{-}}x^2),$$ (50) where we have redefined the normalization as $`\stackrel{~}{\psi }_{\pm \frac{1}{2}}(x)dx\psi _{\pm \frac{1}{2}}(\theta )\frac{d\theta }{\sqrt{2\pi R}}`$. These are the same as Eqs. (49). The Witten index is $`\mathrm{Tr}(1)^{\widehat{f}}=1`$ in this case. A simple generalization of the model is an $`N`$-component one. The fermionic operators are defined by $`\widehat{Q}`$ $`=`$ $`{\displaystyle \underset{a=1}{\overset{N}{}}}\left({\displaystyle \frac{1}{\sqrt{2m}R_a}}\widehat{G}_a+i\widehat{V}_a\right)\widehat{\xi }_a{\displaystyle \underset{a=1}{\overset{N}{}}}\widehat{q}_a\widehat{\xi }_a,`$ (51) $`\widehat{\overline{Q}}`$ $`=`$ $`{\displaystyle \underset{a=1}{\overset{N}{}}}\left({\displaystyle \frac{1}{\sqrt{2m}R_a}}\widehat{G}_ai\widehat{V}_a\right)\widehat{\overline{\xi }}{\displaystyle \underset{a=1}{\overset{N}{}}}\widehat{q}_a^{}\widehat{\overline{\xi }}_a,`$ (52) where $`\widehat{V}_a\widehat{V}_a(\theta _1,\mathrm{},\theta _N)`$ and $`\widehat{G}_a=ih\text{-}/\theta _a+h\text{-}\alpha _a`$ in the $`\widehat{W}`$-diagonal representation. Let us assume $$[\widehat{G}_a,\widehat{W}_b]=h\text{-}\delta _{ab}\widehat{W}_b,\{\widehat{\xi }_a,\widehat{\overline{\xi }}_b\}=\delta _{ab},\{\widehat{\xi }_a,\widehat{\xi }_b\}=0,\{\widehat{\overline{\xi }}_a,\widehat{\overline{\xi }}_b\}=0.$$ (53) Then, the Hamiltonian following from these fermionic operators is $$\widehat{H}=\underset{a=1}{\overset{N}{}}\frac{1}{2mR_a^2}\widehat{G}_a\widehat{G}_a+\widehat{V}_a\widehat{V}_a\underset{a,b=1}{\overset{N}{}}\frac{i}{\sqrt{2m}R_a}[\widehat{G}_a,\widehat{V}_b][\widehat{\xi }_a,\widehat{\overline{\xi }}_b].$$ (54) The Hamiltonian (54) may describe the supersymmetric quantum mechanics on the torus $`T^N=S^1\mathrm{}S^1`$. As before, $`\alpha _a(a=1,\mathrm{},N)`$ may be interpreted as the gauge field appearing as a consequence of the quantization on each topological space $`S^1`$. It is difficult to obtain the exact form of the supersymmetric ground state wave functions of the model. If we, however, restrict ourselves to certain sectors of the model, they can be obtained in closed form such as Eq. (57) . In order to see this, let us define $$\|\|0,\|+\underset{a=1}{\overset{N}{}}\widehat{\overline{\xi }}_a\|0\mathrm{with}\widehat{\xi }_a\|0=0(a=1,\mathrm{},N).$$ (55) Then, $`\widehat{Q}\|=\widehat{\overline{Q}}\|+=0`$ is trivially satisfied, so that in these two sectors, the supersymmetric ground state (zero-energy state) wave functions are obtained in closed form by solving simple first-order equations such as Eqs. (40). Aside from the normalization, the solutions are obtained as $`\mathrm{\Psi }^+(\theta _1,\mathrm{},\theta _N)`$ $`=`$ $`\mathrm{exp}{\displaystyle \underset{a=1}{\overset{N}{}}}\left(i\alpha _a\theta _a+{\displaystyle \frac{\sqrt{2m}R_a}{h\text{-}}}{\displaystyle ^{\theta _a}}𝑑\overline{\theta }_aV_a(\overline{\theta }_1,\mathrm{},\overline{\theta }_a,\mathrm{},\overline{\theta }_N)\right)\|+,`$ $`\mathrm{\Psi }^{}(\theta _1,\mathrm{},\theta _N)`$ $`=`$ $`\mathrm{exp}{\displaystyle \underset{a=1}{\overset{N}{}}}\left(i\alpha _a\theta _a{\displaystyle \frac{\sqrt{2m}R_a}{h\text{-}}}{\displaystyle ^{\theta _a}}𝑑\overline{\theta }_aV_a(\overline{\theta }_1,\mathrm{},\overline{\theta }_a,\mathrm{},\overline{\theta }_N)\right)\|.`$ (56) The wave functions have to satisfy the periodic boundary condition $`\mathrm{\Psi }^\pm (\mathrm{},\theta _a+2\pi ,\mathrm{})=\mathrm{\Psi }^\pm (\mathrm{},\theta _a,\mathrm{})(a=1,\mathrm{},N)`$. If the contributions coming from the superpotential in Eqs. (56) do not spoil the normalizablity and the periodicity of the wave functions, the supersymmetry can be broken for noninteger values of $`\alpha _a(a=1,\mathrm{},N)`$. ## 3 Conclusions and Discussion We have applied the Ohknuki-Kitakado formulation of quantum mechanics on $`S^1`$ to the supersymmetric simple pendulum whose Hamiltonian is given by Eq. (37) and satisfies the algebra (1). According to their formulation, an undetermined parameter, which can be interpreted as a constant gauge field, inevitably enters into the theory to yield the various quantum mechanics on $`S^1`$. We have studied the effects of the quantization ambiguity on the supersymmetric ground state wave functions of the model. We have found that supersymmetry can be broken due to the existence of the gauge field $`\alpha `$. The gauge field twists the boundary condition of the supersymmetric ground state wave functions. For noninteger values of $`\alpha `$, they do not satisfy the required periodic boundary condition. As a result, they become unphysical wave functions though they are normalizable. The mechanism of supersymmetry breaking is different from the usual supersymmetry breaking discussed in supersymmetric quantum mechanics. The latter depends crucially on the structure, the leading term, of the superpotential, while the former is entirely due to the quantization ambiguity resulting firmly from quantization on a space with nontrivial topology like $`S^1`$. We have chosen the superpotential $`V(\widehat{W},\widehat{W}^{})`$ in such a way that it becomes a simple pendulum in the classical limit. In principle, one can choose any superpotential as long as it can be written in terms of integer powers of the operators $`\widehat{W}`$ and $`\widehat{W}^{}`$. Thanks to the factorizable property for finding the supersymmetric ground state wave functions, they are given simply by solving the first order equation (31) and are obtained in closed form $$\psi _{\pm \frac{1}{2}}(\theta )=\mathrm{exp}\left(i\alpha \theta \frac{\sqrt{2m}R}{h\text{-}}^\theta 𝑑\overline{\theta }V(\mathrm{e}^{i\overline{\theta }},\mathrm{e}^{i\overline{\theta }})\right),$$ (57) for the general superpotential. Our mechanism of supersymmetry breaking is not altered by the choice of the superpotential if $`\mathrm{exp}(^\theta 𝑑\overline{\theta }V(\mathrm{e}^{i\overline{\theta }},\mathrm{e}^{i\overline{\theta }}))`$ does not violate the periodicity and the normalizability of the wave functions, which is the case for the superpotential satisfying our criterion. Supersymmetry breaking will always occur for noninteger values of $`\alpha `$. Because of the factorizable property, there is no way to prevent the gauge field from entering into the supersymmetric ground state wave functions and twisting their boundary conditions. The gauge field cannot be removed by a regular gauge transformation. One may wonder whether all the eigenfunctions of the Hamiltonian (37) become unphysical, that is, those that do not satisfy the periodic boundary condition due to the existence of the gauge field. This is not true. In order to see this, let us consider a free Hamiltonian, ignoring all terms except for $`\widehat{G}^2`$. The energy eigenvalue depends on $`\alpha `$ like $`(m+\alpha )^2`$ and the gauge field produces an effect on the observable at the quantum level . The corresponding eigenfunction satisfying the periodic boundary condition is easily found to be $`\mathrm{e}^{im\theta }`$. The ground state wave functions and the other eigenfunctions are obtained by solving essentially different types of differential equations in the system satisfying the supersymmetry algebra (1). We have also discussed the limit of $`R\mathrm{}`$. One-dimensional Euclidean space is realized. In the limit with the relation (46), we have obtained the supersymmetric harmonic oscillator with angular frequency $`\omega =\sqrt{g_N/R}`$. There exists one physical supersymmetric ground state wave function. The other one, though it is a zero-energy state, is unphysical because of its non-normalizability. These two wave functions are actually obtained by taking the limit in the solutions of the zero-energy wave functions (41). In the limit all the effects of the gauge field $`\alpha `$ disappear. Then, an infinite number of inequivalent representations is reduced to a unique representation, which is nothing but the representation of the canonical algebra $`[\widehat{p},\widehat{x}]=ih\text{-}`$ . The fermionic operators become supercharges in the same limit, and they generate supersymmetry transformations between bosons $`(\widehat{x})`$ and fermions $`(\widehat{\xi },\widehat{\overline{\xi }})`$. We have also considered the $`N`$-component generalization of Eq. (25) and studied the supersymmetric ground state wave functions of the model. The Hamiltonian (54) describes the supersymmetric quantum mechanics on the torus $`T^N`$. If we restrict ourselves to the two sectors given by $`\|`$ and $`\|+`$, the wave functions can be obtained in closed form (56) by solving the simple first-order equations. We have found, again, that supersymmetry was able to be broken due to the existence of the gauge field $`\alpha _a\mathrm{integer}(a=1,\mathrm{},N)`$, which appeared as a consequence of the quantization on the each topological space $`S^1`$. Finally, let us discuss the similarities between our mechanism of supersymmetry breaking and that through boundary conditions of fields for compactified directions in supersymmetric quantum field theory. Strictly speaking, supersymmetry breaking through boundary conditions is one thing, and that through our mechanism is another. Nevertheless, it may be interesting to discuss the similarities between the two supersymmetry breakings. In the former case, the breaking means that the action is no longer invariant under the supersymmetry transformations. But it does not necessarily mean the nonexistence of the zero-energy state in the system. The supersymmetric harmonic oscillator at finite temperature is one of the examples in which the action is not invariant under the supersymmetry transformations because of the different boundary conditions between the bosons and the fermions; there exists, however, the zero-energy state in the system, which results from $`\mathrm{Tr}(1)^{\widehat{f}}=1`$ . In the latter case, we assume that the Hamiltonian satisfies the supersymmetry algebra (1), so that supersymmetry breaking immediately means the nonexistence of physical zero-energy states in the system, and whether or not supersymmetry is broken is determined definitely by the existence of the zero-energy state. As seen from Eq. (42), the boundary condition for the $`S^1`$ direction is twisted for $`0<\alpha <1`$ by $`\mathrm{e}^{i2\pi \alpha }`$. If we consider a theory at finite temperature, it is equivalent to studying the theory in a space where the Euclidean time direction is compactified on $`S^1`$. It is well known that supersymmetry is broken at finite temperature by the different boundary conditions for the Euclidean time direction between the bosons and the fermions. The boson (fermion) satisfies the (anti)periodic boundary condition. The case of $`\alpha =1/2`$, which actually corresponds to the antiperiodic boundary condition, is similar to the case of supersymmetry breaking at finite temperature. More generally, if one wishes to break supersymmetry through different boundary conditions between the bosons and the fermions such as the finite temperature case, one can use the boundary condition associated with the $`U(1)_R`$ symmetry , in which the $`U(1)_R`$ charges are different between the bosons and the fermions in a supermultiplet. If we regard the factor $`\mathrm{e}^{i2\pi \alpha }`$ in Eq. (42) as a boundary condition that breaks supersymmetry, our mechanism of supersymmetry breaking is quite similar to that through boundary condition associated with the $`U(1)_R`$ symmetry. If one takes this similarity seriously, one says that a possible physical origin of supersymmetry breaking through the boundary condition associated with the $`U(1)_R`$ symmetry has been found. Needless to say, in order to confirm this statement, we need to clarify how the quantization ambiguity is realized in (supersymmetric) quantum field theory . Acknowledgments I would like to thank Professor Jan Ambj$`\varphi `$rn for fruitful discussions, encouraging me, and critical comments on this manuscript, and the Niels Bohr Institute for warm hospitality. I would also like to thank Professor Shogo Tanimura (Kyoto University) for valuable discussions on the quantum mechanics on $`S^1`$.
warning/0003/cond-mat0003075.html	ar5iv	text	# Quasi-Spin-Charge Separation and the Spin Quantum Hall Effect ## I Introduction Recently a number of works have appeared investigating a new universality class of delocalization transition referred to as the spin quantum Hall effect. This transition can occur in certain dirty superconductors with unbroken $`su(2)`$ spin-rotation symmetry. Kagalovsky et. al. constructed an $`su(2)`$ invariant network model for the transition and numerically determined some critical exponents. Senthil et. al. on the other hand modeled the phenomena with a supersymmetric spin chain. Remarkably, the network model was mapped directly on the lattice onto classical percolation by Gruzberg et. al. and exact exponents were computed which agreed very well with the numerical simulations of the super spin chain. In the quantum field theory approach to delocalization transitions, the computation of critical exponents is normally a difficult strong-coupling problem requiring the existence of non-trivial infrared fixed points. The claim that for the spin quantum Hall effect the infrared fixed point of the disorder averaged effective field theory is simply percolation is rather unexpected and for this reason we set out to understand this using quantum field theory methods. Our starting point is the hamiltonian formulation of the network model given in . We carry out the disorder averaging using the supersymmetric method in conjunction with conformal field theory methods. This leads to an effective action which consists of a conformal field theory with an $`osp(4\|4)`$ super-current algebra symmetry perturbed by certain marginal operators which are bilinear in the supercurrents. The one-loop renormalization group (RG) $`\beta `$eta functions we compute for the three independent couplings appear to be as complicated as those for the usual quantum Hall transition. There exists a fine-tuning of the network model couplings wherein two of the couplings are essentially identified and the resulting model is remarkably simpler and can be solved. This is due in part to the fact that the stress tensor for the $`osp(4\|4)`$ current algebra can be written as the sum of the stress tensors for the $`osp(2\|2)`$ level $`2`$ and $`su(2)`$ level $`0`$ current algebras corresponding to the charge and spin degrees of freedom respectively. This “spin-charge separation” is present in the disorder-averaged effective theory and this leads to two decoupled $`\beta `$eta functions. This simplification allows us to determine the non-trivial infrared fixed point: the $`su(2)`$ level $`0`$ degrees of freedom decouple in the flow and the infrared conformal field theory is simply the coset $`osp(4\|4)_1/su(2)_0`$. In this way we recover some of the percolation exponents predicted in . Due to the logarithmic nature of the above conformal field theories, we find that in spite of the separation of the stress tensor into commuting pieces, $`T_{osp(4\|4)}=T_{osp(2\|2)}+T_{su(2)}`$, the Hilbert space does not factorize. This leads to the peculiar result that the coset $`osp(4\|4)_1/su(2)_0`$ is not equivalent to the $`osp(2\|2)_2`$ current algebra theory, even though this coset possesses the current algebra as a symmetry and the conformal dimensions of the coset are the same as for the $`osp(2\|2)_2`$ current algebra. In the last section we study the possible universality classes based on the 1-loop RG equations and suggest that the network model is universally attracted to a so-called “strange direction”. ## II The Models Kagalovsky et. al. gave a hamiltonian formulation of their network model. The result is the $`4\times 4`$ matrix hamiltonian: $$H=(\tau _xp_x+\tau _zp_y)1+1\stackrel{}{\alpha }\stackrel{}{\sigma }$$ (1) where $`\stackrel{}{\tau },\stackrel{}{\sigma }`$ are two copies of the Pauli matrices, $`p_{x,y}=i_{x,y}`$, and $`\stackrel{}{\alpha }`$ is a random spin potential. This hamiltonian has the defining properties to belong the class C of the classification introduced in ref.. Let us perform a unitary transformation $`H(U^{}1)H(U1)`$, where $`U`$ corresponds to a rotation about the $`x`$ axis for the Pauli matrices: $`U^{}\tau _zU=\tau _y`$, $`U^{}\tau _yU=\tau _z`$, $`U^{}\tau _xU=\tau _x`$. After including $`su(2)`$ gauge potentials to $`\stackrel{}{p}`$, one has the following $`2\times 2`$ block structure: $$H=\left(\begin{array}{cc}\stackrel{}{\alpha }\stackrel{}{\sigma }& i_{\overline{z}}+A_{\overline{z}}\\ i_z+A_z& \stackrel{}{\alpha }\stackrel{}{\sigma }\end{array}\right)$$ (2) where $`_z=_xi_y`$, $`_{\overline{z}}=_x+i_y`$, $`A_z=A_xiA_y`$, $`A_{\overline{z}}=A_x+iA_y`$. In the above hamiltonian, $`\stackrel{}{\alpha }(x,y)`$ is a real, random spin potential, and $`A_\mu =_aA_\mu ^a(x,y)\sigma ^a`$ are random $`su(2)`$ gauge potentials, with $`A_{x,y}^a`$ real. In the sequel we will take them to have the following gaussian distributions $`P(\stackrel{}{\alpha })`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{1}{g_\alpha }}{\displaystyle \frac{d^2x}{2\pi }\stackrel{}{\alpha }(x)\stackrel{}{\alpha }(x)}\right)`$ (3) $`P(A_\mu )`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{4}{g_A}}{\displaystyle \frac{d^2x}{2\pi }A_z^a(x)A_{\overline{z}}^a(x)}\right)`$ (4) As we will show in the next section, renormalization of the effective action obtained upon disorder averaging leads to an additional interaction which can be viewed as arising from a random mass $`m(x,y)`$. The complete hamiltonian which leads to a renormalizable effective action is then $$H=\left(\begin{array}{cc}\stackrel{}{\alpha }\stackrel{}{\sigma }+m& i_{\overline{z}}+A_{\overline{z}}\\ i_z+A_z& \stackrel{}{\alpha }\stackrel{}{\sigma }m\end{array}\right)$$ (5) We will take $`m(x)`$ to have the gaussian distribution $$P(m)=\mathrm{exp}\left(\frac{1}{g_m}\frac{d^2x}{2\pi }m(x)^2\right)$$ (6) Note that since $`_z^{}=_{\overline{z}}`$, $`A_z^{}=A_{\overline{z}}`$, the hamiltonian is hermitian if $`m,\stackrel{}{\alpha },A_{x,y}`$ are real. In this situation the couplings $`g_\alpha ,g_m,g_A`$ are the variances of normalizable gaussian distributions if they are real and positive. Negative couplings $`g_{\alpha ,m,A}`$ can be interpreted as corresponding to imaginary random potentials; hermitian hamiltonians can then be constructed by doubling the number of degrees of freedom, as in . The single particle Green functions are defined by the functional integral $`Z^1D\mathrm{\Psi }^{}D\mathrm{\Psi }\mathrm{exp}(S)`$ with $`Z`$ the partition function and $$S=\frac{d^2x}{2\pi }\mathrm{\Psi }^{}(x)i\left(H\right)\mathrm{\Psi }(x)$$ (7) where $`=E+i\epsilon `$. For $`\epsilon =0^+`$, this defines the retarded Green function $$G_R(x,x^{};E)=\underset{\epsilon 0^+}{lim}x\|\frac{1}{H(E+i\epsilon )}\|x^{}=\underset{\epsilon 0^+}{lim}i\mathrm{\Psi }(x)\mathrm{\Psi }^{}(x^{})$$ (8) Letting $$\mathrm{\Psi }=\left(\begin{array}{c}\overline{\psi }_+\\ \psi _+\end{array}\right),\mathrm{\Psi }^{}=(\psi _{},\overline{\psi }_{})$$ (9) where $`\psi _+`$ is a 2-component fermion $`\psi _+^i,i=1,2`$, and similarly for $`\overline{\psi }_\pm `$, one finds $`S`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d^2x}{2\pi }}(\psi _{}(_{\overline{z}}+iA_{\overline{z}})\psi _++\overline{\psi }_{}(_z+iA_z)\overline{\psi }_++i\stackrel{}{\alpha }(\psi _{}\stackrel{}{\sigma }\overline{\psi }_++\overline{\psi }_{}\stackrel{}{\sigma }\psi _+)`$ (10) $`+im(\psi _{}\overline{\psi }_+\overline{\psi }_{}\psi _+)i\mathrm{\Phi }_E)`$ (11) where $$\mathrm{\Phi }_E=\psi _{}\overline{\psi }_++\overline{\psi }_{}\psi _+$$ (We have suppressed the $`su(2)`$ indices, i.e. $`\psi _{}\overline{\psi }_+=_i\psi _{}^i\overline{\psi }_+^i`$, etc.) In the perturbation $`H_ϵ=ϵ\tau _y1`$ was added and critical exponents for $`ϵ`$ were measured numerically. Performing the unitary transformation $`U`$ defined above, this results in the following perturbation of the action: $$S_ϵ=iϵ\frac{d^2x}{2\pi }\mathrm{\Phi }_ϵ,\mathrm{\Phi }_ϵ=\psi _{}\overline{\psi }_+\overline{\psi }_{}\psi _+$$ (12) Non-zero $`ϵ`$ corresponds to a non-zero average random mass $`m`$. Let us attempt to compare this with the model considered by Senthil et. al.. There, one had a one-dimensional lattice in the $`x`$-direction with sites labeled by $`j`$, and a continuous $`y`$ direction. At each site there are fermionic degrees of freedom $`\chi _j(y)`$, where $`j`$ even corresponds to left-movers and $`j`$ odd to right-movers. The hamiltonian is $`H`$ $`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle }dy(()^j\chi ^{}_j(y)(i_y+\stackrel{}{\eta }_j(y)\stackrel{}{\sigma })\chi _j(y)it^0_j(y)(\chi _{j+1}^{}\chi _j\chi _j^{}\chi _{j+1})`$ (13) $`+\stackrel{}{t}_j(y)(\chi _{j+1}^{}\stackrel{}{\sigma }\chi _j+\chi _j^{}\stackrel{}{\sigma }\chi _{j+1}))`$ (14) Taking a continuum limit, $`\chi _j(y)\psi _+(x,y)`$ for $`j`$ even and $`\chi _j(y)\overline{\psi }_+(x,y)`$ for $`j`$ odd, $`\stackrel{}{\eta }_j(y)\stackrel{}{\eta }_y(x,y)`$, and $`t_j(y)t(x,y)`$ , one finds $`iH`$ $`=`$ $`{\displaystyle }dxdy(\psi _{}(_y+i\stackrel{}{\eta }_y\stackrel{}{\sigma })\psi _+\overline{\psi }_{}(_y+i\stackrel{}{\eta }_y\stackrel{}{\sigma })\overline{\psi }_++t_0(\psi _{}\overline{\psi }_+\overline{\psi }_{}\psi _+)`$ (15) $`+i\stackrel{}{t}(\psi _{}\stackrel{}{\sigma }\overline{\psi }_++\overline{\psi }_{}\stackrel{}{\sigma }\psi _+))`$ (16) Rotational invariance of the kinetic terms can be restored by adding $`_x+i\stackrel{}{\eta }_x\stackrel{}{\sigma }`$ to the derivatives. Performing a rotation to euclidean space $`yiy`$, one then obtains the action (11) with $`\stackrel{}{\alpha }=\stackrel{}{t}`$, $`A_z=(\stackrel{}{\eta }_xi\stackrel{}{\eta }_y)\stackrel{}{\sigma }`$ and $`t_0=im`$. In , $`\stackrel{}{t}`$ and $`t_0`$ were taken as real gaussian distributed but with the same variance. For the model defined by (11), this corresponds to imaginary $`m`$. Letting $`mim`$, one sees that for the effective theory obtained after disorder averaging, imaginary $`m`$ corresponds to negative $`g_m`$. Since the variances are the same, this corresponds to $`g_\alpha =g_m`$.This appears to be inconsistent with hermiticity, since the hamiltonian (5) is not hermitian for imaginary $`m`$, suggesting the two models cannot be naively compared in this way. In summary, though we have not described an exact mapping between the network model in and the super spin chain in , it appears that they should be related on the line $`g_\alpha +g_m=0`$. In fact, one obtains directly the Heisenberg type superspin chain upon disorder averaging only along this line. We will provide further support of this statement based on symmetry in the sequel. As we will show, this line has some rather special properties which allow the model to be solved. ## III Effective Action and Renormalization Group Analysis ### A Supersymmetric Disorder Averaging Since we are dealing with a free field theory, the supersymmetric method for disorder averaging can be used. Conformal field theory techniques in conjunction with this method was used for other models in . We augment the theory with bosonic ghosts $`\beta _\pm ^i,\overline{\beta }_\pm ^i`$, $`i=1,2`$, so that the inverse of the fermionic partition function $`Z(\stackrel{}{\alpha },m,A)=D\psi e^{S(\psi )}`$ is represented as a bosonic functional integral: $$Z(\stackrel{}{\alpha },m,A)^1=D\beta e^{S(\psi \beta )}$$ (17) One can then perform the gaussian integrals over the random potentials. The result is the effective action: $$S_{eff}=S_{cft}+\frac{d^2x}{2\pi }\left(g_\alpha 𝒪_\alpha +g_m𝒪_m+g_A𝒪_A\right)$$ (18) The conformal field theory $`S_{cft}`$ has Virasoro central charge $`c=0`$ and the action: $$S_{cft}=\frac{d^2x}{2\pi }\left(\psi _{}_{\overline{z}}\psi _++\overline{\psi }_{}_z\overline{\psi }_++\beta _{}_{\overline{z}}\beta _++\overline{\beta }_{}_z\overline{\beta }_+\right)$$ (19) The first-order action for the bosonic ghosts of conformal dimension $`\pm 1/2`$ can be treated as in . The operators which perturb away from the conformal field theory are: $`𝒪_\alpha `$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(\psi _{}\stackrel{}{\sigma }\overline{\psi }_++\overline{\psi }_{}\stackrel{}{\sigma }\psi _++\beta _{}\stackrel{}{\sigma }\overline{\beta }_++\overline{\beta }_{}\stackrel{}{\sigma }\beta _+\right)^2`$ (20) $`𝒪_m`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(\psi _{}\overline{\psi }_+\overline{\psi }_{}\psi _++\beta _{}\overline{\beta }_+\overline{\beta }_{}\beta _+\right)^2`$ (21) $`𝒪_A`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(\psi _{}\stackrel{}{\sigma }\psi _++\beta _{}\stackrel{}{\sigma }\beta _+\right)\left(\overline{\psi }_{}\stackrel{}{\sigma }\overline{\psi }_++\overline{\beta }_{}\stackrel{}{\sigma }\overline{\beta }_+\right)`$ (22) ### B $`\beta `$eta Functions The one-loop renormalization group $`\beta `$eta functions can be deduced from the operator product expansions of the marginal perturbing operators $`𝒪_i`$. Namely, if $$𝒪_i(z,\overline{z})𝒪_j(0)\frac{1}{z\overline{z}}C_{ij}^k𝒪_k(0)+\mathrm{}.$$ (23) then to lowest order the $`\beta `$eta function $`\beta _k=dg_k/d\mathrm{log}l`$, where $`l`$ is a length scale, is given by $$\beta _k=\underset{i,j}{}C_{ij}^kg_ig_j$$ (24) We normalize the Pauli matrices as follows: $`\stackrel{}{\sigma }\stackrel{}{\sigma }=(\sigma ^3)^2+\sigma ^+\sigma ^{}+\sigma ^{}\sigma ^+`$, with $$\sigma ^3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\sigma ^+=\sqrt{2}\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),\sigma ^{}=\sqrt{2}\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)$$ (25) With this normalization, one has $`Tr\sigma ^a\sigma ^b=2\delta ^{ab}`$, and $`f^{abc}f^{abd}=8\delta ^{cd}`$, where $`[\sigma ^a,\sigma ^b]=f^{abc}\sigma ^c`$. One also needs the operator products: $`\psi _+(z)\psi _{}(0)\psi _{}(z)\psi _+(0)1/z`$ (26) $`\beta _+(z)\beta _{}(0)\beta _{}(z)\beta _+(0)1/z`$ (27) and similarly for the right-movers $`\overline{\psi },\overline{\beta }`$ with $`z`$ replaced by $`\overline{z}`$. One then finds $`𝒪_A(z,\overline{z})𝒪_A(0)`$ $``$ $`{\displaystyle \frac{2}{z\overline{z}}}𝒪_A,𝒪_m(z,\overline{z})𝒪_m(0){\displaystyle \frac{2}{z\overline{z}}}𝒪_m`$ (28) $`𝒪_\alpha (z,\overline{z})𝒪_\alpha (0)`$ $``$ $`{\displaystyle \frac{1}{z\overline{z}}}\left(2𝒪_\alpha 8𝒪_A\right)`$ (29) $`𝒪_A(z,\overline{z})𝒪_\alpha (0)`$ $``$ $`{\displaystyle \frac{1}{z\overline{z}}}\left({\displaystyle \frac{1}{2}}𝒪_\alpha {\displaystyle \frac{3}{2}}𝒪_m\right)`$ (30) $`𝒪_A(z,\overline{z})𝒪_m(0)`$ $``$ $`{\displaystyle \frac{1}{z\overline{z}}}\left({\displaystyle \frac{3}{2}}𝒪_m{\displaystyle \frac{1}{2}}𝒪_\alpha \right)`$ (31) $`𝒪_\alpha (z,\overline{z})𝒪_m(0)`$ $``$ $`{\displaystyle \frac{1}{z\overline{z}}}\left(3𝒪_m𝒪_\alpha 4𝒪_A\right)`$ (32) The resulting $`\beta `$eta functions are $`\beta _\alpha `$ $`=`$ $`2g_\alpha ^2+2g_\alpha g_m+g_A(g_\alpha +g_m)`$ (33) $`\beta _m`$ $`=`$ $`2g_m^26g_\alpha g_m+3g_A(g_\alpha +g_m)`$ (34) $`\beta _A`$ $`=`$ $`2g_A^2+8g_\alpha (g_\alpha +g_m)`$ (35) If one starts with a model with only the random spin potential $`\stackrel{}{\alpha }`$, i.e. only $`g_\alpha 0`$, then $`g_A`$ and $`g_m`$ are generated under renormalization. Note that the $`\beta `$eta functions simplify dramatically when $`g_\alpha +g_m=0`$, which is the subject of the next section. We will return to a general analysis of the possible universality classes contained in these $`\beta `$eta functions in section V. ## IV Spin-Charge Separation ### A Symmetries of the effective action The free conformal action $`S_{cft}`$ has a maximal $`osp(4\|4)_1`$ current algebra symmetry. The conformal currents $`J^A`$ generating this symmetry correspond to all possible bilinears in fermions and ghosts. Suppressing the possible index structure, $`J^A=\{\psi \psi ,\psi \beta ,\beta \beta \}`$. The perturbing operators are bilinear in these currents: $`𝒪_\alpha =d_{AB}^{(\alpha )}J^A\overline{J}^B`$ for some tensor $`d_{AB}^{(\alpha )}`$ and similarly for $`𝒪_m`$ and $`𝒪_A`$. For arbitrary $`g_\alpha ,g_m,g_A`$, the $`osp(4\|4)`$ symmetry of $`S_{cft}`$ is broken but an $`osp(2\|2)`$ symmetry is preserved. To describe the symmetry, let us define the left-moving currents: $`J`$ $`=`$ $`{\displaystyle \underset{i}{}}\psi _+^i\psi _{}^i,J_\pm ={\displaystyle \underset{i,j}{}}ϵ_{ij}\psi _\pm ^i\psi _\pm ^j`$ (36) $`H`$ $`=`$ $`{\displaystyle \underset{i}{}}\beta _+^i\beta _{}^i,S_\pm ={\displaystyle \underset{i}{}}\psi _\pm ^i\beta _{}^i`$ (37) $`\widehat{S}_\pm `$ $`=`$ $`{\displaystyle \underset{i,j}{}}ϵ_{ij}\psi _\pm ^i\beta _\pm ^j`$ (38) where $`ϵ_{ij}=ϵ_{ji}`$, $`ϵ_{12}=1`$, and similarly for the right-movers $`\overline{J}=\overline{\psi }_+\overline{\psi }_{},\overline{H}=\overline{\beta }_+\overline{\beta }_{},\mathrm{}..`$. These currents generate an $`osp(2\|2)_k`$ current algebra at level $`k=2`$. The non-zero operator products are $`J(z)J(0)`$ $``$ $`{\displaystyle \frac{k}{z^2}},H(z)H(0){\displaystyle \frac{k}{z^2}}`$ (39) $`J(z)J_\pm (0)`$ $``$ $`\pm {\displaystyle \frac{2}{z}}J_\pm ,J_+(z)J_{}(0){\displaystyle \frac{2k}{z^2}}{\displaystyle \frac{4}{z}}J`$ (40) $`J(z)S_\pm (0)`$ $``$ $`\pm {\displaystyle \frac{1}{z}}S_\pm ,J(z)\widehat{S}_\pm (0)\pm {\displaystyle \frac{1}{z}}\widehat{S}_\pm `$ (41) $`H(z)S_\pm (0)`$ $``$ $`\pm {\displaystyle \frac{1}{z}}S_\pm ,H(z)\widehat{S}_\pm (0){\displaystyle \frac{1}{z}}\widehat{S}_\pm `$ (42) $`J_\pm (z)S_{}(0)`$ $``$ $`{\displaystyle \frac{2}{z}}\widehat{S}_\pm ,J_\pm (z)\widehat{S}_{}(0){\displaystyle \frac{2}{z}}S_\pm `$ (43) $`S_\pm (z)\widehat{S}_\pm (0)`$ $``$ $`\pm {\displaystyle \frac{1}{z}}J_\pm `$ (44) $`S_+(z)S_{}(0)`$ $``$ $`{\displaystyle \frac{k}{z^2}}+{\displaystyle \frac{1}{z}}(HJ)`$ (45) $`\widehat{S}_+(z)\widehat{S}_{}(0)`$ $``$ $`{\displaystyle \frac{k}{z^2}}+{\displaystyle \frac{1}{z}}(H+J)`$ (46) The currents $`J`$ and $`J_\pm `$ generate a charged $`su(2)`$ subalgebra under which the fermions $`\psi _\pm ^i`$ transform as doublets. The $`osp(2\|2)`$ symmetry is present for arbitrary potentials before disorder averaging, as we now describe. The fermionic generators define nilpotent transformations of the bosonic and fermionic fields. The action (11) is then an exact variation with respect to these transformations, making then clear the fact that this $`osp(2\|2)`$ symmetry is preserved. For example the generator $`S_+`$ induces the transformation $`\delta \psi _+^i=0`$, $`\delta \psi _{}^i=\beta _{}^i`$, $`\delta \beta _+^i=\psi _+^i`$, $`\delta \beta _{}^i=0`$. The symmetry acts left-right diagonally, so for the right movers $`\delta \overline{\psi }_+^i=0,\delta \overline{\psi }_{}^i=\overline{\beta }_{}^i`$, $`\delta \overline{\beta }_+^i=\overline{\psi }_+^i,\delta \overline{\beta }_{}^i=0`$. The action (11) may be written as $$S=S_{cft}+\delta \frac{d^2x}{2\pi }\mathrm{\Theta }$$ where $$\mathrm{\Theta }=im(\psi _{}\overline{\beta }_+\overline{\psi }_{}\beta _+)+i\stackrel{}{\alpha }(\psi _{}\stackrel{}{\sigma }\overline{\beta }_++\overline{\psi }_{}\stackrel{}{\sigma }\beta _+)+i\stackrel{}{A}(\psi _{}\stackrel{}{\sigma }\beta _++\overline{\psi }_{}\stackrel{}{\sigma }\overline{\beta }_+)$$ (47) Hence $`\delta S=0`$ since $`\delta S_{cft}=0`$ and $`\delta ^2=0`$. This holds similarly, but with another operator $`\widehat{\mathrm{\Theta }}`$, for $`\widehat{S}_{}`$ which generates the transformations $`\widehat{\delta }\psi _+^i=ϵ^{ij}\beta _{}^j`$, $`\widehat{\delta }\psi _{}^i=0`$ and $`\widehat{\delta }\beta _+^i=ϵ^{ij}\psi _{}^j`$, $`\widehat{\delta }\beta _{}^i=0`$. This $`osp(2\|2)`$ symmetry would not be preserved if we had chosen the random gauge potential in $`u(2)`$ or if had added extra scalar potential randomness. This implies that after disorder averaging the perturbations by $`𝒪_\alpha `$, $`𝒪_m`$ and $`𝒪_A`$ in the effective theory preserve a global, left-right diagonal, $`osp(2\|2)`$ symmetry. The current conservation law takes the left-right diagonal form $$_{\overline{z}}J^a+_z\overline{J}^a=0$$ (48) for $`J^a,\overline{J}^a`$ any of the eight $`osp(2\|2)_2`$ currents. Let us describe more explicitly how the $`osp(2\|2)`$ symmetry is manifested in the effective theory. In the sequel we will set $`g_\alpha =g_m=g`$ and the effective action will contain the operator $`𝒪_g𝒪_\alpha 𝒪_m`$. Many terms cancel in the combination $`𝒪_\alpha 𝒪_m`$ and the result can be written as an $`osp(2\|2)`$ current-current perturbation. This simplification is analogous to the $`g_V+g_M=0`$ line for the random $`U(1)`$ fermions which has $`gl(1\|1)`$ symmetry, studied in . By repeated use of the identity $$\sigma _{ij}^a\sigma _{nm}^a=2\delta _{im}\delta _{jn}\delta _{ij}\delta _{nm}$$ (49) one can express $`𝒪_g`$ in terms of the above $`osp(2\|2)`$ currents: $`𝒪_g`$ $``$ $`𝒪_\alpha 𝒪_m`$ (50) $`=`$ $`J\overline{J}+H\overline{H}+{\displaystyle \frac{1}{2}}(J_{}\overline{J}_++J_+\overline{J}_{})+S_{}\overline{S}_+S_+\overline{S}_{}\widehat{S}_{}\overline{\widehat{S}}_++\widehat{S}_+\overline{\widehat{S}}_{}`$ (51) The operator $`𝒪_g`$ has the structure of the quadratic Casimir for $`osp(2\|2)`$. Namely, $`𝒪_g=_{a,b}C_{ab}J^a\overline{J}^b`$, where $`J^a`$ are $`osp(2\|2)`$ currents and $`C_{ab}`$ corresponds to the quadratic Casimir. Thus $`𝒪_g`$ is $`osp(2\|2)`$ invariant. The operator $`𝒪_\alpha +𝒪_m`$ cannot be written only in terms of the $`osp(2\|2)`$ currents; one needs some of the $`osp(4\|4)`$ currents. Let $`B_\pm ^{ij}=\beta _\pm ^i\beta _\pm ^j`$ ; $`\stackrel{}{L}_f=\psi _{}\stackrel{}{\sigma }\psi _+;\stackrel{}{L}_b=\beta _{}\stackrel{}{\sigma }\beta _+`$ (52) $`U_\pm ^{ij}={\displaystyle \frac{1}{2}}(\beta _\pm ^i\psi _\pm ^j+\beta _\pm ^j\psi _\pm ^i)`$ ; $`\stackrel{}{V}_{}=\psi _{}\stackrel{}{\sigma }\beta _+;\stackrel{}{V}_+=\beta _{}\stackrel{}{\sigma }\psi _+`$ (53) Then, $$𝒪_\alpha +𝒪_m=\frac{1}{2}𝒪_g+\stackrel{~}{𝒪}$$ with $`\stackrel{~}{𝒪}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(B_{}^{ij}\overline{B}_+^{ij}+B_+^{ij}\overline{B}_{}^{ij}\right)+{\displaystyle \frac{1}{2}}\left(\stackrel{}{L}_f\stackrel{}{\overline{L}}_f\stackrel{}{L}_b\stackrel{}{\overline{L}}_b\right)`$ (54) $`+U_{}^{ij}\overline{U}_+^{ij}U_+^{ij}\overline{U}_{}^{ij}+{\displaystyle \frac{1}{2}}\left(\stackrel{}{V}_+\stackrel{}{\overline{V}}_{}\stackrel{}{V}_{}\stackrel{}{\overline{V}}_+\right)`$ (55) We have already seen that $`𝒪_g`$ is $`osp(2\|2)`$ invariant. One can check explicitly that $`\stackrel{~}{𝒪}`$ is also $`osp(2\|2)`$ invariant. Let us turn now to the perturbation $`𝒪_A`$. This operator as defined is a left-right current-current perturbation: $$𝒪_A=\underset{a=1}{\overset{3}{}}L^a\overline{L}^a$$ (56) where $$L^a=L_f^a+L_b^a=\psi _{}\sigma ^a\psi _++\beta _{}\sigma ^a\beta _+$$ (57) The central extension (level) cancels between the fermions and ghosts and the result is that $`L^a`$ generate an $`su(2)_k`$ current algebra at level $`k=0`$: $$L^a(z)L^b(0)\frac{1}{z}f^{abc}L^c(0)+\mathrm{reg}.$$ (58) Again, since $`𝒪_A`$ takes the form of the Casimir of $`su(2)`$, $`𝒪_A`$ preserves a global $`su(2)`$ symmetry. Furthermore it is easy to check that this global $`su(2)`$ commutes with the $`osp(2\|2)`$ symmetry: $$[L^a,J^b]=0$$ where $`J^a`$ are the $`osp(2\|2)`$ currents defined in eq. (37). In summary, the model has a global $`osp(2\|2)su(2)`$ symmetry for general $`g_\alpha ,g_m,g_A`$. ### B Quasi-spin-charge separation in the conformal field theory Since $`𝒪_g`$ only involves the $`osp(2\|2)_2`$ currents and $`𝒪_A`$ the $`su(2)_0`$ currents, it is important to understand how the conformal field theory decomposes in terms of these current algebras. It is straightforward to check that the $`su(2)_0`$ and $`osp(2\|2)_2`$ current algebras commute, so that the conformal field theory defined by $`S_{cft}`$ in eq. (19) has an $`osp(2\|2)_2su(2)_0`$ symmetry. We now show the much stronger result that the full stress tensor, which is the Sugawara stress tensor for the $`osp(4\|4)_1`$ supercurrent algebra, separates into two commuting pieces: $$T_{osp(4\|4)_1}=T_{osp(2\|2)_2}+T_{su(2)_0}$$ (59) All $`T`$’s in the above equation are Sugawara stress tensors and have Virasoro central charge equal to zero. The above equation is proven in the appendix. Since $`osp(2\|2)_2`$ contains charged currents whereas $`su(2)_0`$ does not, the equation (59) implies a kind of spin-charge separation and we will use this terminology in the sequel. Based on the separation (59) one would expect that the Hilbert space of $`osp(4\|4)_1`$ factorizes. However by trying to perform explicitly this factorization in some simple cases, we found that it is not possible: $$_{osp(4\|4)_1}_{osp(2\|2)_2}_{su(2)_0}$$ This non-factorization is intimitely related to the way the fields that appear in $`\stackrel{~}{𝒪}`$ transform under $`osp(2\|2)`$. This is described in Figure 2. In particular, the $`su(2)_0`$ currents $`L^a`$ turn out to be susy exact, e.g.: $$S_+(z)V_{}^a(0)\frac{1}{z}L^a$$ (60) Let us argue for the non-factorazibility of the Hilbert space by contradiction. The $`24`$ fields of $`\stackrel{~}{𝒪}`$ form three copies of an eight dimensional reducible but indecomposable representation of $`osp(2\|2)`$, that we shall denote $`[\stackrel{~}{8}]`$. Consider the state $`\|L_A^a(L_{f,1}^aL_{b,1}^a)\|0`$, where $`L_1^a=𝑑zL^a(z)/2i\pi z`$. If the factorization of $`_{osp(4\|4)_1}`$ holds, then $$\|L_A^a=\underset{i}{}\|osp_i\|su_i[\stackrel{~}{8}][3]$$ where $`[3]`$ is the adjoint representation of $`su(2)`$. Let us now act on $`\|L_A^a`$ with the $`su(2)_0`$ current. As a state in $`_{osp(4\|4)_1}`$ one finds that $`L_1^b\|L_A^a=\delta ^{ab}\|0_{\mathrm{Fock}}`$ with $`\|0_{\mathrm{Fock}}`$ the vacuum of the Fock space. As a state in $`_{osp(2\|2)_2}_{su(2)_0}`$, one would find $$L_1^b\|L_A^a=\|osp_iL_1^b\|su_i[\stackrel{~}{8}][1]$$ where $`[1]`$ is the $`su(2)`$ singlet. Thus if the factorization holds, one would deduce that $$\|0_{\mathrm{Fock}}[\stackrel{~}{8}][1]$$ Since $`\|0_{\mathrm{Fock}}`$ is annihilated by all $`osp(2\|2)`$ generators, this would mean that $`\|0_{\mathrm{Fock}}`$ is in the image of susy generators, i.e. there would exist a state $`\|\mathrm{\Omega }`$ such that for example $`\|0_{\mathrm{Fock}}=S_{}\|\mathrm{\Omega }`$. This is clearly a contradiction. A similar argument shows that eq. (60) leads to a contradiction in the factorization of $`V_{}^a`$. Another argument for the non-factorization is based on the fact that the current algebras $`osp(2\|2)_2`$ and $`su(2)_0`$ lead to logarithmic correlation functions, whereas the $`osp(4\|4)_1`$ is a free theory with no logarithms. It appears impossible for these logarithms to cancel in the product of two logarithmic correlation functions. The logarithmic nature of these theories can in fact be traced to transformations such as (60). We now describe a few of the results we need concerning these current algebra conformal field theories. Representations of $`osp(2\|2)`$ are characterized by the quantum numbers of the $`su(2)u(1)`$ subalgebra generated by $`(J,J_\pm )`$ and $`H`$. Highest weights can be labeled $`(j,b)`$ where $`j=0,1/2,1,..`$ is the spin of the charge-$`su(2)`$ and $`b=H/2`$. These representations have dimension $`8j`$. The conformal scaling dimension (left or right-moving) of the corresponding primary fields, as determined from the Sugarawa construction, is $$\mathrm{\Delta }_{(j,b)}^{osp(2\|2)}=\frac{2(j^2b^2)}{2k}$$ (61) On the other hand, the primary fields of the $`su(2)_k`$ current algebra are characterized by the spin $`j`$ of the spin-$`su(2)`$ only and have conformal dimension $$\mathrm{\Delta }_j^{su(2)}=\frac{j(j+1)}{k+2}$$ (62) A simple check of eq. (59) is based on the dimensions of the eight original fermion and ghost fields $`\psi _\pm ^i,\beta _\pm ^i`$. In the original conformal field theory these fields have conformal dimension $`\mathrm{\Delta }=1/2`$. Under the $`osp(2\|2)`$ they transform according to the two four-dimensional representations $`(\psi _+^1,\psi _{}^2,\beta _{}^2,\beta _+^1)`$ and $`(\psi _{}^1,\psi _+^2,\beta _+^2,\beta _{}^1)`$ corresponding to $`(j=1/2,b=0)`$ and its conjugate. At level $`k=2`$, these have dimension $`\mathrm{\Delta }^{osp(2\|2)}=1/8`$. The same fields transform as spin $`j=1/2`$ doublets according to the $`su(2)_0`$, and have dimension $`\mathrm{\Delta }^{su(2)}=3/8`$ at level zero. The exact decomposition (59) implies that these dimensions must add up properly: $`1/2=1/8+3/8`$. ### C Infrared fixed points and critical exponents We now study the model when $`g_\alpha +g_m=0`$. Setting $`g_\alpha =g_m=g`$, the effective action contains the current-current operators $`𝒪_g`$ and $`𝒪_A`$. The operator $`\stackrel{~}{𝒪}`$ which couples the $`osp(2\|2)_2`$ and $`su(2)_0`$ current algebras is not present. The model in was mapped onto a super-spin chain with $`osp(2\|2)`$ symmetry and Heisenberg-type hamiltonian. This further supports the identification of this super spin chain with our model on the line $`g_\alpha +g_m=0`$, since, at least for ordinary bosonic algebras, Heisenberg hamiltonians correspond to symmetry preserving current-current perturbations built on the quadratic casimir in the continuum limit. (As Lie superalgebras, $`osp(2\|2)`$ and $`sl(2\|1)`$ are identical.) The effective action now contains interactions that do not couple the spin and charge currents: $$S_{eff}=S_{\mathrm{cft}}+\frac{d^2x}{2\pi }\left(g𝒪_g+g_A𝒪_A\right)$$ (63) This decoupling is of course consistent with the more general $`\beta `$eta functions (34). Setting $`g_\alpha =g_m=g`$ one finds $$\beta _g=4g^2,\beta _{g_A}=2g_A^2$$ (64) As explained in section II, the model of Senthil et. al. corresponds to $`g`$ and $`g_A`$ positive. From the above $`\beta `$eta functions one sees that $`g_A`$ is marginally relevant whereas $`g`$ is marginally irrelevant. In order to deduce something exact concerning the infrared (IR) fixed point theory from the one-loop $`\beta `$eta functions, one needs to make a hypothesis concerning the role of the higher loop corrections. Consider for comparison the $`su(N)`$ Gross-Neveu models, which are $`su(N)`$ current-current perturbations, with $`\beta `$eta functions as in eq. (64), i.e. $`\beta _g=g^2`$ in a certain convention. When $`g>0`$ the perturbation is marginally relevant, i.e. $`g`$ grows at large distances. Higher loop corrections do not modify this, i.e. the flow to the IR does not stop at some finite value of $`g`$ corresponding to a non-trivial fixed point. Rather, $`g`$ eventually flows to infinity. The theory is thus massive, and in the infrared all these massive modes disappear leaving no massless degrees of freedom, i.e. the infrared theory is an empty theory. When $`g<0`$, the perturbation is marginally irrelevant, i.e. $`g`$ flows back to zero and the unperturbed conformal current algebra is recovered in the IR. We will make the hypothesis that this is the only possible behavior for general current-current perturbations, i.e. the only possible fixed points are $`g=0`$ or $`\mathrm{}`$. For our model we then have the following picture. In the IR, $`g_A`$ flows to infinity and the spin degrees of freedom are massive and decouple. The coupling $`g`$ on the other hand is marginally irrelevant. Therefore the IR fixed point is the coset $`osp(4\|4)_1/su(2)_0`$, and the theory arrives in the IR via the operator $`𝒪_g`$. Due to the non-factorizability of the Hilbert space, this coset conformal field theory is not precisely the $`osp(2\|2)_2`$ current algebra, even though it possesses this current algebra as a symmetry and the conformal dimensions of the coset theory are the same as the $`osp(2\|2)`$ current algebra because of (59). This identification of the IR fixed point allows the computation of certain critical exponents. The density of states $`\rho (E)`$ is $$\rho (E)=\frac{1}{V}\mathrm{Tr}\delta (HE)=\frac{1}{\pi V}\underset{\epsilon 0^+}{lim}\mathrm{Im}Tr\frac{1}{HEi\epsilon }$$ (65) where $`V`$ is the two-dimensional volume. This implies that the disorder averaged density of states is proportional to the one-point correlation function: $$\overline{\rho (E)}\mathrm{Re}\mathrm{\Phi }_E$$ (66) Let $`\mathrm{\Gamma }_E`$ equal the scaling dimension of $`\mathrm{\Phi }_E`$. Then, since the action (11) is dimensionless, viewing $`E`$ as a coupling, $`\mathrm{dim}(E)=2\mathrm{\Gamma }_E`$. Since $`E`$ is the only dimensionful coupling in the theory one deduces $$\overline{\rho (E)}E^{\mathrm{\Gamma }_E/(2\mathrm{\Gamma }_E)}\mathrm{as}E0$$ (67) One can also define a correlation length $`\xi _E`$, $$\xi _EE^{\nu _E},\nu _E=1/(2\mathrm{\Gamma }_E)$$ (68) (In , $`\nu _E`$ was referred to as $`\nu _B`$.) For our theory, $`\mathrm{\Phi }_E=\psi _+\overline{\psi }_{}+\mathrm{}`$ and $`\mathrm{\Gamma }_E`$ is the scaling dimension in the IR, which follows from the $`osp(2\|2)_2`$ conformal dimension of $`\psi _\pm `$, which as explained above is $`\mathrm{\Delta }_{(1/2,0)}^{osp(2\|2)}=1/8`$. Thus $`\mathrm{\Gamma }_E=1/4`$, and: $$\overline{\rho (E)}E^{1/7},\nu _E=4/7$$ (69) Numerical simulations of the $`osp(2\|2)`$ invariant spin-chain agree very well with $`\nu _E=4/7`$ . Since the system flows toward the infrared fixed point along a marginal direction, the scaling (69) is up to computable logarithmic corrections. In the above scenario, the $`su(2)_0`$ spin degrees of freedom are becoming localized, hence the terminology “spin quantum hall effect”. The above model can also exhibit charge localization by changing the signs of the couplings. Namely, if $`g`$ and $`g_A`$ are negative, then $`g`$ is marginally relevant, i.e. it flows to $`\mathrm{}`$ in the IR and $`g_A`$ is marginally irrelevant. Following the same reasoning as above, in this case the charge degrees of freedom decouple in the flow. The IR fixed point is now the coset $`osp(4\|4)_1/osp(2\|2)_2`$. Here, $`\mathrm{\Gamma }_E=3/4`$, and: $$\overline{\rho (E)}E^{3/5},\xi _EE^{4/5}\mathrm{for}\mathrm{charge}\mathrm{delocalization}$$ (70) Finally, if $`g<0,g_A>0`$, then $`g,g_A`$ flow to $`\mathrm{}`$ and $`\mathrm{}`$ respectively, and both sectors are massive and decoupled in the IR. This implies $`\mathrm{\Gamma }_E=0`$, and a constant density of states at $`E=0`$. ### D Path Integral Factorization Let us now present a path integral derivation of the spin-charge separation which has the advantage of being non perturbative and thus valid all along the RG trajectory. One can view this formulation as a way of defining the coset $`osp(4\|4)_1/su(2)_0`$. It consists in decoupling the random gauge field $`A`$ using chiral gauge transformations. The spin-charge separation will then be a consequence of the fact that the $`osp(2\|2)`$ currents are invariant under the chiral gauge transformations. This decoupling is similar to the solution of the random gauge potential described in . Let $`S`$ be the fermionic action (11) and $`Z(A,\alpha ,m)`$ its partition function $`Z(A,\alpha ,m)=D\mathrm{\Psi }e^{S(\mathrm{\Psi })}`$. At a fixed realization of disorder, the gauge potential $`A`$ may be gauged away by a chiral gauge transformation by parameterizing $`A`$ as: $`i\overline{A}=G^1_{\overline{z}}G;iA=G^{}_zG^1`$ (71) with $`G`$ an element of the complex $`su(2)^C`$ group, i.e. $`G`$ is a two by two complex matrix with determinant one. This is always possible on the sphere. This parametization is such that $`\psi _{}(_{\overline{z}}+i\overline{A})\psi _+=(\psi _{}G^1)_{\overline{z}}(G\psi _+)`$. Let us now denote by $`\psi ^{}`$ the chiral gauge-transformed fermions: $`\psi _{}^{}=\psi _{}G^1`$ ; $`\overline{\psi }_{}^{}=\overline{\psi }_{}G^{}`$ (72) $`\psi _+^{}=G\psi _+`$ ; $`\overline{\psi }_+^{}=G^1\overline{\psi }_+`$ (73) Gauge transformed bosons $`\beta ^{}`$ are defined similarly. The fermionic action (11) may be rewritten in terms of $`G`$ and $`\mathrm{\Psi }^{}`$. It becomes: $`\mathrm{\Gamma }(\mathrm{\Psi }^{}\|G,\alpha ,m)`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d^2x}{2\pi }}(\psi _{}^{}_{\overline{z}}\psi _+^{}+\overline{\psi }_{}^{}_z\overline{\psi }_+^{}`$ (74) $`+i\psi _{}^{}G(\stackrel{}{\alpha }\stackrel{}{\sigma }m)G^{}\overline{\psi }_+^{}+i\overline{\psi }_{}^{}G^1(\stackrel{}{\alpha }\stackrel{}{\sigma }+m)G^1\psi _+^{})`$ (75) This is the effective action for the Dirac operator coupled to the disorder variables $`G\stackrel{}{\alpha }\stackrel{}{\sigma }G^{}`$, $`G^1\stackrel{}{\alpha }\stackrel{}{\sigma }G^1`$ and $`GmG^{}`$. As we shall see below, the field $`G`$ will be coupled to the spin degrees of freedom whereas the fermions $`\mathrm{\Psi }^{}`$ shall be coupled to the charged sector. The density of states has a factorized expression in terms of these new spin and charged variables: $`\rho (E)\mathrm{Re}\overline{\mathrm{\Psi }}^{}(GG^{})\mathrm{\Psi }^{}`$ (76) We must not neglect to take into account the jacobians of the transformations $`AG`$ and $`\mathrm{\Psi }\mathrm{\Psi }^{}`$ as well as to represent the partition function $`Z(A,\alpha ,m)`$. The jacobians, which are determinants of Dirac operators, are computed using the chiral anomaly. They may be expressed with the help of the WZW action: $`\left\|{\displaystyle \frac{D\mathrm{\Psi }}{D\mathrm{\Psi }^{}}}\right\|`$ $`=`$ $`\mathrm{Det}\left(i\text{/}+A\right)=\mathrm{exp}\left(S_{wzw}(GG^{})\right)`$ (77) $`\left\|{\displaystyle \frac{DA}{DG}}\right\|`$ $`=`$ $`\mathrm{exp}\left(4S_{wzw}(GG^{})\right)`$ (78) with $`S_{wzw}`$ the WZW action. To represent the partition function, note that $`\mathrm{Det}(i\text{/}+A)=Z(A,\alpha =0,m=0)`$. Therefore, chiral gauge transformations applied to the bosonic beta system give: $`{\displaystyle \frac{\mathrm{Det}\left(i\text{/}+A\right)}{Z(A,\alpha ,m)}}={\displaystyle D\beta ^{}\mathrm{exp}\left(\mathrm{\Gamma }(\beta ^{}\|G,\alpha ,m)\right)}`$ (79) Gathering the jacobians we get the action $`S=4S_{wzw}(GG^{})+\mathrm{\Gamma }(\mathrm{\Psi }^{}\|G,\alpha ,m)+\mathrm{\Gamma }(\beta ^{}\|G,\alpha ,m)`$ (80) with $`\mathrm{\Gamma }`$ defined in eq.(75). This is the action to compute correlations using $`G`$, $`\mathrm{\Psi }^{}`$ and $`\beta ^{}`$ as the path integral variables at fixed disorder. One still has to add the disorder measure (3,6) to compute averaged correlations. In eq.(80) the random spin variables $`A`$ or $`G`$ are not yet decoupled. This decoupling only appears when $`g_\alpha +g_m=0`$. Indeed, in that case integrating over the disorder $`\alpha `$ and $`m`$ yields to current-current type interactions as in eq.(50) but with $`osp(2\|2)`$ currents bilinear in the fermions $`\mathrm{\Psi }`$ instead of the fermions $`\mathrm{\Psi }^{}`$. The crucial point is now to remark that these $`osp(2\|2)`$ currents are invariant under the chiral gauge transformations (72). For example, $`\psi _{}^{}\psi _+^{}`$ $`=`$ $`\psi _{}GG^1\psi _+=\psi _{}\psi _+=J`$ (81) $`ϵ_{ij}\psi _{+i}^{}\psi _{+j}^{}`$ $`=`$ $`\psi _{+n}\psi _{+m}G_{in}G_{jm}ϵ_{ij}=\psi _{+n}\psi _{+m}ϵ_{nm}=J^{}`$ (82) where we used that $`G_{in}G_{jm}ϵ_{ij}=\mathrm{detG}ϵ_{nm}=ϵ_{nm}`$ since $`G`$ has unit determinant. Similarly, it is easily checked that all $`osp(2\|2)`$ currents are invariant under chiral gauge transformations. Note that this is true because the spin disorder variables belong to $`su(2)`$. This would not be valid if for example the spin disorder variables were taking values in $`u(2)`$ instead. Hence, after integrating over the disorder at $`g_\alpha +g_m=0`$, the spin random variables $`G`$ decouple from the $`\mathrm{\Psi }^{}`$ and $`\beta ^{}`$ system. This is the spin-charge separation. We are thus left with the action: $`S=4S_{wzw}(GG^{})+\mathrm{\Gamma }(\mathrm{\Psi }^{}\|G=1,\alpha ,m)+\mathrm{\Gamma }(\beta ^{}\|G=1,\alpha ,m)`$ (83) with the same disorder measure for $`G`$ and $`\alpha ,m`$ as in eq.(3,6) with $`g_A`$ arbitrary but $`g_\alpha +g_m=0`$. The first term describes a WZW theory on the coset space $`su(2)^C/su(2)`$ at level $`k=4`$. This may be thought of as the theory ‘inverse’ to the $`su(2)`$ WZW theory at level $`k=0`$ , since the conformal dimensions of primary field in the two theories have opposite sign. The second and third terms describe the $`osp(2\|2)`$ current-current perturbation of the $`\mathrm{\Psi }^{}\beta ^{}`$ system. As in previous section this $`\mathrm{\Psi }^{}\beta ^{}`$ system may also be described as an $`osp(4\|4)`$ WZW model at level one and the current-current perturbation only couples to the sub-sector generated by the $`osp(2\|2)`$ currents. In the last section we argued that $`g_A`$ flows to infinity under RG flow. At $`g_A=\mathrm{}`$, the measure (3) on $`G`$ is flat and only the WZW action $`S_{wzw}(GG^{})`$ at level $`k=4`$ remains. This describes the spin sector. If $`g_\alpha =g_m>0`$, the $`osp(2\|2)`$ current perturbation is irrelevant and we recover the previous description. Recall the original fermion $`\mathrm{\Psi }`$, from which the density of states is computed, are related to $`G`$ and $`\mathrm{\Psi }^{}`$ by the chiral gauge transformation (72), so that the density of states is factorized as in eq.(76). If $`g_\alpha =g_m<0`$, the current-current perturbation is relevant. We may then propose that in the infrared the $`\mathrm{\Psi }^{}\beta ^{}`$ system is described by the coset theory $`osp(4\|4)_{k=1}/osp(2\|2)_{k=2}`$, which is equivalent to an $`su(2)`$ model at level zero. As a consequence, the fermions $`\mathrm{\Psi }^{}`$ shall flow in the infrared to fields with scaling dimension opposite to that of $`G`$, and the original fermions $`\mathrm{\Psi }`$ shall flow to fields with zero scaling dimension. This suggests that in this case the density of states is finite and regular at zero energy. ## V RG Phase diagram In this section, we return to the general model with $`g_\alpha +g_m0`$, and describe the global features of the phase diagram based on the one-loop RG equations (34). Our method consists in extracting the asymptotics of the RG trajectories by looking for directions in the coupling constant space which are preserved by the RG flow. Then, to analyze whether these asymptotic trajectories are attractive or not, stable or unstable, we project the RG flow onto the sphere and study the vector field thus obtained. This will allow us to point out the special role played by perturbations along a so-called strange direction. Let us first look for lines in the coupling constant space, with coordinates $`g=(g_A,g_\alpha ,g_m)`$, which are preserved by the RG flow and which pass through the origin. These correspond to trajectories which are straight lines and therefore for which the RG velocity field $`\dot{g}=\beta (g)`$ is co-linear to the vector $`g`$. The equations for these fixed line trajectories are thus: $`\beta (g)g=0`$ (84) The above equation is equivalent to finding solutions of the RG equations of the form $`g^i(t)=x^i\lambda (t)`$ where $`x^i`$ are constants independent of the RG time $`t=\mathrm{log}l`$ and $`d\lambda (t)/dt=\lambda ^2`$. Substituting this into (24), one finds $$x^i=C_{jk}^ix^jx^k$$ (85) where as before $`C_{jk}^i`$ are the operator product coefficients. With the linear relation $`g^i=x^i\lambda `$ among the couplings, the effective action contains a single running coupling constant: $$S_{\mathrm{eff}}=S_{cft}+\lambda \frac{d^2x}{2\pi }𝒪_x,𝒪_x=\underset{i}{}x^i𝒪_i=x^A𝒪_A+x^\alpha 𝒪_\alpha +x^m𝒪_m$$ (86) For $`x^i`$ a solution to (85), the operator $`𝒪_x`$ closes on itself under operator product expansion: $$𝒪_x(z,\overline{z})𝒪_x(0)\frac{1}{z\overline{z}}𝒪_x(0)$$ (87) To lowest order the $`\beta `$eta function for this specific perturbation is $`\beta _\lambda =\lambda ^2+\mathrm{}`$ There are six solutions to eqs.(84) (or eqs. (85)). They are: $`A`$ $``$ $`\left(g_A0,g_\alpha =0,g_m=0\right)`$ (88) $`B`$ $``$ $`\left(g_A=0,g_\alpha =0,g_m0\right)`$ (89) $`C`$ $``$ $`\left(g_A=0,g_\alpha +g_m=0\right)`$ (90) $`D`$ $``$ $`\left(g_\alpha +g_m=0,g_A+2g_\alpha =0\right)`$ (91) $`E`$ $``$ $`\left(6g_\alpha =3(1+\sqrt{2})g_A=2(1+2\sqrt{2})g_m\right)`$ (92) $`F`$ $``$ $`\left(6g_\alpha =3(\sqrt{2}1)g_A=2(2\sqrt{2}1)g_m\right)`$ (93) The two first lines correspond to the simple models with only random gauge potential or only random mass . The third and fourth ones correspond to the case we discuss at length in section IV. Contrary to the first four solutions, the last two do not seem to have an obvious simple algebraic interpretation The operator $`𝒪_x`$ is bilinear in the currents, $`𝒪_x=J^Ad_{AB}\overline{J}^B`$ with $`J^A`$ the $`osp(4\|4)`$ currents and $`d_{AB}`$ some bilinear form. The condition (87) is then: $`f_I^{AB}d_{AN}d_{BM}f_K^{NM}=d_{IK}`$. The algebraic interpretation of that equation is not clear to us.. These six solutions correspond to twelve possible stable directions because we have to choose an orientation on the line. Only three of them are in the domain of positive coupling constants $`g_A>0`$, $`g_\alpha >0`$ and $`g_m>0`$. The first two, which are the solution $`A`$ with $`g_A>0`$ and the solution $`B`$ with $`g_m>0`$ are in the border of that domain. The third one, which is the solution $`F`$ with $`g_\alpha >0`$, sits in the middle of the domain of positive $`g`$’s. We shall call this solution the strange direction. To analyze deeper the RG flow let us now project it onto the sphere. This is worth doing since the beta functions are homogeneous. Since there are three coupling constants, we may parameterize them with two angles, that we shall denote $`\theta _1`$ and $`\theta _2`$, and the radial coordinates $`\rho `$, $`\rho ^2=g_A^2+g_\alpha ^2+g_m^2`$. The RG equations (34) may then be written as two equations for the angular variables $`\dot{\theta }_j=\rho \beta _j(\theta _1,\theta _2)`$ (94) together with one equation for the radial variable: $`\dot{\rho }=\rho ^2\beta _\rho (\theta _1,\theta _2)`$ (95) The explicit expressions for the vector fields $`\beta _j`$ or $`\beta _\rho `$ are easy to find. In the angular equations (94) we may absorb the factor $`\rho `$ into a redefinition of the parametrization of the RG trajectories. This does not change the topology of the RG curves but only the speed at which the RG trajectories flow on these curves. So we shall analyze the vector field $`\dot{\theta }_j=\beta _j(\theta _1,\theta _2)`$ on the sphere. See figure 1. The vector field $`\beta _j`$ has twelve zeroes which correspond to the twelve fixed directions (88-93). We can compute the sign of $`\dot{\rho }`$ in these directions to know whether these straight RG trajectories are escaping to infinity ($`\dot{\rho }>0`$) or are flowing back to the origin ($`\dot{\rho }<0`$). To decipher whether the fixed directions are attractive or not we have to analyze whether the corresponding zero of the vector field $`\beta _j`$ on the sphere is attractive or not. For that we linearize the vector field $`\beta _j`$ on the sphere at its zeroes and compute its eigenvalues. Positive eigenvalues correspond to repulsive fixed directions, zero eigenvalues to locally marginal directions. For a fixed direction to be generically attractive, the two eigenvalues have to be non positive; otherwise to be attracted to the fixed direction requires fine tuning. The result is: $`A_+(g_A>0)`$ $``$ $`\dot{\rho }>0\mathrm{and}\mathrm{repulsive}`$ (96) $`A_{}(g_A<0)`$ $``$ $`\dot{\rho }<0\mathrm{and}\mathrm{repulsive}`$ (97) $`B_{}(g_m>0)`$ $``$ $`\dot{\rho }<0\mathrm{and}\mathrm{repulsive}`$ (98) $`B_+(g_m<0)`$ $``$ $`\dot{\rho }>0\mathrm{and}\mathrm{attractive}`$ (99) $`C_{}(g_\alpha >0)`$ $``$ $`\dot{\rho }<0\mathrm{and}\mathrm{repulsive}`$ (100) $`C_+(g_\alpha <0)`$ $``$ $`\dot{\rho }>0\mathrm{and}\mathrm{attractive}`$ (101) $`D_{}(g_\alpha >0)`$ $``$ $`\dot{\rho }<0\mathrm{and}\mathrm{attractive}`$ (102) $`D_+(g_\alpha <0)`$ $``$ $`\dot{\rho }>0\mathrm{and}\mathrm{repulsive}`$ (103) $`E_{}(g_\alpha >0)`$ $``$ $`\dot{\rho }<0\mathrm{and}\mathrm{repulsive}`$ (104) $`E_+(g_\alpha <0)`$ $``$ $`\dot{\rho }>0\mathrm{and}\mathrm{repulsive}`$ (105) $`F_+(g_\alpha >0)`$ $``$ $`\dot{\rho }>0\mathrm{and}\mathrm{attractive}`$ (106) $`F_{}(g_\alpha <0)`$ $``$ $`\dot{\rho }<0\mathrm{and}\mathrm{repulsive}`$ (107) The indices $`\pm `$ refer to trajectories flowing to infinity or back to the origin. There are only four fixed directions which are generically attractive: they stand on the directions $`B_+`$, $`C_+`$, $`D_{}`$ and on the strange direction $`F_+`$. They all are asymptotes to RG trajectories flowing to infinity except $`D_{}`$ which corresponds to asymptotic direction of trajectories looping back to the origin. Let us pause to reconsider the case $`g_\alpha +g_m=0`$ is this language. In that case the solutions of the RG equations are $`1/g_A1/g_A^0=2t`$ and $`1/g_\alpha 1/g_\alpha ^0=4t`$. If $`g_A^0>0`$ and $`g_\alpha ^0>0`$, then $`g_A`$ increases whereas $`g_\alpha `$ decreases. These trajectories are asymptotic to the direction $`A`$ since $`g_A`$ blows up at a finite value of $`t`$ at which $`g_\alpha `$ is finite. However the blow up times are not physical since for $`g_A`$ large enough the one-loop analysis is no more valid and one has to relies on the non-perturbative analysis done in previous section. If $`g_A^0>0`$ but $`g_\alpha ^0<0`$, then both couplings increase in magnitude. Which one blows up first depends on the initial data. The asymptotic trajectories are then either the direction $`A`$ or $`C`$ depending if $`g_A^0+2g_\alpha ^0`$ is positive or negative, so that the line $`g_A+2g_\alpha =0`$ is a separatrix in the phase diagram. On contrary, if $`g_A^0<0`$ but $`g_\alpha ^0>0`$, the couplings decrease and flow back towards the origin along the direction $`D`$. The strange direction $`F_+`$ is the asymptotic direction for all trajectories starting initially with positive coupling constants $`g_{\alpha ,m,A}`$, except those which are fine tuned to be in the direction $`A`$ or $`B`$. The domain $`g_j>0`$ is stable under one-loop RG: no trajectories can escape from it. The fact that there is one and only one asymptotic direction in the domain of positive coupling constants reflects the universality of that behavior: whatever the values of the initial coupling constants, the system flows along this direction in the infrared. The strange direction is hence the one which should be relevant to the description of the low energy behavior of the $`su(2)`$ random Dirac operators, e.g. the network model described in section II. Furthermore, since $`\dot{\rho }>0`$ for the strange direction, it is a strongly coupled system. In summary, based on the one-loop RG equations, it appears the generic network model is in a different universality class than the one fine-tuned to the line $`g_\alpha +g_m=0`$. However it certainly remains a possibility that the higher loop corrections can spoil the above analysis. ## VI Discussion In summary, we have shown that the network model with the identification of couplings $`g_\alpha +g_m=0`$ exhibits a spin-charge separation in the effective disorder averaged theory, and this allows a precise identification of the infrared fixed point as the coset $`osp(4\|4)_1/su(2)_0`$. This coset conformal field theory has some novel features in that it possesses an $`osp(2\|2)_2`$ supercurrent algebra symmetry, but due to the non-factorizability of the Hilbert space, it is not identical to the $`osp(2\|2)`$ theory defined by the Sugawara construction. The resulting critical exponents agree with the predictions based on percolation and the numerical simulations of the super spin chain studied in . Based on the one-loop renormalization group, we have argued that the network model without the constraint $`g_\alpha +g_m=0`$ is in a different universality class. How our analysis relates to percolation remains an interesting open question. There appear to be two possibilities. One is that the 1-loop strange direction we described in section V survives to higher loops and indeed corresponds to percolation; our coset theory $`osp(4\|4)_1/su(2)_0`$ is then the wrong fixed point. Another possibility is that the higher loop corrections actually restore the symmetry to the line $`g_\alpha +g_m=0`$, and the coset $`osp(4\|4)_1/su(2)_0`$ is a new description of percolation with $`osp(2\|2)`$ symmetry. The latter possibility can be investigated by comparing the conformal blocks of percolation with those computed in . It would be interesting to construct explicitly the continuum field theory corresponding to the super spin chain along the lines of . It is important to understand the role of the global $`osp(2\|2)`$ symmetry of the spin chain if this continuum limit indeed corresponds to percolation. ## VII Acknowledgments We thank A. Ludwig, H. Saleur and D. Serban for discussions, and A. La Belière for spirited encouragement. D.B. was supported in part by the CNRS, by the CEA and the European TMR contract ERBFMRXCT960012. ## VIII Appendix: spin-charge separation of the stress tensor. We give here some details on the proof of eq.(59) for the energy momentum tensor. Recall that the Sugawara construction of the energy momentum tensor in terms of the currents $`J^a`$ is: $`T(z)=\kappa \underset{wz}{lim}\left(J^a(w)C_{ab}J^b(z)k{\displaystyle \frac{C_{ab}\delta ^{ab}}{(zw)^2}}\right)`$ (108) with $`C_{ab}`$ the quadratic Casimir and the normalization constant $`\kappa `$ is chosen such that the currents have conformal dimension one. In practice $`T`$ is computed by extracting the regular term in product $`J^a(z)C_{ab}J^b(0)`$. For the $`osp(2\|2)`$ algebra at level $`2`$ this gives: $`T_{osp(2\|2)_2}={\displaystyle \frac{1}{8}}\left(J^2H^2{\displaystyle \frac{1}{2}}(J_{}J_++J_+J_{})+(S_+S_{}S_{}S_+)+(\widehat{S}_{}\widehat{S}_+\widehat{S}_+\widehat{S}_{})\right)`$ (109) The $`osp(2\|2)`$ currents are given in eq.(37) in terms of the $`\beta \psi `$ system. This Sugawara energy momentum tensor can thus be expressed in terms of these bosonic and fermionic fields. A simple computation yields: $`T_{osp(2\|2)_2}`$ $`=`$ $`{\displaystyle \frac{1}{8}}(_z\psi _{}\psi _++_z\psi _+\psi _{}+_z\beta _{}\beta _+_z\beta _+\beta _{}4(\psi _{}\beta _+)(\beta _{}\psi _+)`$ (111) $`+3(\psi _{}\psi _+)^2(\beta _{}\beta _+)^2+2(\psi _{}\psi _+)(\beta _{}\beta _+))`$ Similarly, with the normalization (25) for the Pauli matrices, the $`su(2)`$ Sugawara tensor is: $$T_{su(2)_0}=\frac{1}{8}L^aL^a$$ with $`L^a`$ defined in eq.(57). Again with the help of the identity (49), this may be written in terms of the $`\beta \psi `$ system as: $`T_{su(2)_0}`$ $`=`$ $`{\displaystyle \frac{1}{8}}(3(_z\psi _{}\psi _++_z\psi _+\psi _{}+_z\beta _{}\beta _+_z\beta _+\beta _{})+4(\psi _{}\beta _+)(\beta _{}\psi _+)`$ (113) $`3(\psi _{}\psi _+)^2+(\beta _{}\beta _+)^22(\psi _{}\psi _+)(\beta _{}\beta _+))`$ Adding both pieces we get: $`T_{osp(2\|2)_2}+T_{su(2)_0}={\displaystyle \frac{1}{2}}\left(_z\psi _{}\psi _++_z\psi _+\psi _{}+_z\beta _{}\beta _+_z\beta _+\beta _{}\right)`$ (114) This is the energy momentum tensor $`T`$ for the conformal field theory $`S_{cft}`$.
warning/0003/quant-ph0003054.html	ar5iv	text	# Information transfer and fidelity in quantum copiers ## I Introduction Quantum copying has attracted considerable interest in recent years, ever since the discovery of the no-cloning theorem, and the universal quantum copying machine which copies arbitrary unknown qubits with the best fidelity. To date, most treatments have used fidelity to characterize the quality of the copies produced. The fidelity between two quantum states characterized by density operators $`\widehat{\rho }_1`$ and $`\widehat{\rho }_2`$ is $$F(\widehat{\rho }{}_{1}{}^{},\widehat{\rho }{}_{2}{}^{})=\left\{\text{Tr}\left[\sqrt{(\widehat{\rho }{}_{1}{}^{})^{1/2}\widehat{\rho }{}_{2}{}^{}(\widehat{\rho }{}_{1}{}^{})_{}^{1/2}}\right]\right\}^2.$$ (1) A good summary of its properties is given in Ref.. In the case where one of the states is pure, the fidelity is simply the square of the overlap between the two states. Many authors have made use of two fidelity measures for quantum copiers: the global fidelity of the combined output (both copies) of the copier, with respect to a product state of (unentangled) perfect copies, and the local fidelity of one copy with respect to the original input state. Here, we will concentrate on a different indicator of copying success: mutual information content between the copies and the originals. One finds that which copier is optimal depends greatly on which indicator is used. In practice, this will mean that what sort of quantum copier is best depends on what one wants to do with the copies afterward. This article proceeds in the following fashion: After commenting on some drawbacks of fidelity, and why one might want to use different indicators, we outline exactly what we mean by information content between copies and originals in Sec. II. General features of the copiers that will be considered are mentioned in Sec. III. Copiers optimized for maximum copied information are given in Sec. IV (and derivations are given in Appendixes A and B) for three cases: (1) when the information is decoded from the copies one state at a time; (2) when efficient block-coding schemes are used to transmit as much information as is allowed by the Holevo bound; and (3) when the copies are an unentangled product state. In Sec. V the performance of these copiers is assessed according to information transfer and fidelity criteria, and compared to the performance of fidelity-optimized copiers known previously. ## II Mutual Information and Fidelity Measures ### A Fidelity, and some of its drawbacks Fidelity is used in many fields as an indicator of closeness between two states, and is often quite useful. It is probably also one of the easiest such indicators to calculate. However, it sometimes suffers from a number of drawbacks (examples of which are given below) when used as a measure of closeness over broad classes of systems, so there will be times when one wants to use a different indicator. While a fidelity of $`1`$ obviously implies identical states, and $`0`$ implies orthogonal states, what intermediate values mean is highly dependent on the particular states that are being compared, particularly if both states are impure. Thus a statement such as “The fidelity between the two states was $`x`$,” to be unambiguous, often needs considerable additional information on the states that were compared. To give an example: For standard optical coherent states of complex amplitude $`\alpha `$, given by $$\|\alpha =e^{(1/2)\|\alpha \|^2}\underset{n=0}{\overset{\mathrm{}}{}}\frac{\alpha ^n}{n!}\|n,$$ (2) the fidelity between two pure coherent states $`\|\alpha `$ and $`\|\alpha +1`$ is always constant: $$F(\|\alpha \alpha \|,\|\alpha +1\alpha +1\|)=\frac{1}{e}.$$ (3) Now if $`\alpha =0`$, the two states are the vacuum and a low-photon-number coherent state — states with qualitatively different properties. However, if $`\alpha `$ is large, then $`\|\alpha `$ and $`\|\alpha +1`$ are macroscopic, and experimentally indistinguishable, but the fidelity between them is still $`1/e`$. Another drawback of fidelity is that it is not directly related to other quantities commonly measured in experiments. While the fidelity is an expectation value of an observable (the observable being either one of the two states), it cannot usually be calculated from the results of experiments whose aim is to do something other than measure fidelity. It is not in general directly related to expectation values or measurement probabilities of other quantities, so it does not say much about the usefulness of a copy. In this sense, fidelity characterizes the closeness of the mathematical representation of physical states more than the closeness of the physical properties of those states. Of course, in many situations, these two types of closeness are equivalent, but not always. For the specific case of quantum copiers, global or local fidelities are not robust to unitary transformations made on the copies individually after all copying has been completed, and also can be very high even though the copies are uncorrelated with the originals. For example, suppose a message is encoded in a binary alphabet of orthogonal states $`\|0,\|1`$, and sent through a lossless communication channel that interchanges the states, i.e., they undergo the transformation $`\|0`$ $`\|1,`$ (4) $`\|1`$ $`\|0,`$ (5) then the fidelity of the transmitted with respect to the initial state is zero, but nothing of interest has been lost. It is sufficient for an observer receiving the message to relabel the states which they receive to recover the original message. Conversely, consider the situation where very nonorthogonal states $`\|a`$ and $`\|b`$ are used to encode a message. Using appropriate error-correction schemes, some information can be reliably transmitted with this encoding. However, now suppose that the message is intercepted by an eavesdropper, who simply sends the same state $`\sqrt{1/2}(\|a+\|b)`$ on to the intended receiver every time. The fidelity between sent and received states is still very high, but the received message carries no information from the sender. Global fidelity measures are often particularly removed from experimental results, since they compare the combined state of both copies with a perfect copy state that is generally unattainable due to the no-cloning theorem. However, in practice, one usually makes copies so they can subsequently be considered only individually. ### B Mutual information measures A different, natural, measure of copying efficiency that can be used is the amount of mutual (Shannon) information shared between the original states, and the copies. This mutual information does away with some of the drawbacks of fidelity, as discussed below. Consider two observers: one of them, the sender (labeled $`A`$), is sending states chosen from some ensemble, where the a priori probability of sending the $`i`$th variety of state is $`P_i^A`$. The other observer, the receiver (labeled $`B`$), makes measurements on one of the copies, obtaining the $`j`$th measurement result with probability $`P_{j\|i}^B`$, given that the $`i`$th state was sent into the copier. The amount of information (in bits per sent state) that the receiver has obtained from the sender is the Shannon mutual information, given by $$I(A:B)=\underset{i,j}{}P_i^AP_{j\|i}^B\mathrm{log}_2\frac{P_{j\|i}^B}{P_j^B},$$ (6) where $`P_j^B`$ is the overall probability of the receiver obtaining the $`j`$th measurement result, averaged over the input states. To use this measure to characterize a copying machine, rather than the specific message encoding or the ingenuity of the receiver in constructing a measuring apparatus, three points should be noted. First, even if a perfect copier is used, the amount of information that can be transmitted from originals to copy depends on the ensemble of states that is used to encode the message. Thus, the information about the original extractable from the copy $`I(A:B)`$ must be compared to the amount of information extractable from the original $`I(A:A)`$. Secondly, if observer $`B`$ makes a suboptimal (in terms of recovering the original message) set of measurements, then $`B`$’s stupidity will affect the mutual information. To eliminate the effect of $`B`$’s ingenuity (or lack of it), it has to be assumed that optimal measurements are made to recover the encoded message. Thirdly, a characterisation of the copier would usually involve examining its information-copying performance for a given set of input states. However, these may occur with various a priori probabilities $`P_i^A`$. We will take the case where these probabilities are chosen to encode the maximum amount of information in the signal states to be most representative of the behavior of information in the copier. Thus the mutual information quantities that will be used in later sections of this article are $`I_m(A:B)`$ and $`I_m(A:A)`$, given by $$I_m(A:B)=\underset{\left\{P_i^A\right\}}{\mathrm{max}}\left[\underset{\left\{_{}\right\}}{\mathrm{max}}I(A:B)\right],$$ (7) where $`\left\{P_i^A\right\}`$ denotes the set of a priori probabilities of $`A`$ using the $`i`$th state in the encoding of the message, and $`\left\{_{}\right\}`$ is the set of all positive valued operator measures. We will call $`I_m`$ the copied information. While this quantity can be more laborious to calculate, it has some advantages over fidelity. It is unchanged by relabeling or by local unitary transformations on the copies after they have left the copier, as well as always being zero if the copies are independent of the originals. Also, such mutual information is a physical quantity of interest in its own right, and is in fact what one is interested in in many fields (such as cryptography, for example). Even where this is not the case, mutual information between originals and copies can often be calculated from probability distributions of experimental measurements. Furthermore, it is clear what the statement “the mutual information transfer from $`A`$ to $`B`$ is $`x`$ ” means physically, with no further knowledge of the actual quantum states that were sent. It could be said that the information-copying capacity of a quantum cloner quantifies the practical usefulness, in many situations, of the copies produced by it. There is a qualitative difference between information-theoretic quantities such as copied information, and quantities such as fidelity. Fidelity, and similar quantities such as the Hilbert-Schmidt norm or the Bures distance, are quantifications of relations between two quantum states (or, more precisely, between their mathematical representations), while information-theoretic quantities deal with the relations between ensembles of states. This is the reason that they are robust to such postcopying effects as relabeling of the copy states. ### C Ultimate and one-state copied information Consider the situation discussed in the previous subsection. Observer $`A`$ encodes a message into a sequence of quantum states, chosen from a set of states $`\{\widehat{\rho }{}_{i}{}^{A}\}`$ labeled by the index $`i`$. Each of the sent states has an a priori probability $`P_i^A`$ of being the $`i`$th one in the set. When the copying machine acts on the signal state $`\widehat{\rho }_i^A`$ , it produces a copy state $`\widehat{\rho }_i^B`$, which is usually different from the original. It has been shown that the mutual information between $`A`$ and $`B`$ can be no more than $`I_H(A:B)`$, given by $`I_m(A:B)`$ $``$ $`I_H(A:B)`$ (9) $`=`$ $`S({\displaystyle \underset{i}{}}P_i^A\widehat{\rho }{}_{i}{}^{A}){\displaystyle \underset{i}{}}P_i^AS(\widehat{\rho }{}_{i}{}^{A}),`$ (10) where $`S(\widehat{\rho })`$ is the von Neumann quantum entropy of state $`\widehat{\rho }`$: $$S(\widehat{\rho })=\text{Tr}\left[\widehat{\rho }\mathrm{log}_2\widehat{\rho }\right],$$ (11) a result known as the Holevo theorem. In practice, the transmitted information will usually be significantly less than $`I_H(A:B)`$. However, it has been shown that if $`A`$ encodes the message using only certain sequences of states out of all the possible ones (although still respecting the a priori probabilities of individual states), and $`B`$ makes measurements on whole such sequences rather than on individual states, then as the length of these sequences increases, the information capacity per state can approach arbitrarily close to the Holevo bound $`I_H(A:B)`$. This is called a block-coding scheme, and such a communication setup is analogous to sending and distinguishing only whole “words” at a time in the message, rather than individual “letters.” In this analogy, letters correspond to individual quantum states, and words to sequences of them. Naturally, only special choices of the “words” to be used will approach the Holevo bound, Eq. (II C). With this in mind, there are two obvious candidates for a mutual information quantity with which to characterize copiers: the ultimate copied information given by $`I_H`$, and the one-state copied information $`I_1`$, which is the maximum information obtainable if measurements are made on only one state at a time. Both will be considered in what follows. ## III General Properties of the Copying Setups Considered In the interest of clarity and simplicity (and, one must admit, ease of analysis), only the most basic relevant copying setups have been investigated. This should make the principles involved easier to see, without introducing too much complexity. Thus, we will consider the case where observer $`A`$ encodes a message into a binary sequence of pure quantum states $`\widehat{\rho }{}_{i}{}^{A}=\|\psi _i^A\psi _i^A\|(i=1,2)`$ with equal a priori probabilities of being sent ($`P_i^A=\frac{1}{2}`$). The $`P_i^A`$ are chosen to be one-half for two reasons: (1) this is the simplest case; (2) this is the situation where the maximum amount of information is encoded in the input states. Since there are only two input states, the dimension of the relevant Hilbert space can be reduced to 2 by appropriate unitary transformations, because the states span at most a two-dimensional manifold in Hilbert space. Any such can be written (discarding an irrelevant phase factor) in an orthogonal basis $`\{\|+,\|\}`$ as $`\|\psi _1^A`$ $`=`$ $`\mathrm{cos}\theta \|++e^{i\mu }\mathrm{sin}\theta \|,`$ (13) $`\|\psi _2^A`$ $`=`$ $`\mathrm{sin}\theta \|++e^{i\mu }\mathrm{cos}\theta \|,`$ (14) where the parameter $`\theta `$ ranges from $`0`$ to $`\pi /4`$ (other values of $`\theta `$ are equivalent to a relabeling of the two states). In the rest of the article, $`\mu `$ will be taken to be zero for simplicity, although all results can easily be extended to the nonzero case. This, then, gives a one-parameter family of input states: $`\|\psi _1^A`$ $`=`$ $`\mathrm{cos}\theta \|++\mathrm{sin}\theta \|,`$ (16) $`\|\psi _2^A`$ $`=`$ $`\mathrm{sin}\theta \|++\mathrm{cos}\theta \|.`$ (17) These can be fully labeled by the fidelity between them, $$f=F(\widehat{\rho }{}_{1}{}^{A},\widehat{\rho }{}_{2}{}^{A})=\mathrm{sin}^2(2\theta ).$$ (18) In similar fashion, by taking the least complex case, the copiers considered will be unitary, create only two copies, and be symmetric. By symmetric we mean that the reduced quantum states of both copies by themselves are equal. The unitarity of the copying process implies a “black box” process: no external disturbance is required during the copying. Probabilistic copiers are not considered here. Physically, there are two subsystems $`o`$ and $`c`$ (which can be considered two dimensional for reasons outlined above) put into the unitary copying machine, and two come out. At the input, the subsystem $`o`$ contains the original state to be copied, while $`c`$ contains a “blank” state that is always the same, irrespective of what enters at $`o`$. Both subsystems contain the (usually imperfect) copies when they exit the copier, while an ancillary machine state subsystem ($`x`$) is also used in some of the copiers. At the input, all three subsystems are unentangled, while at the output, entanglement is usually present. Due to unitarity, the full entangled output states consisting of all three subsystems $`o`$, $`c`$, and $`x`$ are pure, but the states of individual subsystems are in general mixed. ## IV Three Information-Optimized Quantum Copiers In this section, we present transformations for several copiers optimized for information transfer to the copies, given a binary sequence of equiprobable input states. All these copiers are symmetric. The input states are in general nonorthogonal, and the degree of orthogonality is characterized by $`f`$, the square of the overlap between the two input states $`\widehat{\rho }_1^A`$ and $`\widehat{\rho }_2^A`$. These will be compared to known fidelity-optimized copiers in the next section. ### A Copiers that optimize the one-state copied information Rather than carry out a tedious optimisation, it stands to reason that if any unitary copier allows one to extract as much information about the originals from the copies as from the originals themselves, then it achieves the optimum. Is there such a copier? Perhaps surprisingly, one finds that the Wootters-Zurek (WZ) quantum copying machine (used in the original proof of the no-cloning theorem) allows one to extract as much information (using a one state at a time extraction) from either of the copies as from the original. One can imagine that the same information transfer could be achieved by making measurements on the originals, and sending the results classically, but that a simple unitary transformation with no coupling to the external environment can achieve the same is perhaps less obvious. What is more, the WZ copier does much better than any fidelity-optimized copiers, as will be seen later. Explicitly, the transformation of the input states (III) is given by $`\|\psi _1^A`$ $``$ $`\mathrm{sin}\theta \|+++\mathrm{cos}\theta \|,`$ (20) $`\|\psi _2^A`$ $``$ $`\mathrm{cos}\theta \|+++\mathrm{sin}\theta \|,`$ (21) where the basis vectors $`\|+`$, etc., indicate tensor products $`\|+_o\|_c`$ of the basis vectors for the $`o`$ and $`c`$ copy subsystems, respectively. The combined state of the copies is highly entangled, but the reduced density matrices of the copies (the full output density matrices traced over all subsystems except one copy) are in the classically mixed states $`\widehat{\rho }_1^B`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}^2\theta & 0\\ 0& \mathrm{sin}^2\theta \end{array}\right),`$ (25) $`\widehat{\rho }_2^B`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{sin}^2\theta & 0\\ 0& \mathrm{cos}^2\theta \end{array}\right).`$ (28) The one-state copied information, which is the same as can be extracted from the originals, is $$I_1^{\text{WZ}}=\frac{1}{2}\left[(1+q)\mathrm{log}_2(1+q)+(1q)\mathrm{log}_2(1q)\right],$$ (30) where $`q`$, which we will call the distinguishability parameter, is $$q=\sqrt{1f}.$$ (31) From the purely classical nature of $`\widehat{\rho }_i^B`$, it follows that the ultimate copied information $`I_H^{\text{WZ}}`$ is no bigger than $`I_1^{\text{WZ}}`$. In fact, applying more WZ copying machines to the copies made by the first one, in a cascade effect, creates larger numbers of copies, each of which still carries the same amount of (one-state) information as the original message. In this way, arbitrary numbers of optimal copies can be made — similarly to how one can make arbitrary numbers of copies of classical information. The local fidelity between a copy and the originals is $$F(\widehat{\rho }{}_{i}{}^{A},\widehat{\rho }{}_{i}{}^{B})=1\frac{f}{2}.$$ (32) There are other copiers related to the WZ copier which allow the same optimal one-state information transfer. One example is the family of copying transformations created by applying identical local unitary transformations on both copies after they come out of the WZ copier. The particular transformation presented above in Eq. (IV A) is the one that gives the best local fidelity out of this family of transformations. ### B Copiers without ancilla that optimize the ultimate copied information It is also of interest how well information can be transmitted when the possibility of complicated block-coding schemes is allowed, as discussed in Sec. II C. To make the calculations relatively tractable analytically, we have made two restrictions on the copiers that we considered for this task. First, only copiers that do not use an ancillary subsystem $`x`$, entangled with the copies, have been considered. It is probably possible to obtain somewhat better performance in ultimate information copying by using such helper subsystems, since discarding $`x`$ after copying is completed partially relaxes the conditions that the copy states $`o`$ and $`c`$ must satisfy to preserve unitarity (since one then has more parameters left to optimize over). It is not clear how much better one could do with such helper states, but we suspect not much better, since from Fig. 2, the copier considered here is only marginally better than several others obtained by optimising over different indicators such as fidelity and one-state copied information. Secondly, for similar reasons, we have assumed that since both possible input states $`\widehat{\rho }_i^A`$ are of equal purity $`\text{Tr}\left[(\widehat{\rho }{}_{i}{}^{A})^2\right]`$ (totally pure, in fact), then both reduced copy states $`\widehat{\rho }_i^B`$ will be of equal purity also: $$\text{Tr}\left[(\widehat{\rho }{}_{1}{}^{B})^2\right]=\text{Tr}\left[(\widehat{\rho }{}_{2}{}^{B})^2\right].$$ (33) This is also a property shared by all other copiers mentioned in this article. The usual assumptions of Sec. III, such as both copies being equal, apply also. So, an ancillaless copier, that produces two identical (usually imperfect) copies of any of two possible pure signal states, that makes copies of the same purity whichever of the two input states is sent, and that (given the above) maximizes the amount of information that can be transmitted to each of the copies by any block-coding scheme when the two input states are equiprobable, is given by the somewhat lengthy characterisation below. The details of how this was obtained have been left for Appendix A. There is a whole family of copying transformations, related by local unitary transformations on the copies after they have stopped interacting with each other, which give the same ultimate information copied $`I_H^u`$. Of these, we will specify that particular one in this family which gives the greatest local fidelity between the copies and originals. The transformation can be written in terms of the parameters $`r_m`$ and $`\varphi _m`$, which have to be determined numerically. In terms of the initial states (III), $`\|\psi _1^A`$ $``$ $`\sqrt{{\displaystyle \frac{1+r_m}{2}}}\|b_1+\sqrt{{\displaystyle \frac{1r_m}{2}}}\|b_2,`$ (35) $`\|\psi _2^A`$ $``$ $`\sqrt{{\displaystyle \frac{x}{2}}}\|b_1+\sqrt{{\displaystyle \frac{x}{2}}r_m\mathrm{cos}\varphi _m}\|b_2`$ (37) $`+\sqrt{{\displaystyle \frac{1x+r_m\mathrm{cos}\varphi _m}{2}}}\left(\|b_3+\|b_4\right),`$ where $$x=\frac{1}{2}\left(1+\mathrm{cos}^2\varphi _m+2r_m\mathrm{cos}\varphi _m+\sqrt{1r_m^2}\mathrm{sin}^2\varphi _m\right),$$ (38) and the four $`\|b_j`$ are orthogonal basis states, given in terms of the usual $`\|+`$ and $`\|`$ basis states used in Eqs. (III) and (IV A) by the matrix equation $$\left(\begin{array}{c}\|b_1\\ \|b_2\\ \|b_3\\ \|b_4\end{array}\right)=U\left(\begin{array}{c}\|++\\ \|+\\ \|+\\ \|\end{array}\right),$$ (39) where the unitary matrix $`U`$ is $$U=\frac{1}{2}\left(\begin{array}{cccc}1+\mathrm{sin}\varphi _m/2& 1\mathrm{sin}\varphi _m/2& \mathrm{cos}\varphi _m/2& \mathrm{cos}\varphi _m/2\\ 1\mathrm{sin}\varphi _m/2& 1+\mathrm{sin}\varphi _m/2& \mathrm{cos}\varphi _m/2& \mathrm{cos}\varphi _m/2\\ \mathrm{cos}\varphi _m/2& \mathrm{cos}\varphi _m/2& 1+\mathrm{sin}\varphi _m/2& \mathrm{sin}\varphi _m/21\\ \mathrm{cos}\varphi _m/2& \mathrm{cos}\varphi _m/2& \mathrm{sin}\varphi _m/21& 1+\mathrm{sin}\varphi _m/2\end{array}\right).$$ (40) As can be seen from the above, the basis states $`\|b_j`$ are entangled over the two copies. The parameter $`\varphi _m`$ is actually the angle between the Bloch vectors of the two possible reduced copy states $`\widehat{\rho }_i^B`$, which can be written $`\widehat{\rho }_1^B`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1+q& q_H\\ q_H& 1q\end{array}\right),`$ (44) $`\widehat{\rho }_2^B`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1q& q_H\\ q_H& 1+q\end{array}\right),`$ (47) where the parameters $`q`$ and $`q_H`$ are $`q`$ $`=`$ $`r_m\mathrm{sin}{\displaystyle \frac{\varphi _m}{2}},`$ (49) $`q_H`$ $`=`$ $`r_m\mathrm{cos}{\displaystyle \frac{\varphi _m}{2}},`$ (50) and appear in the expressions for $`I_1`$ and $`I_H`$. Now $`\mathrm{cos}\varphi _m`$ is dependent on $`r_m`$, and is given in terms of it as the second largest real root of the following quartic polynomial in $`\mathrm{cos}\varphi _m`$: $`0`$ $`=`$ $`\mathrm{cos}^4\varphi _m\left[r_m^2(2r_m^22\sqrt{1r_m^2})\right]+\mathrm{cos}^3\varphi _m\left[4r_m^2(1\sqrt{1r_m^2})\right]`$ (53) $`+\mathrm{cos}^2\varphi _m\left\{2[r_m^4+2r_m^2+4f(\sqrt{1r_m^2}1)]\right\}+\mathrm{cos}\varphi _m\left[4r_m^2(1+\sqrt{1r_m^2}4f)\right]`$ $`+\left[(4f1)^2(1r_m^2)^2+2(r_m^24f)\sqrt{1r_m^2}\right].`$ The ultimate copied information is given by $`I_H^u`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(1+r_m)\mathrm{log}_2(1+r_m)+(1r_m)\mathrm{log}_2(1r_m)\right]`$ (55) $`{\displaystyle \frac{1}{2}}\left[(1+q_H)\mathrm{log}_2(1+q_H)+(1q_H)\mathrm{log}_2(1q_H)\right],`$ which can be made a function of $`r_m`$ only, using Eq. (53). To obtain the optimum copier, we find numerically the value of $`r_m`$ that maximizes $`I_H^u`$ on $`r_m[\sqrt{1f},1]`$. The one-state copied information $`I_1^u`$ is given by the same expression in the distinguishability parameter $`q`$ as for the WZ copier \[Eq. (30)\], with $`q`$ now given by Eq. (49). It is interesting to note that, for input states which are sufficiently nonorthogonal ($`f0.206`$), the copier given here is just the WZ copier described in Sec. IV A. In these cases, $`\varphi _m=\pi `$ and $`r_m=\sqrt{1f}`$. This sudden change in behavior (particularly evident in Fig. 1 and Fig. 3 ) may be due to excluding the use of ancillary subsystems. Allowing these may make the $`I_H`$ optimal copier consistently better (although possibly not by much) than the Wootters-Zurek for all values of $`f`$, even the small ones. The local fidelity between copies and originals for this copier is $$F(\widehat{\rho }{}_{i}{}^{A},\widehat{\rho }{}_{i}{}^{B})=\frac{1}{2}(1+q\sqrt{1f}+q_H\sqrt{f}).$$ (56) ### C An optimal copier that gives unentangled copies As has been remarked by many previously, optimal quantum copiers typically produce highly entangled copies. This also applies to the two quantum copiers given in Secs. IV A and IV B. Nevertheless, copies of some quality can be made without entanglement between them. This might be desirable in some situations. Once again two simplifying assumptions have been made to make the calculation easier. It has been assumed that the copies are, again, unentangled with ancillary machine states, and that the output state of the copier is simply a product state of the two identical copies, rather than a classical mixture of several such product states. The case with additional machine states present might allow somewhat higher information transmission $`I_H`$ with block-coding methods, for the same reasons as in Sec. IV B. This would be interesting to check, but we have not done this to date. Allowing classical correlations between copies and a machine state subsystem $`x`$ does not, however, improve information transmission. Given the above two restrictions, a copier that optimizes both the one-state and ultimate copied information, while keeping the copies unentangled, is given by $`\|\psi _1^A`$ $``$ $`{\displaystyle \frac{1+\sqrt{1\sqrt{f}}}{2}}\|+++{\displaystyle \frac{1\sqrt{1\sqrt{f}}}{2}}\|`$ (59) $`+{\displaystyle \frac{1}{2}}f^{1/4}\left(\right\|++\|+),`$ $`\|\psi _2^A`$ $``$ $`{\displaystyle \frac{1\sqrt{1\sqrt{f}}}{2}}\|+++{\displaystyle \frac{1+\sqrt{1\sqrt{f}}}{2}}\|`$ (61) $`+{\displaystyle \frac{1}{2}}f^{1/4}\left(\right\|++\|+),`$ with notation identical to Eqs. (IV A). See Appendix B for details of the optimisation. This gives pure state copies (they must be pure from the unitarity of the transformation, since the input states are pure, and the output state is $`\widehat{\rho }{}_{i}{}^{B}\widehat{\rho }_i^B`$) $`\widehat{\rho }_1^B`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1+\sqrt{1\sqrt{f}}& f^{1/4}\\ f^{1/4}& 1\sqrt{1\sqrt{f}}\end{array}\right),`$ (65) $`\widehat{\rho }_2^B`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1\sqrt{1\sqrt{f}}& f^{1/4}\\ f^{1/4}& 1+\sqrt{1\sqrt{f}}\end{array}\right).`$ (68) A family of copiers which do as well in the information measures, but worse in local fidelity between originals and copies, is given by making unitary transformations on the copies individually. The one-state copied information $`I_1^{\text{NE}}`$ is given by the same expression in $`q`$ as for the WZ copier (30), with $`q`$ now given by $$q=\sqrt{1\sqrt{f}}.$$ (69) The ultimate copied information is $$I_H^{\text{NE}}=1\frac{1+f^{1/4}}{2}\mathrm{log}_2(1+f^{1/4})\frac{1f^{1/4}}{2}\mathrm{log}_2(1f^{1/4}).$$ (70) The local fidelity of copies with respect to originals is $$F(\widehat{\rho }{}_{i}{}^{A},\widehat{\rho }{}_{i}{}^{B})=\frac{1}{2}[f^{3/4}+1+\sqrt{(1f)(1\sqrt{f})}].$$ (71) It turns out that this copier also gives the best local fidelity out of such unentangling copiers (see Appendix B). ## V A Comparison of the Copiers To see how well the copiers rate in terms of the information measures $`I_H`$ and $`I_1`$, we first need to determine how much information could be extracted from the input states if they were not copied. Since the input states are not orthogonal for $`f>0`$, then a full bit of information cannot be extracted from each state even though they are equiprobable. One finds that the information extractable one state at a time is $$I_1^o=\frac{1}{2}\left[(1+q)\mathrm{log}_2(1+q)+(1q)\mathrm{log}_2(1q)\right],$$ (72) with the distinguishability parameter $`q=\sqrt{1f}`$. This is the same as with the Wootters-Zurek copier (30). The ultimate information extractable from the signal if block-coding methods are used is, however, unlike that for the WZ copier, much larger: $$I_H^o=1\frac{1+\sqrt{f}}{2}\mathrm{log}_2(1+\sqrt{f})\frac{1\sqrt{f}}{2}\mathrm{log}_2(1\sqrt{f}).$$ (73) It is interesting to compare the performance of the copiers given in Sec. IV to previously known fidelity-optimized ones. Three will be considered here, and a brief summary of the copies they produce is given in Appendix C in terms of the input state overlap parameter $`f`$. These three copiers are as follows. (1) The universal quantum copying machine (UQCM), which copies arbitrary qubits with a local fidelity of $`5/6`$. This is the maximum possible if it is to copy all with equal fidelity. (2) A copier found by Bruß et al. that optimizes the global fidelity when copying one of two nonorthogonal input states. (3) A copier also found by Bruß et al. that optimizes the local fidelity when copying one of two nonorthogonal input states. So let us see how they compare in performance. ### A One-state copied information The one-state copied information is a good indicator of the efficiency of communicating classical data to the two copies. The recovery and coding of the information in this case relies only on measurement of one-qubit states, and classical error-correction schemes. Looking at Fig. 1, one sees that the Wootters-Zurek copier, apart from achieving the optimum and transmitting as much one-state information to both copies as was encoded originally, is also far better at it than any of the other copiers shown (except for the small-$`f`$ region, where the ultimate-information optimized copier becomes the WZ). The WZ copier has by far the simplest transformation out of these copiers, so it seems that for basic information transmission the simplest copier is the best. The fidelity-optimized copiers do not do as well as the WZ, which in itself is to be expected, as after all they were optimized for fidelity, not information transfer. However, they do very much worse, causing the loss of much information that could be regained if better copiers were used. This shows quite clearly that fidelity is not necessarily a good measure of the quality of the copies for all situations. It is perhaps also surprising that even though we are considering information transmitted to one copy here, the copier that has been optimized for global fidelity between the combined output state and perfect copies, does significantly better than the one that has been optimized for local fidelity between a single copy and original. The UQCM gives much less information transfer than the other copiers, since all the others have been specifically tailored for the two signal states, whereas the UQCM must handle any arbitrary states with equal fidelity. The copiers that give optimum unentangled copies do generally significantly worse than the other copiers apart from the UQCM, but one sees that all the copiers apart from the WZ copier and UQCM converge to the same efficiency (much worse than the optimum) for high values of $`f`$, i.e., when the signal states are not very orthogonal. Note that a plot of the actual (rather than relative) amount of information extractable from the original signal $`I_1`$ is shown in Fig. 2 as the Wootters-Zurek curve, since $`I_H^{\text{WZ}}=I_1^o`$. ### B Ultimate copied information The ultimate (Holevo bound) copied information $`I_H`$ gives an absolute maximum on how much information could possibly be transmitted by a given copier, with the best signaling scheme that is possible. In general, to achieve this bound, the encoding/decoding scheme has to be very elaborate, and it is often not achievable in practice due to complexity. In the case of qubit systems being transmitted here, this would entail making measurements of many-qubit observables to decode the information: a difficult task at present. As can be seen in Fig. 2, most of the copiers cluster just below the optimal capacity achieved by the copier of Sec. IV B. While this is not necessarily the absolute optimum that can be achieved, as there remains the possibility that introducing helper machine states may increase this bound, this bunching makes it seem plausible that no large gains can be achieved beyond this. This ultimate-information optimal copier is quantitatively not much better than the Wootters-Zurek copier. Its greatest gains, which are still quite modest, come when the overlap between signal states is high, where the absolute information content in the signal is small. It can be seen that, while the no-cloning theorem did not stop one from perfectly copying information contained in one state at a time, its effect is strong where block-coding schemes are allowed. This is because, if we restrict ourselves to the one case at a time situation, we are not utilising those properties of the states that are affected by the no-cloning theorem. The difference between what can be extracted from a copy and from the originals is quite striking, and for highly overlapping input states, over 60% of the information in the originals is unavailable from a copy. The behavior of the copiers for high overlap between states is as one would expect. That is, the Wootters-Zurek copier becomes much less efficient than the others when block-coding schemes are used, as the other copiers do not fully entangle the copies with each other, thus allowing one to extract some extra information by looking at several sequential states together. Since the Wootters-Zurek copier has $`I_1=I_H`$, by comparing the values of $`I_H`$ for the local and global-fidelity-optimized copiers to the WZ copier, one can see that for these fidelity-optimal copiers, much more information than $`I_1`$ can be sent to the copies by allowing complicated block-coding schemes which use correlations between subsequent signal states. This approach, however, is unhelpful with the Wootters-Zurek copier, and is of very little help when using the the UQCM. As for the other information measure, the global-fidelity-optimized copier does slightly better than the local fidelity one. The unentangled copier does slightly worse than the rest, except for the UQCM which is consistently worse on all counts, as it is not tailored to the input states like the others. ### C Local fidelity This is shown for various copiers in Fig. 3. The UQCM is absent from the plot, as its local fidelity lies far below the others shown there. Figures 1 and 3 show quite clearly that fidelity and information transfer quantify quite different properties of the copying transformation, and one has to keep in mind which properties are desired, before deciding on a quantity to characterize efficiency. As expected, the best local fidelity occurs for the copier that was optimized for this, and the global fidelity optimal copier is almost as good. The WZ copier is no good at fidelity at all for significantly overlapping states. The unentangled copier is once again slightly worse than most of the others. The sharp change in behavior for the ultimate-information optimal copier is particularly evident in this plot. ## VI Comments and Conclusions It was seen in the previous section that quantum copiers optimized for fidelity measures are far from optimal for basic information transmission to the copies, and, vice versa, information-optimized copiers are far from optimized for fidelity between copies and originals. This indicates that various measures of quality should be used for quantum devices, depending on what final use is to be made of the states created. Some other general trends that were seen for the quantum copying devices that were considered, include the following. The ultimate-copied-information-optimized copier behaves more similarly to the fidelity-optimized ones than to the one-state optimized WZ copier (where it differs from the WZ). The fidelity-optimized copiers are not bad when one allows multiparticle measurements on the copies, but are far from optimal if one does not. This may be because the fidelity-optimized copiers preserve some of the quantum superposition of the input states (as evidenced by the off-diagonal terms in the density matrices of the copies), whereas the WZ copier makes the copies purely classical mixtures when they are considered individually. To get extra information transmission by making measurements on multistate observables, one needs some quantum effects between the successive copy states, and these effects are lacking with the WZ copier. A small, but perhaps surprising feature was that the global-fidelity-optimized copier gave better performance in the information measures than the local-fidelity-optimized one, even though only information flowing to one copy was considered. Other features seen include the poor performance of the UQCM relative to the other copiers — unsurprising, since the other ones are tailored specifically to the two signal states, and the poorer performance when the copies are made unentangled with each other. For all copiers considered, when the input signal states are nonorthogonal, the information carrying capacity of a channel between two observers is significantly greater when the receiver gets undisturbed states ($`I_H^o`$) than when the receiver gets one copy, even when the copier is highly optimized ($`I_H^u`$). This is an information-theoretic manifestation of the no-cloning theorem. ## A Derivation of Ultimate-Information Optimal Copier The copier sought has the following properties: it takes one of two ($`i=1,2`$) pure input states $`\widehat{\rho }_i^A`$ of Hilbert space dimension 2, and by a unitary transformation creates a state $`\widehat{\rho }_i`$ consisting of two (possibly entangled) copies ($`\widehat{\rho }_i^o`$ and $`\widehat{\rho }_i^c`$), again of Hilbert space dimension 2. The state of each copy, when the other copy is ignored, is identical, and both possible copy states (corresponding to input states) have equal purity, as measured by their self-fidelity $`\text{Tr}[\widehat{\rho }{}_{}{}^{2}]`$. Assuming all states considered are normalized, these conditions can be written as normalization: $`\text{Tr}[\widehat{\rho }{}_{i}{}^{A}]=1,`$ (A1) input pure: $`\widehat{\rho }{}_{i}{}^{A}=\|\psi _i^A\psi _i^A\|,`$ (A2) unitarity: $`\widehat{\rho }{}_{i}{}^{}=\|\psi _i\psi _i\|,`$ (A3) $`\text{Tr}[\widehat{\rho }{}_{1}{}^{}\widehat{\rho }{}_{2}{}^{}]=\text{Tr}[\widehat{\rho }{}_{1}{}^{A}\widehat{\rho }{}_{2}{}^{A}]=f,`$ (A4) symmetry: $`\widehat{\rho }{}_{i}{}^{o}=\text{Tr}_c[\widehat{\rho }{}_{i}{}^{}]=\widehat{\rho }{}_{i}{}^{c}=\text{Tr}_o[\widehat{\rho }{}_{i}{}^{}]=\widehat{\rho }{}_{i}{}^{B},`$ (A5) equal purity: $`\text{Tr}\left[(\widehat{\rho }{}_{1}{}^{B})^2\right]=\text{Tr}\left[(\widehat{\rho }{}_{2}{}^{B})^2\right].`$ (A6) And, of course, on top of these conditions, the Holevo bound on ultimate information copied $`I_H`$ is to be maximized. The output states can be written in terms of a vector of complex expansion coefficients in some basis as $$\|\psi _j=\frac{1}{\sqrt{2}}[\alpha _j,\beta _je^{i\varphi _{\beta j}},\gamma _je^{i\varphi _{\gamma j}},\delta _je^{i\varphi _{\delta j}}],$$ (A7) where $`\alpha _j`$,$`\beta _j`$,$`\gamma _j`$,$`\delta _j`$ $`[0,\sqrt{2}]`$, and the angles $`\varphi _{\mathrm{}}[0,2\pi )`$. Normalization gives $`\alpha _j^2+\beta _j^2+\gamma _j^2+\delta _j^2=2`$. One of the expansion coefficients can be made real and positive, without affecting the final bound, by multiplying by appropriate unphysical phase factors, so let us do this to the $`\alpha _j`$. Now, any two states in a two-dimensional Hilbert space (such as the reduced states of the two possible copies $`\widehat{\rho }_1^B`$ and $`\widehat{\rho }_2^B`$), can be described by two Bloch vectors $`𝐫_i`$. The states are then given by $$\widehat{\rho }{}_{i}{}^{}(𝐫_1)=\frac{1}{2}(\widehat{I}+𝝈𝐫_i)\text{where}𝝈=[\widehat{\sigma }{}_{1}{}^{},\widehat{\sigma }{}_{2}{}^{},\widehat{\sigma }{}_{3}{}^{}]$$ (A8) and $`\widehat{\sigma }_j`$ are the Pauli matrices. By an appropriate choice of basis, one of the two Bloch vectors can be chosen to lie in an arbitrary direction, while the other is separated by some angle $`\varphi _r`$ from the first, both of them lying in a plane of our choosing. Thus there are only three parameters for these two states that are not arbitrary, depending on the choice of basis: the lengths of the Bloch vectors $`r_i`$, and the angle between them $`\varphi _r`$. Also, since $$\text{Tr}\left[\widehat{\rho }{}_{i}{}^{}(𝐫_i)_{}^{2}\right]=\frac{1}{2}(1+\|𝐫_i\|^2),$$ (A9) and we are assuming equal copy purity (A6), both Bloch vectors are of equal length $`r=\|𝐫_i\|`$. Let us choose these Bloch vectors to be $$𝐫_1=r[0,0,1]\text{and}𝐫_2=r[\mathrm{sin}\varphi _r,0,\mathrm{cos}\varphi _r].$$ (A10) Thus, without any loss of generality, the copies can be written in an appropriate basis as $`\widehat{\rho }_1^B`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1+r& 0\\ 0& 1r\end{array}\right),`$ (A14) $`\widehat{\rho }_2^B`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1+r\mathrm{cos}\varphi _r& r\mathrm{sin}\varphi _r\\ r\mathrm{sin}\varphi _r& 1r\mathrm{cos}\varphi _r\end{array}\right).`$ (A17) Using Eqs. (A10), (A7), and conditions (A1), (A5), one obtains the following restrictions on the expansion coefficients of the total output states $`\widehat{\rho }_i`$: $$\begin{array}{ccccccc}\hfill \gamma _1& =& \beta _1,\hfill & & \hfill \gamma _2& =& \beta _2,\hfill \\ \hfill \beta _1^2& =& 1+r\alpha _1^2,\hfill & & \hfill \beta _2^2& =& 1+r\mathrm{cos}\varphi _r\alpha _2^2,\hfill \\ \hfill \delta _1^2& =& \alpha _1^22r,\hfill & & \hfill \delta _2^2& =& \alpha _2^22r\mathrm{cos}\varphi _r,\hfill \end{array}$$ (A19) $`\beta _1\left(\alpha _1e^{i\varphi _{\beta 1}}+\delta _1e^{\varphi _{\delta 1}\varphi _{\beta 1}}\right)`$ $`=`$ $`0,`$ (A20) $`\beta _2\left(\alpha _2e^{i\varphi _{\beta 2}}+\delta _2e^{\varphi _{\delta 2}\varphi _{\beta 2}}\right)`$ $`=`$ $`r\mathrm{sin}\varphi _r.`$ (A21) Now Eq. (A20) implies that either $`\beta _1=0`$ or ($`\alpha _1=\delta _1`$ and $`2\varphi _{\beta 1}=\varphi _{\delta 1}+\pi `$). The second possibility is uninteresting, as it immediately leads to $`r=0`$, which gives $`I_H=0`$ — certainly not the optimum case, one hopes! Also, using the unitarity condition (A4) and the equal purity condition (A6), one obtains the restrictions $`2f`$ $`=`$ $`x+r(r1)\mathrm{cos}\varphi _r`$ (A23) $`+C\sqrt{1r^2}\sqrt{x(x2r\mathrm{cos}\varphi _r)},`$ $`r^2(1\mathrm{cos}^2\varphi _r)`$ $`=`$ $`2(1+r\mathrm{cos}\varphi _rx)`$ (A25) $`\times [xr\mathrm{cos}\varphi _r+K\sqrt{x(x2r\mathrm{cos}\varphi _r)}],`$ respectively. For brevity, the mutually independent parameters $`x,K,C`$ have been introduced, where $`x`$ $`=`$ $`\alpha _2^2,`$ (A27) $`K`$ $`=`$ $`\mathrm{cos}(\varphi _{\beta 2}+\varphi _{\gamma 2}\varphi _{\delta 2}),`$ (A28) $`C`$ $`=`$ $`\mathrm{cos}(\varphi _{\gamma 2}\varphi _{\gamma 1}).`$ (A29) Note that the condition (A25) is equivalent to Eq. (A21). Using Eqs. (II C), (A8), and (A10) leads to $`I_H`$ being given by the expression $`I_H`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(1+r)\mathrm{log}_2(1+r)+(1r)\mathrm{log}_2(1r)\right]`$ (A31) $`{\displaystyle \frac{1}{2}}\left[(1+q_H)\mathrm{log}_2(1+q_H)+(1q_H)\mathrm{log}_2(1q_H)\right],`$ with $$q_H=r\mathrm{cos}\frac{\varphi _r}{2}.$$ (A32) One finds that $`I_H(r,\mathrm{cos}\varphi _r)`$ is a monotonically decreasing function of $`\mathrm{cos}\varphi _r`$ — thus, to maximize $`I_H`$ for a given value of $`r=r_o`$, it suffices to minimize $`\mathrm{cos}\varphi _r`$ (i.e. make the angle between the possible copy Bloch vectors as close to $`\pi `$ as possible). $`I_H(r,\mathrm{cos}\varphi _r)`$ is also a monotonically increasing function of $`r`$. For any particular values of $`r`$ and $`\mathrm{cos}\varphi _r`$, there are three parameters left to vary to try to satisfy Eqs. (A23) and (A25), after the relations (A) have been used: $`x`$, $`K`$, and $`C`$. Each of the two Eqs. (A23), (A25) will give an allowable range for $`x`$ (exactly which point in these ranges is satisfied by the copier then depends on $`C`$ and $`K`$). The ends of these ranges are given by $$\frac{\mathrm{cos}\varphi _r}{C}=0\text{ or }C=\pm 1$$ (A34) for Eq. (A23), and $$\frac{\mathrm{cos}\varphi _r}{K}=0\text{ or }K=\pm 1$$ (A35) for Eq. (A25). Only those values of $`\mathrm{cos}\varphi _r`$ for which the two $`x`$ ranges partially overlap give allowable copiers. Now, for any particular $`r=r_o`$, if we vary $`\mathrm{cos}\varphi _r`$, the $`x`$ ranges will vary also. In particular, at that value of $`\mathrm{cos}\varphi _r`$ which lies at the boundary of allowed $`\mathrm{cos}\varphi _r(r_o)`$ values, at least one extremity of the first $`x`$ range, due to Eq. (A23), will coincide with an extremity of the second $`x`$ range due to Eq. (A25). Of course, not all cases where $`x`$ range extremities coincide will correspond to a $`\mathrm{cos}\varphi _r(r_o)`$ extremity, but any parameters for which such $`x`$ extremities coincide will give viable copiers \[they could be well within a region of allowed $`\mathrm{cos}\varphi _r(r_o)`$ values\]. Hence, if we look at all the parameters \[given by Eqs. (A23), (A25), and (A) \] where $`x`$ range extremities occur, then one of them will give the desired minimum $`\mathrm{cos}\varphi _r(r_o)`$ value. It turns out that this $`\mathrm{cos}\varphi _r(r_o)`$ minimum corresponds to $`K=C=1`$ when $`r[\sqrt{1f},1]`$. For $`r<\sqrt{1f}`$, $`\mathrm{cos}\varphi _r(r_o)`$ can reach its absolute minimum value of $`1`$, but since $`I_H`$ is also monotonically increasing in $`r`$, the optimum $`I_H`$ copier must have $`r\sqrt{1f}`$, so these low values of $`r`$ can be ignored. This leads to the second largest real root of polynomial (53) as the expression for $`\mathrm{cos}\varphi _r(r)`$ that maximizes $`I_H`$ for a given $`r\sqrt{1f}`$. The final value of $`r`$ that maximizes $`I_H`$ out of all the copiers considered, $`r_H`$, is given now by a straightforward, one-parameter maximization of $`I_H(r,\mathrm{cos}\varphi _r(r))`$ over $`r[\sqrt{1f},1]`$. Because this calculation is simple, straightforward, and accurate numerically, but not so simple analytically, an analytical solution has not been attempted. Now, to find the particular transformation which, given input states (III), not only maximizes $`I_H`$ but also makes the local copy-original fidelity as large as possible, first make the Bloch vectors of the copies be in the same plane as the Bloch vectors of the input states, and then make both pairs symmetric about a common axis. The Bloch vectors of the input states are $$𝐬_1=[\sqrt{f},0,\sqrt{1f}]\text{and}𝐬_2=[\sqrt{f},0,\sqrt{1f}].$$ (A36) These are in the $`(\widehat{\sigma }{}_{1}{}^{}\widehat{\sigma }{}_{3}{}^{})`$ plane, and symmetrically spaced about $`[1,0,0]`$. So, to achieve the desired optimum local fidelity copier, the appropriate transformation of the input states is found to be $$\|\psi _i^A(U_HU_H)\|\psi _i,$$ (A38) where $`\|\psi _i`$ is given by Eq. (A7), and the unitary transformations are $$U_H=\left(\begin{array}{cc}\mathrm{cos}\xi _H& \mathrm{sin}\xi _H\\ \mathrm{sin}\xi _H& \mathrm{cos}\xi _H\end{array}\right)\text{where}\xi _H=\frac{\varphi _r(r_H)\pi }{4}.$$ (A39) This can be written as Eqs. (39) and (40). ## B Derivation of Unentangled Optimal Copier Consider copiers producing product states of the copies. This transformation can be written $$\widehat{\rho }{}_{i}{}^{A}\widehat{\rho }{}_{i}{}^{B}\widehat{\rho }{}_{i}{}^{B}\widehat{\rho }{}_{i}{}^{x},$$ (B1) where $`\widehat{\rho }_i^B`$ are the copies and $`\widehat{\rho }_i^x`$ is a helper machine state. The only other constraint on the copier is that it must be unitary, which means that traces are preserved. This immediately leads to $`\widehat{\rho }_i^B`$ and $`\widehat{\rho }_i^x`$ being pure because the input states are pure (via $`\text{Tr}[\widehat{\rho }{}_{}{}^{2}]`$). Furthermore, $$f=\left(\text{Tr}[\widehat{\rho }{}_{1}{}^{B}\widehat{\rho }{}_{2}{}^{B}]\right)^2\text{Tr}[\widehat{\rho }{}_{1}{}^{x}\widehat{\rho }{}_{2}{}^{x}]=f_{12}^2f_x,$$ (B2) where $`f_{12}`$ and $`f_x`$ are the fidelities between, respectively, the two copy and two machine states produced after input of originals. Thus, since $`f_x1`$, it follows that $`\sqrt{f}f_{12}1`$. Let us start with optmizing for one-state information transfer $`I_1`$. It is easily shown that for equiprobable input states, $`I_1`$ satisfies Eq. (30) with the distinguishability parameter given by $$q=\sqrt{1f_{12}}.$$ (B3) This is most straightforward to show using the Bloch vectors of the copies. Since $`I_1`$ is monotonically increasing with $`q`$, it will be maximized when $`q`$ is maximized. This is when $`f_{12}=\sqrt{f}`$. Now let us look at $`I_H`$. For qubit copy states, this is again given by Eq. (A31), and since the copies are pure, $`r=1`$, and one finds $`q_H=\sqrt{f_{12}}`$. With $`r=1`$, $`I_H`$ depends only on $`q_H`$, and will reach extreme values either when $$\frac{dI_H}{dq_H}=\frac{1}{2}\mathrm{log}_2\left(\frac{1q_H}{1+q_H}\right)=0,$$ (B4) or at the end points of the $`q_H`$ range: $`q_H=(f^{1/4}\text{ or }1)`$. One sees that Eq. (B4) is only satisfied for $`q_H=f_{12}=f=0`$, so for general $`f`$, extreme values of $`I_H`$ are reached at $`f_{12}=1`$ or $`f_{12}=\sqrt{f}`$. $`f_{12}=1`$ leads to $`I_H=0`$, so the optimal value for $`f_{12}`$ is again $`\sqrt{f}`$. Thus the same copiers that are optimal in $`I_1`$ are also optimal in $`I_H`$. Lastly, let us look at local fidelity. The fidelity between any two pure states is given by $$F(\widehat{\rho }{}_{1}{}^{},\widehat{\rho }{}_{2}{}^{})=\frac{1}{2}(1+\mathrm{cos}\varphi ),$$ (B5) in terms of $`\varphi `$, the angle between their Bloch vectors. To minimize the average over both possible inputs of this Bloch angle between originals and copies, we choose the Bloch vectors of the copies to lie in the same plane as the Bloch vectors of the originals, and to be symmetric about the same axis. Obviously, in this case, the local fidelity will be maximized if the Bloch angle between the copies is as similar to the Bloch angle between the originals as possible (since the Bloch angle between original and copy is half the difference between these). Since $`f_{12}=\sqrt{f}f`$, this means that we want $`f_{12}=\sqrt{f}`$ again. Hence, the unentangled optimal copier given in Sec. IV C is optimal in all three indicators considered in this article. Choosing Bloch vector parameters such that Eq. (B1) holds, $`f_{12}=\sqrt{f}`$, and local fidelity is optimized, easily leads to the copier given in Eq. (IV C). It is simplest to use Bloch vectors for this calculation. ## C Some Fidelity-Optimized Copiers This section gives a brief summary of the fidelity-optimized copiers that are compared to the information-optimized ones in Sec. V. Expressions are given in terms of $`f`$, the square overlap between the two input states. Much more detail is given in the literature. ### 1 The copier that optimizes the global fidelity The quantum copying machine that optimizes the global fidelity between the combined state of both copies and a state consisting of unentangled perfect copies has been found by Bruß et al. The copies produced are (with the help of a little algebra) $`\widehat{\rho }_1^B`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1+\sqrt{{\displaystyle \frac{1f}{1+f}}}& {\displaystyle \frac{f+\sqrt{f}}{1+f}}\\ {\displaystyle \frac{f+\sqrt{f}}{1+f}}& 1\sqrt{{\displaystyle \frac{1f}{1+f}}}\end{array}\right),`$ (C4) $`\widehat{\rho }_2^B`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1\sqrt{{\displaystyle \frac{1f}{1+f}}}& {\displaystyle \frac{f+\sqrt{f}}{1+f}}\\ {\displaystyle \frac{f+\sqrt{f}}{1+f}}& 1+\sqrt{{\displaystyle \frac{1f}{1+f}}}\end{array}\right).`$ (C7) The local fidelity is \[from Eq. (47) in Ref. \] $$F(\widehat{\rho }{}_{i}{}^{A},\widehat{\rho }{}_{}{}^{B},i)=\frac{1}{2}(1+\frac{(1f)\sqrt{1+f}+f(1+\sqrt{f})}{1+f}),$$ (C8) and the one-state copied information is given by Eq. (30) with distinguishability parameter $$q=\sqrt{\frac{1f}{1+f}}.$$ (C9) The ultimate copied information is given by the expression (A31), where $`r`$, the magnitude of the Bloch vectors of the copies, is in this case $$r=\frac{\sqrt{1+f(1+2\sqrt{f})}}{1+f},$$ (C11) and the parameter $`q_H`$ is $$q_H=\frac{f+\sqrt{f}}{1+f}.$$ (C12) ### 2 The copier that optimizes the local fidelity As in Appendix C 1, Bruß et al. have found the copier that optimizes the local fidelity between a copy and the originals. From Eqs. (C1)-(C6), and(C12) and subsequent discussion in Ref. , the copies are in the states $`\widehat{\rho }_1^B`$ $`=`$ $`{\displaystyle \frac{\mathrm{sec}2\varphi }{2}}\left(\begin{array}{cc}\mathrm{cos}2\varphi +\sqrt{1f}& (1+\sqrt{f})\mathrm{sin}2\varphi \\ (1+\sqrt{f})\mathrm{sin}2\varphi & \mathrm{cos}2\varphi \sqrt{1f}\end{array}\right),`$ (C16) $`\widehat{\rho }_2^B`$ $`=`$ $`{\displaystyle \frac{\mathrm{sec}2\varphi }{2}}\left(\begin{array}{cc}\mathrm{cos}2\varphi \sqrt{1f}& (1+\sqrt{f})\mathrm{sin}2\varphi \\ (1+\sqrt{f})\mathrm{sin}2\varphi & \mathrm{cos}2\varphi +\sqrt{1f},\end{array}\right),`$ (C19) where the angle $`\varphi `$ is defined by $$\mathrm{sin}2\varphi =\frac{\sqrt{f}1+\sqrt{12\sqrt{f}+9f}}{4\sqrt{f}}.$$ (C20) The local fidelity is \[rearranging Eq. (C11) of Ref.\] $$F(\widehat{\rho }{}_{i}{}^{A},\widehat{\rho }{}_{i}{}^{B})=\frac{1}{2}\{1+\mathrm{cos}2\varphi [1f+\sqrt{f}(1+\sqrt{f})\mathrm{sin}2\varphi ]\}.$$ (C21) After some algebra, one finds that $`q`$ $`=`$ $`\sqrt{1f}\mathrm{cos}2\varphi ,`$ (C23) $`r`$ $`=`$ $`\mathrm{cos}2\varphi \sqrt{1f+(1+\sqrt{f})^2\mathrm{sin}^22\varphi },`$ (C24) $`q_H`$ $`=`$ $`\mathrm{sin}2\varphi \mathrm{cos}2\varphi (1+\sqrt{f}),`$ (C25) which can be used in expressions (30) and (A31), respectively, to find $`I_1`$ and $`I_H`$. ### 3 The UQCM The universal quantum copying machine copies any two-dimensional input states with an equal, optimal, local fidelity of $`5/6`$. This copier is unique among those mentioned in this article, in that it uses a machine helper state which becomes entangled with both copies after the process is complete. Given the input states (III) used in this article, the UQCM will create the copies $`\widehat{\rho }_1^B`$ $`=`$ $`{\displaystyle \frac{1}{6}}\left(\begin{array}{cc}3+\sqrt{1f}& 2\sqrt{f}\\ 2\sqrt{f}& 32\sqrt{1f}\end{array}\right),`$ (C29) $`\widehat{\rho }_1^B`$ $`=`$ $`{\displaystyle \frac{1}{6}}\left(\begin{array}{cc}3\sqrt{1f}& 2\sqrt{f}\\ 2\sqrt{f}& 3+2\sqrt{1f}\end{array}\right).`$ (C32) To calculate $`I_1`$ and $`I_H`$, use $`q`$ $`=`$ $`{\displaystyle \frac{2}{3}}\sqrt{1f},`$ (C34) $`r`$ $`=`$ $`{\displaystyle \frac{2}{3}},`$ (C35) $`q_H`$ $`=`$ $`{\displaystyle \frac{2}{3}}\sqrt{f}`$ (C36) in expressions (30) and (A31).
warning/0003/nlin0003029.html	ar5iv	text	# Contents ## 1 Introduction The superalgebra $`s\mathrm{}(2\|1)`$ appears in various quantum systems as underlying their symmetry and dynamics. Finite-dimensional representations describe spin states. For example, a lattice site which is allowed to be empty or occupied by an electron with the two spin states $`\pm \frac{1}{2}`$ , but not by two electrons, corresponds to the three-dimensional fundamental representation. Chains consisting of sites carrying the fundamental representation with integrable short-range interaction have been constructed. The Hamiltonian of $`tJ`$ model is obtained from the transfer matrix of the integrable model based on the three dimensional fundamental representation of $`s\mathrm{}(2\|1)`$ . The superalgebra $`s\mathrm{}(2\|1)`$ has a series of four-dimensional representations parametrized by a parameter $`b\pm \frac{1}{2}`$. Integrable models built from R-matrices defined on tensor products of two different four-dimensional representations have been considered in . The simplest integrable chain structure is the homogeneous periodic spin chain; most applications make use of this case. There exist some important modifications. The construction of open spin chains is well known . The treatment of integrable inhomogeneous spin chains is more involved. An integrable model has been constructed in view of the relevance for systems with impurities and in particular for the Kondo effect. The representations of $`s\mathrm{}(2\|1)`$ accommodate $`s\mathrm{}(2)`$ representations of spin $`s`$ and $`s\pm \frac{1}{2}`$ and allow in this way to construct chains with mobile impurities . The integrable chains turn out to describe approximately the effective interaction in four-dimensional gauge theories in the Regge limit and in the Bjorken limit . Unlike the above examples here the sites carry infinite-dimensional representations of $`s\mathrm{}(2)`$ accommodating all the momentum states of reggeons and partons. Besides of the case of homogeneous periodic chains also the case of open chains is encountered both in the Bjorken limit and in the Regge limit . A particular $`s\mathrm{}(2\|1)`$ representation of interest are the infinitesimal conformal transformations in one dimension together with their twofold supersymmetric extensions. This symmetry applies to the Bjorken limit of four-dimensional supersymmetric Yang-Mills theory at least up to one loop . This means, the one-loop renormalization of quasipartonic composite operators can be obtained by $`s\mathrm{}(2\|1)`$ symmetric pairwise interactions of partons, the states (light-cone momenta, helicity,fermion number) of which form an infinite-dimensional lowest weight module of this algebra. In the present paper we consider the algebra $`s\mathrm{}(2\|1)`$, its lowest weight modules and construct on this basis the solutions of the Yang-Baxter equation, i.e. the R-matrix acting on the tensor product of two those modules. R-matrices acting on tensor products of two fundamental $`s\mathrm{}(2\|1)`$-representations and of two four-dimensional representations have been constructed in the above mentioned papers . General integrable models based on R-matrices acting on tensor product of two arbitrary finite-dimensional atypical(chiral) representations of $`s\mathrm{}(2\|1)`$ have been constructed quite recently by generalizing the known approach from $`s\mathrm{}(2)`$ to the case of $`s\mathrm{}(2\|1)`$. We propose the alternative approach and generalize these results. Motivated by possible applications to the Bjorken limit in QCD, we represent the lowest weight modules by polynomials in one even ($`z`$) and two odd ($`\theta ,\overline{\theta }`$) variables. The general R-matrices are in fact operators acting on two-point functions, i.e. (polynomial) functions of two sets ($`z_1,\theta _1,\overline{\theta }_1`$) and ($`z_2,\theta _2,\overline{\theta }_2`$) representing the tensor product. We construct the R-operators acting on the tensor product of two arbitrary (finite or infinite-dimensional) $`s\mathrm{}(2\|1)`$-modules. This is done by calculating the two-point eigenfunctions of the lowest weight and the eigenvalues of the R-operator. From the particular result for arbitrary but isomorphic representations we derive the integrable nearest neighbour interaction Hamiltonian for homogeneous periodic chains with sites carrying arbitrary isomorphic representations. In the case of $`s\mathrm{}(2)`$ there exist integrable nearest neighbour interactions in open chains homogeneous inside but with arbitrary different representations corresponding to the end points . We extend this result to the case of $`s\mathrm{}(2\|1)`$ and construct the corresponding Hamiltonians. The presentation is organized as follows. In Section 2 we introduce definitions and summarize the standard facts about the superalgebra $`s\mathrm{}(2\|1)`$ and its representations. We represent the lowest weight modules by polynomials in one even ($`z`$) and two odd variables ($`\theta ,\overline{\theta }`$) and the $`s\mathrm{}(2\|1)`$-generators as first order differential operators. In Section 3 we derive the defining relation for the general R-matrix, i.e. the solution of the Yang-Baxter equation acting on tensor products of two arbitrary representations, the elements of which are polynomial functions of ($`z_1,\theta _1,\overline{\theta }_1`$) and ($`z_2,\theta _2,\overline{\theta }_2`$). We solve this defining relation in the space of lowest weights. In Section 4 we construct the local integrable Hamiltonians in the simplest case of homogeneous periodic chain carrying arbitrary isomorphic representations on the sites. In Section 5 we construct the local integrable Hamiltonians for the open chain with arbitrary isomorphic representations inside and other arbitrary representations at the end points. Finally, in Section 6 we summarize. Appendix A contains some technical details of calculations. In Appendix B we give the expression for the R-matrix acting in the tensor product of chiral modules. In Appendix C we discuss shortly the case of finite-dimensional representations and show that obtained general R-matrix reduces to the known ones for the tensor product of modules with minimal dimensions. ## 2 Algebra $`s\mathrm{}(2\|1)`$ ### 2.1 Commutators and Casimir Operators The superalgebra $`s\mathrm{}(2\|1)`$ has eight generators: four odd $`V^\pm ,W^\pm `$ and four even ones $`S,S^\pm `$ and $`B`$. The commutation relations have the following form : anticommutators $$\{V^\pm ,V^\pm \}=0;\{V^\pm ,V^{}\}=0;\{W^\pm ,W^\pm \}=0;\{W^\pm ,W^{}\}=0$$ $$\{V^\pm ,W^\pm \}=\pm S^\pm ;\{V^\pm ,W^{}\}=S\pm B,$$ commutators $$[S,S^\pm ]=\pm S^\pm ;[S^+,S^{}]=2S$$ $$[B,S^\pm ]=0;[B,S]=0;[S^\pm ,V^\pm ]=0;[S^\pm ,W^\pm ]=0$$ $$[B,V^\pm ]=\frac{1}{2}V^\pm ;[B,W^\pm ]=\frac{1}{2}W^\pm ;[S,V^\pm ]=\pm \frac{1}{2}V^\pm ;[S,W^\pm ]=\pm \frac{1}{2}W^\pm $$ $$[S^\pm ,V^{}]=V^\pm ;[S^\pm ,W^{}]=W^\pm .$$ These generators are linear combinations of the generators $`E_{AB}`$ of the superalgebra $`g\mathrm{}(2\|1)`$ . The commutation relations for the nine generators of $`g\mathrm{}(2\|1)`$ can be written compactly in the form: $$[E_{AB},E_{CD}]=\delta _{CB}E_{AD}()^{(\overline{A}+\overline{B})(\overline{C}+\overline{D})}\delta _{AD}E_{CB}$$ where the graded commutator is defined as: $$[E_{AB},E_{CD}]E_{AB}E_{CD}()^{(\overline{A}+\overline{B})(\overline{C}+\overline{D})}E_{CD}E_{AB}$$ The indices $`A,B,C,D=1,2,3`$ and we choose the grading:$`\overline{1}=\overline{3}=0;\overline{2}=1`$. The connection between both sets of generators is the following: $$S^{}=E_{31};W^{}=E_{21};V^{}=E_{32}$$ $$S^+=E_{13};W^+=E_{23};V^+=E_{12}$$ $$S=\frac{1}{2}E_{11}\frac{1}{2}E_{33};B=\frac{1}{2}E_{11}E_{22}\frac{1}{2}E_{33}$$ (2.1.1) In the fundamental representation all generators $`e_{AB}`$ of $`g\mathrm{}(2\|1)`$ are $`3\times 3`$-matrices and the basis in the space of these matrices can be chosen in the standard way: $$(e_{AB})_{CD}=\delta _{AC}\delta _{BD};e_{AB}e_{CD}=\delta _{CB}e_{AD}$$ (2.1.2) In the fundamental representation the generators $`W^\pm ,V^\pm ,S^\pm ,S,B`$ have the form: $$S^{}=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 1& 0& 0\end{array}\right);W^{}=\left(\begin{array}{ccc}0& 0& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right);V^{}=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 0& 1& 0\end{array}\right)$$ $$S^+=\left(\begin{array}{ccc}0& 0& 1\\ 0& 0& 0\\ 0& 0& 0\end{array}\right);W^+=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 0& 0\end{array}\right);V^+=\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)$$ $$S=\left(\begin{array}{ccc}\frac{1}{2}& 0& 0\\ 0& 0& 0\\ 0& 0& \frac{1}{2}\end{array}\right);B=\left(\begin{array}{ccc}\frac{1}{2}& 0& 0\\ 0& 1& 0\\ 0& 0& \frac{1}{2}\end{array}\right)$$ (2.1.3) There exists a simple construction for the central elements of the enveloping algebra of $`g\mathrm{}(2\|1)`$ . The first step is the construction of covariant operators: suppose we have two covariant operators $`V_{CD}^{(i)}`$, i.e. operators which have the following commutation relations with generators $$[E_{AB},V_{CD}^{(i)}]=\delta _{CB}V_{AD}^{(i)}()^{(\overline{A}+\overline{B})(\overline{C}+\overline{D})}\delta _{AD}V_{CB}^{(i)};i=1,2.$$ It is easy to check that operator $`V_{AB}=_C()^{\overline{C}}V_{AC}^{(1)}V_{CB}^{(2)}`$ is also covariant. This simple observation allows to construct covariant operators using generators $`E_{CD}`$ as simplest building blocks. The second step is the construction of a central element from the covariant operator: for any covariant operator $`V_{CD}`$ the operator $`V=_CV_{CC}`$ belongs to the center of the algebra $$[E_{AB},V]=0.$$ Repeating this construction we obtain central elements $`K_n,n=1,2,3\mathrm{}`$ for the enveloping algebra of $`g\mathrm{}(2\|1)`$: $$K_1=\underset{A}{}E_{AA};K_2=\underset{AB}{}()^{\overline{B}}E_{AB}E_{BA};K_3=\underset{ABC}{}()^{\overline{B}+\overline{C}}E_{AB}E_{BC}E_{CA};\mathrm{}$$ (2.1.4) The eight generators of algebra $`s\mathrm{}(2\|1)`$ may be introduced by defining: $$_{AB}E_{AB}\delta _{AB}()^{\overline{B}}\underset{A}{}E_{AA}$$ (2.1.5) $$_{31}=S^{};_{21}=W^{};_{32}=V^{}$$ $$_{13}=S^+;_{23}=W^+;_{12}=V^+$$ $$_{11}=BS;_{22}=2B;_{33}=B+S.$$ It can be verified that generators $`_{AB}`$ satisfy the same commutation relations as $`E_{AB}`$: $$[_{AB},_{CD}]=\delta _{CB}_{AD}()^{(\overline{A}+\overline{B})(\overline{C}+\overline{D})}\delta _{AD}_{CB}$$ There exists only one restriction for the Cartan generators: $`_{11}+_{22}+_{33}=0`$ and the independent eight generators can be chosen in the form (2.1.1). The center of algebra $`s\mathrm{}(2\|1)`$ is generated by Casimir operators $`C_n,n=2,3\mathrm{}`$ . We shall use only two of them: $$C_2=\frac{1}{2!}\underset{AB}{}()^{\overline{B}}_{AB}_{BA}=S^2B^2+S^+S^{}+V^+W^{}+W^+V^{}$$ (2.1.6) $$C_3=\frac{1}{3!}\underset{ABC}{}()^{\overline{B}+\overline{C}}_{AB}_{BC}_{CA}=B(S^2B^2)+BS^+S^{}+\frac{3}{2}B(V^+W^{}+W^+V^{})+$$ (2.1.7) $$+\frac{1}{4}(W^+V^+V^+W^+)S^{}+\frac{1}{4}S^+(V^{}W^{}W^{}V^{})+\frac{1}{2}(S1)(V^+W^{}W^+V^{}).$$ ### 2.2 Global form of superconformal transformations We represent the generators as first order differential operators, acting on the space of polynomials $`\mathrm{\Phi }(z,\theta ,\overline{\theta })`$. Lowering(decreasing the polynomial degree) operators have the form $$S^{}=;V^{}=_\theta +\frac{1}{2}\overline{\theta };W^{}=_{\overline{\theta }}+\frac{1}{2}\theta $$ (2.2.1) and generate the following global transformations $$e^{\lambda S^{}}\mathrm{\Phi }(z;\theta ,\overline{\theta })=\mathrm{\Phi }(z\lambda ;\theta ,\overline{\theta }),$$ (2.2.2) $$e^{ϵV^{}}\mathrm{\Phi }(z;\theta ,\overline{\theta })=\mathrm{\Phi }(z+\frac{ϵ\overline{\theta }}{2};\theta +ϵ,\overline{\theta });e^{ϵW^{}}\mathrm{\Phi }(z;\theta ,\overline{\theta })=\mathrm{\Phi }(z+\frac{ϵ\theta }{2};\theta ,\overline{\theta }+ϵ).$$ Rising(increasing the polynomial degree) operators $$V^+=\left[z_\theta +\frac{1}{2}\overline{\theta }z+\frac{1}{2}\overline{\theta }\theta _\theta \right](\mathrm{}b)\overline{\theta };W^+=\left[z_{\overline{\theta }}+\frac{1}{2}\theta z+\frac{1}{2}\theta \overline{\theta }_{\overline{\theta }}\right](\mathrm{}+b)\theta ,$$ $$S^+=z^2+z\theta _\theta +z\overline{\theta }_{\overline{\theta }}+2\mathrm{}zb\theta \overline{\theta },$$ (2.2.3) generate the global transformations $$e^{\lambda S^+}\mathrm{\Phi }(z;\theta ,\overline{\theta })=\left[1+\frac{\theta \overline{\theta }\lambda }{(1\lambda z)}\right]^b\frac{1}{(1\lambda z)^2\mathrm{}}\mathrm{\Phi }(\frac{z}{1\lambda z};\frac{\theta }{1\lambda z},\frac{\overline{\theta }}{1\lambda z}),$$ $$e^{ϵV^+}\mathrm{\Phi }(z;\theta ,\overline{\theta })=\frac{1}{(1+ϵ\overline{\theta })^\mathrm{}b}\mathrm{\Phi }(\frac{z}{1+\frac{ϵ\overline{\theta }}{2}};\frac{\theta ϵz}{1+\frac{ϵ\overline{\theta }}{2}},\overline{\theta }),$$ $$e^{ϵW^+}\mathrm{\Phi }(z;\theta ,\overline{\theta })=\frac{1}{(1+ϵ\theta )^{\mathrm{}+b}}\mathrm{\Phi }(\frac{z}{1+\frac{ϵ\theta }{2}};\theta ,\frac{\overline{\theta }ϵz}{1+\frac{ϵ\theta }{2}}).$$ Two remaining elements of the Cartan subalgebra: $$S=z+\frac{1}{2}\theta _\theta +\frac{1}{2}\overline{\theta }_{\overline{\theta }}+\mathrm{};B=\frac{1}{2}\overline{\theta }_{\overline{\theta }}\frac{1}{2}\theta _\theta +b$$ generate the transformations: $$e^{\lambda S}\mathrm{\Phi }(z;\theta ,\overline{\theta })=e^\mathrm{}\lambda \mathrm{\Phi }(e^\lambda z;e^{\frac{\lambda }{2}}\theta ,e^{\frac{\lambda }{2}}\overline{\theta })$$ $$e^{\lambda B}\mathrm{\Phi }(z;\theta ,\overline{\theta })=e^{b\lambda }\mathrm{\Phi }(z;e^{\frac{\lambda }{2}}\theta ,e^{\frac{\lambda }{2}}\overline{\theta })$$ We use a natural convention here and assign scaling dimension one and $`U(1)`$-charge zero to the even variable $`z`$ and the scaling dimension $`\frac{1}{2}`$ and $`U(1)`$-charge $`\frac{1}{2}`$ to the odd variables $`\theta `$ and $`\overline{\theta }`$. ### 2.3 $`s\mathrm{}(2\|1)`$-lowest weight modules The lowest weight $`s\mathrm{}(2\|1)`$-module $`V_{\mathrm{},b}=V_{\stackrel{}{\mathrm{}}},\stackrel{}{\mathrm{}}=(\mathrm{},b)`$ is built on the lowest weight vector $`\psi `$ obeying: $$V_{}\psi =0;W^{}\psi =0;S^{}\psi =0;S\psi =\mathrm{}\psi ;B\psi =b\psi $$ In generic situation $`\mathrm{}\pm b`$ the module is characterised uniquely by the action of the Casimir operators on its elements: $$C_2v=(\mathrm{}^2b^2)v;C_3v=b(\mathrm{}^2b^2)v;vV_{\mathrm{},b}$$ It is a vector space spanned by the following basis with the even vectors $$A_k=(S^+)^k\psi ;B_k=(S^+)^{k1}W^+V^+\psi ,k𝒵_+$$ and the odd vectors $$V_k=(S^+)^kV^+\psi ;W_k=(S^+)^kW^+\psi ,k𝒵_+$$ We shall use the above realization of the $`s\mathrm{}(2\|1)`$-generators as the differential operators of first order acting on the infinite-dimensional(for generic $`\mathrm{}`$) space $`V_{\mathrm{},b}`$ of polynomials $`\mathrm{\Phi }(z,\theta ,\overline{\theta })`$ of variables $`z,\theta ,\overline{\theta }`$. It is easy to obtain the expression for the coherent states $`e^{\lambda S^+}\psi ,e^{\lambda S^+}W^+V^+\psi `$ and $`e^{\lambda S^+}V^+\psi ,e^{\lambda S^+}W^+\psi `$ for the lowest weight $`\psi =1`$ using the formulae of the previous section. They are the generating functions , the power expansion in $`\lambda `$ of which produces the basis: $$A_k=(2\mathrm{})_k\left[z^k\frac{kb}{2\mathrm{}}z^{k1}\theta \overline{\theta }\right];B_k=\frac{\mathrm{}b}{2\mathrm{}}(2\mathrm{})_k\left[z^k+\left(b+\mathrm{}+\frac{k}{2}\right)z^{k1}\theta \overline{\theta }\right]$$ (2.3.1) $$V_k=(\mathrm{}b)(2\mathrm{}+1)_kz^k\overline{\theta };W_k=(\mathrm{}+b)(2\mathrm{}+1)_kz^k\theta ;(2\mathrm{})_k\frac{\mathrm{\Gamma }(2\mathrm{}+k)}{\mathrm{\Gamma }(2\mathrm{})}$$ It is useful to introduce the subspaces of functions with definite chirality. Let us define for this purpose two operators called supercovariant derivatives: $$D^+=_\theta +\frac{1}{2}\overline{\theta };D^{}=_{\overline{\theta }}+\frac{1}{2}\theta $$ (2.3.2) and two subspaces $`V_{\mathrm{},b}^\pm \mathrm{ker}D^\pm V_{\mathrm{},b}`$: $$\mathrm{\Phi }(z,\theta ,\overline{\theta })V_{\mathrm{},b}^+\mathrm{\Phi }(z,\theta ,\overline{\theta })=\mathrm{\Phi }(z_+,\theta _+);z_+z+\frac{1}{2}\theta \overline{\theta },\theta _+\overline{\theta }$$ $$\mathrm{\Phi }(z,\theta ,\overline{\theta })V_{\mathrm{},b}^{}\mathrm{\Phi }(z,\theta ,\overline{\theta })=\mathrm{\Phi }(z_{},\theta _{});z_{}z\frac{1}{2}\theta \overline{\theta },\theta _{}\theta $$ In the generic case the chiral subspaces $`V_{\mathrm{},b}^\pm `$ are not $`s\mathrm{}(2\|1)`$-invariant ones. Indeed, the operators $`D^\pm `$ have the following commutation relations with $`s\mathrm{}(2\|1)`$-generators: $$\{D^\pm ,V^{}\}=0;\{D^\pm ,W^{}\}=0;[D^\pm ,S^{}]=0$$ $$[D^\pm ,S]=\frac{1}{2}D^\pm ;[D^\pm ,B]=\frac{1}{2}D^\pm $$ (2.3.3) $$\{D^+,V^+\}=\overline{\theta }D^+;\{D^+,W^+\}=\mathrm{}+b;[D^+,S^+]=(\mathrm{}+b)\overline{\theta }+z_+D^+$$ $$\{D^{},V^+\}=\mathrm{}b;\{D^{},W^+\}=\theta D^{};[D^{},S^+]=(\mathrm{}b)\theta +z_{}D^{},$$ and it is easy to see that chiral subspaces $`V_{\mathrm{},b}^\pm `$ are $`s\mathrm{}(2\|1)`$-invariant only under the condition $`\mathrm{}=b`$. In this case the whole module $`V_{\mathrm{},\pm \mathrm{}}`$ has definite chirality $`V_{\mathrm{},\pm \mathrm{}}=V_{\mathrm{},\pm \mathrm{}}^{}`$: $$D^{}v=0,vV_\mathrm{},\mathrm{};D^+v=0,vV_\mathrm{},\mathrm{}.$$ The notions of chirality and chiral representations are used here as they are common in supersymmetric field theory. In the mathematical literature about superalgebra representations generic representations are called typical and chiral representations are called atypical . There exist some special values of $`\mathrm{}`$: $`\mathrm{}=n;n\frac{1}{2}𝒵_+`$ , for which the module $`V_{\mathrm{},b}`$ becomes a finite-dimensional vector space . Indeed it is evident from (2.3.1) that all basis vectors are equal zero for $`kn+1`$. There are three cases depending on the relation between $`b`$ and $`n`$. The first case is for generic $`b`$(typical representations): $`b\pm n;dimV_{n,b}=8n`$ $$\mathrm{\Phi }_k^\pm =z_\pm ^k,k=1\mathrm{}2n1,\mathrm{\Phi }_0=1,\mathrm{\Phi }_{2n}=\left(z+\frac{b}{2n}\theta \overline{\theta }\right)^{2n};\mathrm{\Psi }_k^\pm =\theta _\pm z_\pm ^k,k=0\mathrm{}2n1$$ The second and third case appears for $`b=\pm n`$ and here the representation spaces have definite chirality (atypical representations). $`b=\pm n;V_{n,\pm n}=V_{n,n}^\pm ;dimV_{n,\pm n}=4n+1`$ $$\mathrm{\Phi }_k^\pm =z_\pm ^k,k=0\mathrm{}2n;\mathrm{\Psi }_k^\pm =\theta _\pm z_\pm ^k,k=0\mathrm{}2n1$$ Let us introduce the special notation for the fundamental $`s\mathrm{}(2\|1)`$-module: $`V_{\frac{1}{2},\frac{1}{2}}V_\stackrel{}{f},\stackrel{}{f}=(\frac{1}{2},\frac{1}{2})`$. In the basis $$e_1=z_{}\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right);e_2=\theta _{}\left(\begin{array}{ccc}0& & \\ 1& & \\ 0& & \end{array}\right);e_3=1\left(\begin{array}{ccc}0& & \\ 0& & \\ 1& & \end{array}\right)$$ the $`s\mathrm{}(2\|1)`$-generators take their fundamental form (2.1.3). ### 2.4 Tensor products of two $`s\mathrm{}(2\|1)`$-modules The tensor product of two $`s\mathrm{}(2\|1)`$-modules has the following direct sum decomposition : $$V_{\mathrm{}_1,b_1}V_{\mathrm{}_2,b_2}=V_{\mathrm{},b}+2\underset{n=1}{\overset{\mathrm{}}{}}V_{\mathrm{}+n,b}+\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n+\frac{1}{2},b\frac{1}{2}}+\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n+\frac{1}{2},b+\frac{1}{2}};\mathrm{}_i\pm b_i$$ (2.4.1) $$\mathrm{}=\mathrm{}_1+\mathrm{}_2;b=b_1+b_2$$ Note that this formula is applicable in the generic situation $`\mathrm{}_i\pm b_i`$. The direct sum decomposition reduces for the tensor product involving chiral modules: $$V_{\mathrm{}_1,\mathrm{}_1}V_{\mathrm{}_2,b_2}=\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n,b}+\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n+\frac{1}{2},b\frac{1}{2}};\mathrm{}_2\pm b_2$$ $$V_{\mathrm{}_1,\mathrm{}_1}V_{\mathrm{}_2,\pm \mathrm{}_2}=\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n,b};V_{\mathrm{}_1,\mathrm{}_1}V_{\mathrm{}_2,\mathrm{}_2}=\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n+\frac{1}{2},b\frac{1}{2}}$$ In Appendix C we discuss the modifications of (2.4.1) arising for finite-dimensional representations . For the proof of (2.4.1) one has to determine all possible lowest weight vectors appearing in the tensor product $`V_{\mathrm{}_1,b_1}V_{\mathrm{}_2,b_2}`$. In the realization on functions of $`z,\theta ,\overline{\theta }`$ the space $`V_{\mathrm{}_1,b_1}V_{\mathrm{}_2,b_2}`$ is isomorphic to the space of polynomials on two even variables $`z_1,z_2`$ and four odd variables $`\theta _1,\overline{\theta }_1,\theta _2,\overline{\theta }_2`$ called for the sake of brevity two-point functions. The $`s\mathrm{}(2\|1)`$-generators acting on the $`V_{\mathrm{}_1,b_1}V_{\mathrm{}_2,b_2}`$ are just the sums of corresponding generators acting in $`V_{\mathrm{}_i,b_i}`$. The lowest weight vectors of the irreducible representations in the decomposition of $`V_{\mathrm{}_1,b_1}V_{\mathrm{}_2,b_2}`$ are defined as the common solutions of the equations: $$S^{}\mathrm{\Phi }=V^{}\mathrm{\Phi }=W^{}\mathrm{\Phi }=0$$ (2.4.2) which have the form $$\mathrm{\Phi }(z_1,z_2;\theta _1,\theta _2;\overline{\theta }_1,\overline{\theta }_2)=\mathrm{\Phi }(Z_{12},\theta _{12},\overline{\theta }_{12})$$ (2.4.3) where $$Z_{12}z_1z_2+\frac{1}{2}(\overline{\theta }_1\theta _2+\theta _1\overline{\theta }_2);\theta _{12}\theta _1\theta _2;\overline{\theta }_{12}\overline{\theta }_1\overline{\theta }_2.$$ Indeed, from (2.4.2) follows immediately that the function $`\mathrm{\Phi }`$ has to be invariant with respect to global transformations (2.2.2). This invariance predicts the general form of $`\mathrm{\Phi }`$: $$e^{\beta W^{}}e^{\alpha V^{}}e^{aS^{}}\mathrm{\Phi }(z_1,z_2;\theta _1,\theta _2;\overline{\theta }_1,\overline{\theta }_2)=$$ $$=\mathrm{\Phi }(z_1a+\frac{\alpha (\overline{\theta }_1+\beta )}{2}+\frac{\beta \theta _1}{2},z_2a+\frac{\alpha (\overline{\theta }_2+\beta )}{2}+\frac{\beta \theta _2}{2};\theta _1+\alpha ,\theta _2+\alpha ;\overline{\theta }_1+\beta ,\overline{\theta }_2+\beta )$$ and choosing $`a=z_2+\frac{1}{2}\theta _2\overline{\theta }_2,\alpha =\theta _2,\beta =\overline{\theta }_2`$ we obtain (2.4.3). There are additional restrictions of definite chirality for the lowest weights in the tensor product of the chiral modules: $$D_1^\pm \mathrm{\Phi }=0\mathrm{\Phi }=\mathrm{\Phi }(Z_{12}\pm \frac{1}{2}\theta _{12}\overline{\theta }_{12},\theta _{12}^\pm );\theta _{12}^\pm \theta _1^\pm \theta _2^\pm $$ $$D_2^\pm \mathrm{\Phi }=0\mathrm{\Phi }=\mathrm{\Phi }(Z_{12}\frac{1}{2}\theta _{12}\overline{\theta }_{12},\theta _{12}^\pm ).$$ Now, the lowest weight vectors in the decomposition of the tensor product are constructed from functions (2.4.3) being eigenfunctions of generators $`S`$ and $`B`$. The eigenfunctions of the operator $`S`$ are the polynomials with scaling dimension $`n`$ and the eigenfunctions of the operator $`B`$ are the polynomials with one of the possible $`U(1)`$-charges: $`0,\pm \frac{1}{2}`$. Finally we obtain that all lowest weights in the space $`V_{\mathrm{}_1,b_1}V_{\mathrm{}_2,b_2}`$ are divided on two sets, the even lowest weights: $$\mathrm{\Phi }_n^\pm \left(Z_{12}\pm \frac{1}{2}\theta _{12}\overline{\theta }_{12}\right)^n;D_1^\pm \mathrm{\Phi }^\pm =0,S\mathrm{\Phi }_n^\pm =(n+\mathrm{})\mathrm{\Phi }_n^\pm ,B\mathrm{\Phi }_n^\pm =b\mathrm{\Phi }_n^\pm $$ (2.4.4) and the odd lowest weights: $$\mathrm{\Psi }_n^{}\theta _{12}Z_{12}^n;\mathrm{\Psi }_n^+\overline{\theta }_{12}Z_{12}^n;S\mathrm{\Psi }_n^\pm =(n+\mathrm{}+\frac{1}{2})\mathrm{\Psi }_n^\pm ,B\mathrm{\Psi }_n^\pm =(b\pm \frac{1}{2})\mathrm{\Psi }_n^\pm $$ (2.4.5) It is convenient to choose the chiral basis $`D_1^\pm \mathrm{\Phi }_n^\pm =0`$ for the even lowest weights. Thus we have obtained the full set of lowest weights appearing in the expansion of $`V_{\mathrm{}_1,b_1}V_{\mathrm{}_2,b_2}`$ and this proves the direct sum decomposition in generic situation (2.4.1). The modifications in the case of the tensor product of chiral modules are evident: the lowest weights in the tensor product of the chiral modules are obtained by imposing the chirality restrictions. All lowest weights in the space $`V_{\mathrm{}_1,\mathrm{}_1}V_{\mathrm{}_2,b_2},\mathrm{}_2\pm b_2`$ are $`\mathrm{\Phi }_n^\pm `$ and $`\mathrm{\Psi }_n^\pm `$ and the lowest weights in the space $`V_{\mathrm{}_1,\mathrm{}_1}V_{\mathrm{}_2,\pm \mathrm{}_2}`$ are $`\mathrm{\Phi }_n^\pm `$. All lowest weights in the space $`V_{\mathrm{}_1,\mathrm{}_1}V_{\mathrm{}_2,\mathrm{}_2}`$ are $`\theta _{12}(z_1^{}z_2^{})^n=\mathrm{\Psi }_n^{}`$ and all lowest weights in the space $`V_{\mathrm{}_1,\mathrm{}_1}V_{\mathrm{}_2,\mathrm{}_2}`$ are $`\overline{\theta }_{12}(z_1^+z_2^+)^n=\mathrm{\Psi }_n^+`$. ## 3 Yang-Baxter equation and general operator $`\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)`$ Let $`V_{\mathrm{}_i,b_i};i=1,2,3`$ be three lowest weight $`s\mathrm{}(2\|1)`$-modules. We shall use the short-hand notation: $$\stackrel{}{\mathrm{}}=(\mathrm{},b);V_{\stackrel{}{\mathrm{}}}=V_{\mathrm{},b}$$ Let us consider the three operators $`\text{}_{\stackrel{}{\mathrm{}}_i\stackrel{}{\mathrm{}}_j}(u)`$ which are acting in $`V_\stackrel{}{\mathrm{}}_iV_\stackrel{}{\mathrm{}}_j`$ and obey the Yang-Baxter equation in the space $`V_\stackrel{}{\mathrm{}}_1V_\stackrel{}{\mathrm{}}_2V_\stackrel{}{\mathrm{}}_3`$ : $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(uv)\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_3}(u)\text{}_{\stackrel{}{\mathrm{}}_2\stackrel{}{\mathrm{}}_3}(v)=\text{}_{\stackrel{}{\mathrm{}}_2\stackrel{}{\mathrm{}}_3}(v)\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_3}(u)\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(uv)$$ (3.0.1) We are going to find the general solution $`\text{}_{\stackrel{}{\mathrm{}}_i\stackrel{}{\mathrm{}}_j}(u)`$ of Yang-Baxter equation by the following three steps. First one obtains the operator $`\text{}_{\stackrel{}{f},\stackrel{}{f}}(u)`$ in the simplest situation: $`\stackrel{}{\mathrm{}}_i=\stackrel{}{f}`$ for all $`i=1,2,3`$ so that the space $`V_\stackrel{}{f}=V_{\frac{1}{2},\frac{1}{2}}`$ has the minimal possible dimension: $`dimV_\stackrel{}{f}=3`$. In the second step we fix $`\stackrel{}{\mathrm{}}_i=\stackrel{}{f}`$ for $`i=1,2`$ and obtain the solution $`\text{}_\stackrel{}{f}\stackrel{}{\mathrm{}}`$ for arbitrary $`\stackrel{}{\mathrm{}}`$. In the third step we fix $`\stackrel{}{\mathrm{}}_3=\stackrel{}{f}`$ and using the result for the operator $`\text{}_{\stackrel{}{f}\stackrel{}{\mathrm{}}_1}`$ we obtain and solve the defining equation for the general R-matrix $`\text{}_{\stackrel{}{\mathrm{}}_i\stackrel{}{\mathrm{}}_j}(u)`$. It should be noted that the analogous approach was used for the derivation of the $`s\mathrm{}(2)`$-invariant R-matrix . ### 3.1 Fundamental solution $`\text{}_{\stackrel{}{f},\stackrel{}{f}}`$ First one considers the simplest situation: $`\mathrm{}_i=\frac{1}{2},b_i=\frac{1}{2}\stackrel{}{\mathrm{}}_i=\stackrel{}{f}_i`$. We shall prove that the operator: $$\text{}_{\stackrel{}{f}_i\stackrel{}{f}_j}(u)=u+\eta P_{ij};P_{ij}\underset{AB}{}()^{\overline{B}}e_{AB}^ie_{BA}^j$$ (3.1.1) where $`e_{AB}^i`$ are generators acting in the space $`V_{\stackrel{}{f}_i}`$, is the solution of the Yang-Baxter equation : $$\text{}_{\stackrel{}{f}_1\stackrel{}{f}_2}(uv)\text{}_{\stackrel{}{f}_1\stackrel{}{f}_3}(u)\text{}_{\stackrel{}{f}_2\stackrel{}{f}_3}(v)=\text{}_{\stackrel{}{f}_2\stackrel{}{f}_3}(v)\text{}_{\stackrel{}{f}_1\stackrel{}{f}_3}(u)\text{}_{\stackrel{}{f}_1\stackrel{}{f}_2}(uv)$$ Indeed, the proof is the following. The Yang-Baxter equation has the simple form in short notations: $$(uv+\eta P_{12})(u+\eta P_{13})(v+\eta P_{23})=(v+\eta P_{23})(u+\eta P_{13})(uv+\eta P_{12})$$ Comparing operator coefficients of $`u^k`$ on both sides of this equation yields: $$u^0:P_{12}P_{13}P_{23}=P_{23}P_{13}P_{12}$$ (3.1.2) $$u^1:P_{13}P_{23}+P_{12}P_{23}=P_{23}P_{12}+P_{23}P_{13}$$ (3.1.3) Using (2.1.2) one can prove that the operator $`P_{ij}`$ is the permutation: $$P_{ij}e_{AB}^i=e_{AB}^jP_{ij}P_{ij}P_{jk}=P_{ik}P_{ij}$$ and this commutation relation for $`P_{ij}`$ allows to check that eqs. (3.1.2), (3.1.3) hold and this proves that $`\text{}_{\stackrel{}{f}\stackrel{}{f}}`$ (3.1.1) obeys the Yang-Baxter equation. ### 3.2 The solution for the operator $`\text{}_{\stackrel{}{f},\stackrel{}{\mathrm{}}}(u)`$ We fix $`\stackrel{}{\mathrm{}}_i=\stackrel{}{f}_i`$ for $`i=1,2`$ and obtain the solution $`\text{}_\stackrel{}{f}\stackrel{}{\mathrm{}}`$ for arbitrary $`\stackrel{}{\mathrm{}}`$. The operator: $$\text{}_\stackrel{}{f}\stackrel{}{\mathrm{}}(u)=u+\eta \underset{AB}{}()^{\overline{B}}e_{AB}E_{BA},$$ where $`E_{AB}`$ are generators in arbitrary representation $`\stackrel{}{\mathrm{}}`$, is the solution of the Yang-Baxter equation: $$\text{}_{\stackrel{}{f}_1\stackrel{}{f}_2}(uv)\text{}_{\stackrel{}{f}_1\stackrel{}{\mathrm{}}_3}(u)\text{}_{\stackrel{}{f}_2\stackrel{}{\mathrm{}}_3}(v)=\text{}_{\stackrel{}{f}_2\stackrel{}{\mathrm{}}_3}(v)\text{}_{\stackrel{}{f}_1\stackrel{}{\mathrm{}}_3}(u)\text{}_{\stackrel{}{f}_1\stackrel{}{f}_2}(uv)$$ The proof is the following . The Yang-Baxter equation has the simple form in short-hand notations: $$(uv+\eta P_{12})(u+\eta e1E)(v+\eta 1eE)=(v+\eta 1eE)(u+\eta e1E)(uv+\eta P_{12})$$ Matching operator coefficients of $`u^k`$ on both sides of this equality yields: $$u^0:P_{12}(e1E)(1eE)=(1eE)(e1E)P_{12}$$ $$u^1:(e1E)(1eE)+P_{12}(1eE)=(1eE)P_{12}+(1eE)(e1E)$$ The first equation is a simple consequence of the properties of $`P_{12}`$. Using the fact that the matrices $`e_{AB}`$ form a basis we obtain that the second equation is equivalent to the following system of equations: $$E_{AB}E_{CD}()^{(\overline{A}+\overline{B})(\overline{C}+\overline{D})}E_{CD}E_{AB}=\delta _{CB}E_{AD}()^{(\overline{A}+\overline{B})(\overline{C}+\overline{D})}\delta _{AD}E_{CB}.$$ This is nothing else but the commutation relations for generators $`E_{AB}`$. Let us represent the operator under consideration: $$\text{}_\stackrel{}{f}\stackrel{}{\mathrm{}}(u)=u+\eta \underset{AB}{}()^{\overline{B}}e_{AB}E_{BA}$$ in the matrix form in the standard basis with the grading $`\overline{1},\overline{3}=0;\overline{2}=1`$: $$e_1=\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right);e_2=\left(\begin{array}{ccc}0& & \\ 1& & \\ 0& & \end{array}\right);e_3=\left(\begin{array}{ccc}0& & \\ 0& & \\ 1& & \end{array}\right)$$ We shall use the following definition of the matrix of an operator: $$Fe_A=e_BF_{BA},$$ (3.2.1) which leads to the validity of the common rule for the matrix product, i.e. without additional sign factors, $$FGe_A=e_B(FG)_{BA};(FG)_{BA}=F_{BC}G_{CA}.$$ Let us calculate the matrix of R-operator using the definitions (3.2.1) and (2.1.2): $$(e_{AB}E_{BA})e_C=()^{(\overline{B}+\overline{A})\overline{C}}e_{AB}e_CE_{BA}=()^{(\overline{B}+\overline{A})\overline{C}}e_D(e_{AB})_{DC}E_{BA}.$$ Therefore the matrix element of operator $`e_{AB}E_{BA}`$ has the form $$\underset{AB}{}()^{\overline{B}}(e_{AB}E_{BA})_{CD}=()^{\overline{D}\overline{C}}E_{DC}.$$ Finally one obtains $$\text{}_\stackrel{}{f}\stackrel{}{\mathrm{}}(u)=u+\eta \underset{AB}{}()^{\overline{B}}e_{AB}E_{BA}=\left(\begin{array}{ccc}u+\eta E_{11}& \eta E_{21}& \eta E_{31}\\ \eta E_{12}& u\eta E_{22}& \eta E_{32}\\ \eta E_{13}& \eta E_{23}& u+\eta E_{33}\end{array}\right).$$ This is the expression for the $`g\mathrm{}(2\|1)`$-invariant R-matrix. The $`s\mathrm{}(2\|1)`$-invariant R-matrix can be derived from this result in a simple way: the operator $`K_1=E_{11}+E_{22}+E_{33}`$ belongs to the center of the algebra and therefore the R-operator $`\text{}_\stackrel{}{f}\stackrel{}{\mathrm{}}(u\eta K_1)`$ is also a solution of the Yang-Baxter equation. Using the connection between $`E_{AB}`$ and $`s\mathrm{}(2\|1)`$-generators we obtain the $`s\mathrm{}(2\|1)`$-invariant R-matrix $$\text{}_\stackrel{}{f}\stackrel{}{\mathrm{}}(u\eta K_1)=\left(\begin{array}{ccc}u+\eta (S+B)& \eta W^{}& \eta S^{}\\ \eta V^+& u+2\eta B& \eta V^{}\\ \eta S^+& \eta W^+& u+\eta (BS)\end{array}\right).$$ ### 3.3 General R-matrix $`\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)`$ To obtain the defining relation for the general -operator we consider the special case $`\stackrel{}{\mathrm{}}_3=\stackrel{}{f}`$ in (3.0.1). Then one can choose the above matrix realization in $`V_\stackrel{}{\mathrm{}}_3`$ and the operators $`\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{f}}`$,$`\text{}_{\stackrel{}{\mathrm{}}_2\stackrel{}{f}}`$ are linear functions of spectral parameter $`u`$ in this particular case $$\text{}_{\stackrel{}{\mathrm{}}_i\stackrel{}{f}}(u\eta K_1)=u+\eta \text{}_i;\text{}_i=\left(\begin{array}{ccc}S_i+B_i& W_i^{}& S_i^{}\\ V_i^+& 2B_i& V_i^{}\\ S_i^+& W_i^+& B_iS_i\end{array}\right);i=1,2$$ Now the general R-matrix $`\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)`$ acting in the tensor product $`V_\stackrel{}{\mathrm{}}_1V_\stackrel{}{\mathrm{}}_2`$ of arbitrary modules, is fixed by the condition $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(uv)\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{f}}(u\eta K_1)\text{}_{\stackrel{}{\mathrm{}}_2\stackrel{}{f}}(v\eta K_1)=\text{}_{\stackrel{}{\mathrm{}}_2\stackrel{}{f}}(v\eta K_1)\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{f}}(u\eta K_1)\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(uv)$$ (3.3.1) or equivalently: $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(uv)\left(\frac{uv}{\eta ^2}+\frac{u+v}{2\eta }(\text{}_1+\text{}_2)+\frac{uv}{2\eta }(\text{}_2\text{}_1)+\text{}_1\text{}_2\right)=$$ $$=\left(\frac{uv}{\eta ^2}+\frac{v+u}{2\eta }(\text{}_2+\text{}_1)+\frac{vu}{2\eta }(\text{}_1\text{}_2)+\text{}_2\text{}_1\right)\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(uv).$$ After separation of $`u+v`$ and $`uv`$ dependence we obtain two equations($`uvu`$): $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)(\text{}_1+\text{}_2)=(\text{}_1+\text{}_2)\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)$$ (3.3.2) $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)\left(\frac{u}{2\eta }(\text{}_2\text{}_1)+\text{}_1\text{}_2\right)=\left(\frac{u}{2\eta }(\text{}_2\text{}_1)+\text{}_2\text{}_1\right)\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)$$ (3.3.3) The first equation (3.3.2) expresses the fact that $`\text{}(u)`$ has to be invariant with respect to the action of $`s\mathrm{}(2\|1)`$-algebra and the second equation is the wanted defining relation for the operator $`\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)`$. The $`s\mathrm{}(2\|1)`$-invariance of the operator $`\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)`$ allows to simplify the problem. Due to $`s\mathrm{}(2\|1)`$-invariance any eigenspace of the operator $`\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}`$ is a lowest weight $`s\mathrm{}(2\|1)`$-module generated by some lowest weight eigenvector. Therefore without loss of generality we can solve the defining relation (3.3.3) in the space of lowest weights. Let us consider in more details the structure of eigenspace of the $`s\mathrm{}(2\|1)`$-invariant operator acting on the tensor product $`V_{\mathrm{}_1,b_1}V_{\mathrm{}_2,b_2}`$. As we have seen from direct sum decomposition: $$V_{\mathrm{}_1,b_1}V_{\mathrm{}_2,b_2}=V_{\mathrm{},b}+2\underset{n=1}{\overset{\mathrm{}}{}}V_{\mathrm{}+n,b}+\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n+\frac{1}{2},b\frac{1}{2}}+\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n+\frac{1}{2},b+\frac{1}{2}}$$ for every fixed $`n`$ the space of lowest weight vectors with eigenvalue $`b`$ is two-dimensional and the ones with eigenvalues $`b\pm \frac{1}{2}`$ are one-dimensional. Therefore the operator $`\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}`$ is diagonal on odd lowest weight vectors $`\mathrm{\Psi }_n^+`$ and $`\mathrm{\Psi }_n^{}`$ but acts non-trivially on the two-dimensional subspace of even lowest weight vectors spanned on $`\mathrm{\Phi }_n^+`$ and $`\mathrm{\Phi }_n^{}`$. In matrix form we have: $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)\left(\begin{array}{c}\mathrm{\Phi }_n^+\\ \mathrm{\Phi }_n^{}\\ \mathrm{\Psi }_n^+\\ \mathrm{\Psi }_n^{}\end{array}\right)=\left(\begin{array}{cccc}A_n(u)& B_n(u)& 0& 0\\ C_n(u)& D_n(u)& 0& 0\\ 0& 0& F_n(u)& 0\\ 0& 0& 0& E_n(u)\end{array}\right)\left(\begin{array}{c}\mathrm{\Phi }_n^+\\ \mathrm{\Phi }_n^{}\\ \mathrm{\Psi }_n^+\\ \mathrm{\Psi }_n^{}\end{array}\right)$$ (3.3.4) The matrix relation (3.3.3) leads to a set of recurrence relations for the coefficients $`A_n,\mathrm{},E_n`$. Some details of calculations can be found in Appendix and here we present the final expression for the general solution of these recurrence relations: $$A_n(u)=(1)^{n+1}\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)}\frac{\text{u}+b_1b_2}{(\mathrm{}_1b_1)(\mathrm{}_2+b_2)}$$ (3.3.5) $$B_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}\frac{(\mathrm{}_1+b_1)(\mathrm{}_2b_2)}{(\mathrm{}_1b_1)(\mathrm{}_2+b_2)}$$ $$C_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)};D_n(u)=(1)^{n+1}\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}$$ $$\frac{(\mathrm{}_2b_2)(\mathrm{}_2+b_2)\left(\text{u}b_1b_2\right)\left(\text{u}+b_1+b_2\right)\left(\text{u}b_2\mathrm{}_1\right)\left(\text{u}b_2+\mathrm{}_1\right)}{(\mathrm{}_1b_1)(\mathrm{}_2+b_2)}$$ $$E_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}\frac{(\text{u}+b_1\mathrm{}_2)(\text{u}+b_1+\mathrm{}_2)}{(\mathrm{}_1b_1)(\mathrm{}_2+b_2)}$$ $$F_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}\frac{(\text{u}b_2\mathrm{}_1)(\text{u}b_2+\mathrm{}_1)}{(\mathrm{}_1b_1)(\mathrm{}_2+b_2)}$$ where we used the notations: $$\mathrm{}_nn+\mathrm{}_1+\mathrm{}_2;\text{u}\frac{u}{\eta }+b_1b_2.$$ As usual the obtained general solution of the Yang-Baxter equation is fixed up to overall normalization. We choose the normalization such that the R-matrix coincides with the permutation operator for $`u=0`$ and $`\stackrel{}{\mathrm{}}_1=\stackrel{}{\mathrm{}}_2`$. The obtained R-matrix (3.3.4) acts on the space of two-point functions which are polynomials in $`z_i,\theta _i,\overline{\theta }_i,i=1,2`$. This holds also if the representation parameters $`\mathrm{}_i,b_i,i=1,2`$ correspond to chiral or antichiral cases. If one or both modules in the tensor product are chiral or antichiral then the tensor product representation space is a proper subspace of the space of all two-point polynomials (compare (2.4.1)). It is important to observe that in these cases the action of R-matrix can be consistently restricted to the corresponding subspace. Indeed, after multiplying with the overall factor $`(\mathrm{}_1b_1)(\mathrm{}_2+b_2)`$, the matrix becomes triangular in these cases in the way as expected. We list the reduced R-matrices in all special cases involving chiral representations in Appendix B. ## 4 Homogeneous periodic chain ### 4.1 Commuting transfer matrices Let us construct the set of commuting $`s\mathrm{}(2\|1)`$-invariant operators the generating function of which is the transfer-matrix $`\text{𝕋}_\stackrel{}{m}(u)`$. We construct $`\text{𝕋}_\stackrel{}{m}(u)`$ as the supertrace of a monodromy matrix built of the elementary -matrix blocks . We introduce the $`N`$ spaces $`V_\stackrel{}{\mathrm{}}_i`$ and $`N`$ operators $`\text{}_{\stackrel{}{m},\stackrel{}{\mathrm{}}_i}(u)`$: $$\text{}_{\stackrel{}{m},\stackrel{}{\mathrm{}}_i}(u):V_\stackrel{}{m}V_\stackrel{}{\mathrm{}}_iV_\stackrel{}{m}V_\stackrel{}{\mathrm{}}_i$$ The periodicity convention $`N+11`$ is implied. The monodromy matrix $$\text{}_\stackrel{}{m}(u)\text{}_{\stackrel{}{m},\stackrel{}{\mathrm{}}_1}(uc_1)\mathrm{}\text{}_{\stackrel{}{m},\stackrel{}{\mathrm{}}_N}(uc_N)$$ (4.1.1) acts then on the space $`V_\stackrel{}{m}V_\stackrel{}{\mathrm{}}_1\mathrm{}V_\stackrel{}{\mathrm{}}_N`$, and $`\text{𝕋}_\stackrel{}{m}(u)`$ is obtained by taking the supertrace in the auxiliary space $`V_\stackrel{}{m}`$: $$\text{𝕋}_\stackrel{}{m}(u)=\text{str}_{V_\stackrel{}{m}}\text{}_\stackrel{}{m}(u).$$ These monodromy matrices form the commutative family: $$\text{𝕋}_{\stackrel{}{m}_1}(u)\text{𝕋}_{\stackrel{}{m}_2}(v)=\text{𝕋}_{\stackrel{}{m}_2}(v)\text{𝕋}_{\stackrel{}{m}_1}(u)$$ (4.1.2) The relation (4.1.2) follows from the fact that there exists the operator $`\text{}_{\stackrel{}{m}_1,\stackrel{}{m}_2}`$ such that $$\text{}_{\stackrel{}{m}_1,\stackrel{}{m}_2}(uv)\text{}_{\stackrel{}{m}_1,\stackrel{}{\mathrm{}}_i}(u)\text{}_{\stackrel{}{m}_2,\stackrel{}{\mathrm{}}_i}(v)=\text{}_{\stackrel{}{m}_2,\stackrel{}{\mathrm{}}_i}(v)\text{}_{\stackrel{}{m}_1,\stackrel{}{\mathrm{}}_i}(u)\text{}_{\stackrel{}{m}_1,\stackrel{}{m}_2}(uv)$$ The -operator is even. Therefore by using standard arguments one derives an analogous equation for the monodromy matrices: $$\text{}_{\stackrel{}{m}_1,\stackrel{}{m}_2}(uv)\text{}_{\stackrel{}{m}_1}(u)\text{}_{\stackrel{}{m}_2}(v)=\text{}_{\stackrel{}{m}_2}(v)\text{}_{\stackrel{}{m}_1}(u)\text{}_{\stackrel{}{m}_1,\stackrel{}{m}_2}(uv)$$ (4.1.3) From this one can derive easily that corresponding traces and supertraces are commuting operators separately: $$\text{𝕋}_\stackrel{}{m}^+(u)=\text{tr}_{V_\stackrel{}{m}}\text{}_\stackrel{}{m}(u);\text{𝕋}_{\stackrel{}{m}_1}^+(u)\text{𝕋}_{\stackrel{}{m}_2}^+(v)=\text{𝕋}_{\stackrel{}{m}_2}^+(v)\text{𝕋}_{\stackrel{}{m}_1}^+(u)$$ $$\text{𝕋}_\stackrel{}{m}^{}(u)=\text{str}_{V_\stackrel{}{m}}\text{}_\stackrel{}{m}(u);\text{𝕋}_{\stackrel{}{m}_1}^{}(u)\text{𝕋}_{\stackrel{}{m}_2}^{}(v)=\text{𝕋}_{\stackrel{}{m}_2}^{}(v)\text{𝕋}_{\stackrel{}{m}_1}^{}(u)$$ but only $`\text{𝕋}_\stackrel{}{m}(u)\text{𝕋}_\stackrel{}{m}^{}(u)`$ is the generating function for the $`s\mathrm{}(2\|1)`$-invariant operators. Instead of giving the general proof we demonstrate all this on the example of $`\text{𝕋}_\stackrel{}{f}(u)`$, where the auxiliary space corresponds to the fundamental representation. Let us represent $`\text{𝕋}_\stackrel{}{f}(u)`$ in the form: $$\text{𝕋}_\stackrel{}{m}(u)=e_{AB}\text{𝕋}_{AB}(u)$$ where operators $`\text{𝕋}_{AB}(u)`$ act in tensor product $`V_\stackrel{}{\mathrm{}}_1\mathrm{}V_\stackrel{}{\mathrm{}}_n`$ and we assume the summation over repeated indices. The general equation (4.1.3) has the form in this case $$\left[uv+\eta ()^{\overline{G}}e_{FG}^ie_{GF}^j\right]e_{AB}^i\text{𝕋}_{AB}(u)e_{CD}^j\text{𝕋}_{CD}(v)=$$ $$=e_{CD}^j\text{𝕋}_{CD}(v)e_{AB}^i\text{𝕋}_{AB}(u)\left[uv+\eta ()^{\overline{G}}e_{FG}^ie_{GF}^j\right].$$ The traces and supertraces of generators $`e_{CD}`$ are calculated as follows: $$\text{tr}e_{AB}(e_{AB})_{CC}=\delta _{AB};\text{str}e_{AB}()^{\overline{C}}(e_{AB})_{CC}=()^{\overline{A}}\delta _{AB}$$ Using $`e_{AB}e_{CD}=\delta _{CB}e_{AD}`$ and taking (super-)traces in corresponding spaces one easily obtains: $$\text{tr}:\text{𝕋}_{AA}(u)\text{𝕋}_{CC}(v)=\text{𝕋}_{CC}(v)\text{𝕋}_{AA}(u),$$ $$\text{str}:()^{\overline{A}}\text{𝕋}_{AA}(u)()^{\overline{C}}\text{𝕋}_{CC}(v)=()^{\overline{C}}\text{𝕋}_{CC}(v)()^{\overline{A}}\text{𝕋}_{AA}(u),$$ The $`g\mathrm{}(2\|1)`$-invariance can be demonstrated in the simplest example $`N=2`$. The generalization to arbitrary $`n`$ is straightforward. $$\text{𝕋}_\stackrel{}{f}^+(u)=\text{tr}\left(u+\eta ()^{\overline{B}}e_{AB}E_{BA}^1\right)\left(u+\eta ()^{\overline{D}}e_{CD}E_{DC}^2\right)=$$ $$=3u^2+\eta u()^{\overline{A}}\left(E_{AA}^1+E_{AA}^2\right)+\eta ^2E_{AB}^1E_{BA}^2$$ $$\text{𝕋}_\stackrel{}{f}^{}(u)=\text{str}\left(u+\eta ()^{\overline{B}}e_{AB}E_{BA}^1\right)\left(u+\eta ()^{\overline{D}}e_{CD}E_{DC}^2\right)=$$ $$=u^2+\eta u\left(E_{AA}^1+E_{AA}^2\right)+\eta ^2()^{\overline{B}}E_{AB}^1E_{BA}^2.$$ We have seen ((2.1.4) and discussion before) that only operators entering $`\text{𝕋}_\stackrel{}{f}^{}(u)`$ are $`g\mathrm{}(2\|1)`$-invariant. ### 4.2 Local Hamiltonians If one fixes the arbitrary representation $`\stackrel{}{m}=\stackrel{}{\mathrm{}}`$ and the same representations $`\stackrel{}{\mathrm{}}_1=\stackrel{}{\mathrm{}}_2=\mathrm{}=\stackrel{}{\mathrm{}}_N=\stackrel{}{\mathrm{}}`$ in remaining spaces $`V_\stackrel{}{\mathrm{}}_i`$ we obtain the generating functions of Hamiltonians. The local Hamiltonians can be obtained for the homogeneous ($`c_i=0`$) chain.The whole construction is quite general . Let us calculate the first two coefficients in the Taylor expansion of the operator $$\text{𝕋}_{\stackrel{}{\mathrm{}}}(u)=\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}(u)\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_N}(u)$$ around the point $`u=0`$. It is easy to see from the derived expression for the R-matrix that by condition $`\stackrel{}{\mathrm{}}_i=\stackrel{}{\mathrm{}}_j=\stackrel{}{\mathrm{}}`$ the point $`u=0`$ is regular: $$\text{}_{\stackrel{}{\mathrm{}}_i,\stackrel{}{\mathrm{}}_j}(0)=\text{}_{\stackrel{}{\mathrm{}}_i,\stackrel{}{\mathrm{}}_j}$$ (4.2.1) where is permutation operator. Indeed, the permutation operator acts as follows on the lowest weight basis: $$\text{}\left(\begin{array}{c}\mathrm{\Phi }_n^+\\ \mathrm{\Phi }_n^{}\\ \mathrm{\Psi }_n^+\\ \mathrm{\Psi }_n^{}\end{array}\right)=(1)^n\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right)\left(\begin{array}{c}\mathrm{\Phi }_n^+\\ \mathrm{\Phi }_n^{}\\ \mathrm{\Psi }_n^+\\ \mathrm{\Psi }_n^{}\end{array}\right)$$ The expression for matrix coefficients of R-operator takes the simple form in homogeneous case ($`\mathrm{}_1=\mathrm{}_2,b_1=b_2`$): $$A_n(u)=(1)^{n+1}\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)}\frac{\text{u}}{(\mathrm{}+b)(\mathrm{}b)}$$ (4.2.2) $$B_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)};C_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)}$$ $$D_n(u)=(1)^{n+1}\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}\frac{\text{u}\left(\text{u}^22\mathrm{}^22b^2\right)}{(\mathrm{}+b)(\mathrm{}b)}$$ $$E_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}\frac{(\text{u}+b\mathrm{})(\text{u}+b+\mathrm{})}{(\mathrm{}+b)(\mathrm{}b)}$$ $$F_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}\frac{(\text{u}b\mathrm{})(\text{u}b+\mathrm{})}{(\mathrm{}+b)(\mathrm{}b)},$$ where $$\mathrm{}_1=\mathrm{}_2=\mathrm{};b_1=b_2=b;\text{u}\frac{u}{\eta };\mathrm{}_nn+2\mathrm{}.$$ The equality (4.2.1) can be easily checked. The first coefficient in the Taylor expansion of the operator $`\text{𝕋}_{\stackrel{}{\mathrm{}}}(u)`$ is proportional to the operator of cyclic shift: $$\text{𝕋}_{\stackrel{}{\mathrm{}}}(0)=\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_2}\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_N}=const\text{}_{\stackrel{}{\mathrm{}}_N,\stackrel{}{\mathrm{}}_{N1}}\text{}_{\stackrel{}{\mathrm{}}_{N1},\stackrel{}{\mathrm{}}_{N2}}\mathrm{}\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}$$ This is readily checked. First we move $`\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}`$ to the right, then $`\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}`$ to the right an so on and after all we use $`\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}=const`$. The second coefficient has the form: $$\text{𝕋}_{\stackrel{}{\mathrm{}}}^{}(0)=\underset{i}{}\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_i}^{}(0)\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_N}$$ In order to simplify this expression let us consider the i-th term and move the $`\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_i}^{}(0)`$ to the right: $$\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_{i1}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_i}^{}(0)\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_{i+1}}\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_n}=\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_{i1}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_{i+1}}\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_n}\text{}_{\stackrel{}{\mathrm{}}_{i+1},\stackrel{}{\mathrm{}}_i}^{}(0)$$ By moving first $`\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}`$ to the right, then $`\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}`$ and so on we transform the remaining term $$\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_{i1}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_{i+1}}\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_N}=const\text{}_{\stackrel{}{\mathrm{}}_{N1},\stackrel{}{\mathrm{}}_N}\mathrm{}\text{}_{\stackrel{}{\mathrm{}}_{i1},\stackrel{}{\mathrm{}}_{i+1}}\mathrm{}\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}$$ On the last stage we multiply the obtained expression by the operator $$\text{𝕋}_{\stackrel{}{\mathrm{}}}^1(0)=(const)^1\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}\mathrm{}\text{}_{\stackrel{}{\mathrm{}}_{N1},\stackrel{}{\mathrm{}}_N}$$ from the left and obtain the following expression: $$\text{𝕋}_{\stackrel{}{\mathrm{}}}^1(0)\text{𝕋}_{\stackrel{}{\mathrm{}}}^{}(0)=\underset{i}{}\text{}_{\stackrel{}{\mathrm{}}_{i+1},\stackrel{}{\mathrm{}}_i}\text{}_{\stackrel{}{\mathrm{}}_{i+1},\stackrel{}{\mathrm{}}_i}^{}(0)=\underset{i}{}\text{}_{\stackrel{}{\mathrm{}}_i,\stackrel{}{\mathrm{}}_{i+1}}^{}(0)\text{}_{\stackrel{}{\mathrm{}}_i,\stackrel{}{\mathrm{}}_{i+1}}=\underset{i}{}_{\stackrel{}{\mathrm{}}_i,\stackrel{}{\mathrm{}}_{i+1}}.$$ (4.2.3) The resulting operator can be chosen as the Hamiltonian. It commutes with the integrals of motions generated by $`\text{𝕋}_{\stackrel{}{\mathrm{}}}(u)`$ and is a sum of operators acting on two adjacent sites only. The two-particle Hamiltonians in the sum are $$_{\stackrel{}{\mathrm{}}_i,\stackrel{}{\mathrm{}}_{i+1}}=\text{}_{\stackrel{}{\mathrm{}}_i,\stackrel{}{\mathrm{}}_{i+1}}^{}(0)\text{}_{\stackrel{}{\mathrm{}}_i,\stackrel{}{\mathrm{}}_{i+1}}$$ and have the following matrix elements: $$_{\stackrel{}{\mathrm{}}_i,\stackrel{}{\mathrm{}}_{i+1}}=\eta ^1\left(\begin{array}{cccc}2\psi (\mathrm{}_n+1)& \frac{\mathrm{}_n}{(\mathrm{}b)(\mathrm{}+b)}& 0& 0\\ \frac{2}{\mathrm{}_n}\frac{\mathrm{}^2+b^2}{(\mathrm{}b)(\mathrm{}+b)}& 2\psi (\mathrm{}_n)& 0& 0\\ 0& 0& 2\psi (\mathrm{}_n+1)+\frac{2b}{(\mathrm{}b)(\mathrm{}+b)}& 0\\ 0& 0& 0& 2\psi (\mathrm{}_n+1)\frac{2b}{(\mathrm{}b)(\mathrm{}+b)}\end{array}\right).$$ (4.2.4) The eigenvalues of this matrix and corresponding eigenvalues can be easily calculated. ## 5 Inhomogeneous open chain ### 5.1 Integrals of motion In this section we shall consider the open spin chain . To start with we introduce the operator $`_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}(u)`$: $$_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}(u)\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}(u)\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}^1(u).$$ It is possible to derive the following commutation relation for $`(u)`$: $$\text{}_{\stackrel{}{m}_1,\stackrel{}{m}_2}(uv)_{\stackrel{}{m}_1,\stackrel{}{\mathrm{}}_i}(u)\text{}_{\stackrel{}{m}_1,\stackrel{}{m}_2}(u+v)_{\stackrel{}{m}_2,\stackrel{}{\mathrm{}}_i}(v)=_{\stackrel{}{m}_2,\stackrel{}{\mathrm{}}_i}(v)\text{}_{\stackrel{}{m}_1,\stackrel{}{m}_2}(u+v)_{\stackrel{}{m}_1,\stackrel{}{\mathrm{}}_i}(u)\text{}_{\stackrel{}{m}_1,\stackrel{}{m}_2}(uv).$$ It is evident that the operator $`_\stackrel{}{m}(u)`$ constructed from the monodromy matrix (4.1.1) $$_\stackrel{}{m}(u)\text{}_\stackrel{}{m}(u)\text{}_\stackrel{}{m}^1(u),$$ satisfies the analogous equation: $$\text{}_{\stackrel{}{m}_1,\stackrel{}{m}_2}(uv)_{\stackrel{}{m}_1}(u)\text{}_{\stackrel{}{m}_1,\stackrel{}{m}_2}(u+v)_{\stackrel{}{m}_2}(v)=_{\stackrel{}{m}_2}(v)\text{}_{\stackrel{}{m}_1,\stackrel{}{m}_2}(u+v)_{\stackrel{}{m}_1}(u)\text{}_{\stackrel{}{m}_1,\stackrel{}{m}_2}(uv).$$ Using this equation it is possible to show that corresponding supertraces are commuting operators: $$𝒯_{\stackrel{}{m}_1}(u)𝒯_{\stackrel{}{m}_2}(v)=𝒯_{\stackrel{}{m}_2}(v)𝒯_{\stackrel{}{m}_1}(u);𝒯_\stackrel{}{m}(u)=\text{str}_{V_\stackrel{}{m}}_\stackrel{}{m}(u)=\text{str}_{V_\stackrel{}{m}}\text{}_\stackrel{}{m}(u)\text{}_\stackrel{}{m}^1(u).$$ If one fixes the representation $`\stackrel{}{m}=\stackrel{}{f}`$ we obtain the generating function of integrals of motions for the open chain. ### 5.2 Local Hamiltonian Let us suppose that representations $`\stackrel{}{\mathrm{}}_i`$ and shifts $`c_i`$ in the product: $$\text{}_{\stackrel{}{\mathrm{}}}(u)\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}(u+c_1)\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_2}(u+c_2)\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_N}(u+c_N)$$ are fixed as follows: $$\stackrel{}{\mathrm{}}_2=\stackrel{}{\mathrm{}}_3=\mathrm{}\stackrel{}{\mathrm{}}_{N1}=\stackrel{}{\mathrm{}};c_2=c_3=\mathrm{}c_{N1}=0\text{}_{\stackrel{}{\mathrm{}}_i,\stackrel{}{\mathrm{}}}(0)=\text{}_{\stackrel{}{\mathrm{}}_i,\stackrel{}{\mathrm{}}},i=2,\mathrm{},n1,$$ where is simply the permutation. In the case of $`s\mathrm{}(2)`$ there exist integrable nearest neighbour interactions for this slightly inhomogeneous chain . We extend this result to the case of $`s\mathrm{}(2\|1)`$ and construct the corresponding Hamiltonians. The R-matrix (3.3.4) obeys the equation: $$\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}(u)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}(u)=P(u)1$$ (5.2.1) where right side of equation is proportional of unit operator and $`P(u)`$ is the function: $$P(u)\frac{(\text{u}+b_1\mathrm{}_2)(\text{u}+b_1+\mathrm{}_2)(\text{u}b_2+\mathrm{}_1)(\text{u}b_2\mathrm{}_1)}{(\mathrm{}_1b_1)(\mathrm{}_1+b_1)(\mathrm{}_2b_2)(\mathrm{}_2+b_2)},$$ which follows from the expression for matrix elements of the R-matrix (3.3.5). The explicit form of the matrix $`\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}(u)`$ is the following: $$\text{}_{\stackrel{}{\mathrm{}}_2\stackrel{}{\mathrm{}}_1}(u)=\left(\begin{array}{cccc}\overline{D}_n(u)& \overline{C}_n(u)& 0& 0\\ \overline{B}_n(u)& \overline{A}_n(u)& 0& 0\\ 0& 0& \overline{F}_n(u)& 0\\ 0& 0& 0& \overline{E}_n(u)\end{array}\right)$$ where the matrix elements $`\overline{A}_n,\overline{B}_n,\mathrm{}`$ are obtained from the elements $`A_n,B_n,\mathrm{}`$ by formal change of variables $`\mathrm{}_1\mathrm{}_2`$ and $`b_1b_2`$. Due to equation (5.2.1) we have $$\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}^1(u)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}(u)$$ (5.2.2) so that we can use the operator $`\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}(u)`$ instead of operator $`\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}^1(u)`$. This changes the normalization of the operator $`𝒯_\stackrel{}{m}(u)`$ only. Let us calculate the first two coefficients in the Taylor expansion of the operator $$𝒯_{\stackrel{}{\mathrm{}}}(u)=\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}(u+c_1)\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_N}(u+c_N)\text{}_{\stackrel{}{\mathrm{}}_N,\stackrel{}{\mathrm{}}}(uc_N)\mathrm{}\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(uc_1).$$ The first coefficient is proportional to the unit operator for arbitrary representations $`\stackrel{}{\mathrm{}}_i`$ and shifts $`c_i`$ due to the property (5.2.2), $$𝒯_{\stackrel{}{\mathrm{}}}(0)=\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}(c_1)\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_N}(c_N)\text{}_{\stackrel{}{\mathrm{}}_N,\stackrel{}{\mathrm{}}}(c_N)\mathrm{}\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(c_1)\text{str}_V_{\stackrel{}{\mathrm{}}}1.$$ Let us introduce the short-hand notation $$_{\stackrel{}{\mathrm{}},\stackrel{}{m}}(c)=\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{m}}^{}(c)\text{}_{\stackrel{}{m},\stackrel{}{\mathrm{}}}(c)+\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{m}}(c)\text{}_{\stackrel{}{m},\stackrel{}{\mathrm{}}}^{}(c)$$ (5.2.3) and calculate the expression for $`𝒯_{\stackrel{}{\mathrm{}}}^{}(0)`$ which contains several terms , $$𝒯_{\stackrel{}{\mathrm{}}}^{}(0)=\text{str}_V_{\stackrel{}{\mathrm{}}}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}(c_1)+\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}(c_1)_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_2}(0)\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(c_1)+$$ $$+\underset{i=2\mathrm{}N}{}\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}(c_1)\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_{i1}}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_i}(c_i)\text{}_{\stackrel{}{\mathrm{}}_{i1},\stackrel{}{\mathrm{}}}\mathrm{}\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(c_1).$$ Note that this formula is true for $`c_2=\mathrm{}c_{N1}=0`$ only. Let us consider each term separately. The first term is constant $$\text{str}_V_{\stackrel{}{\mathrm{}}}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_i}(c_i)=const$$ The $`i`$-th term in the sum can be transformed easily to the simpler expression $$\text{tr}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}(c_1)\mathrm{}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_{i1}}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_i}(c_i)\text{}_{\stackrel{}{\mathrm{}}_{i1},\stackrel{}{\mathrm{}}}\mathrm{}\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(c_1)=_{\stackrel{}{\mathrm{}}_{i1},\stackrel{}{\mathrm{}}_i}(c_i)\text{str}_V_{\stackrel{}{\mathrm{}}}1$$ It turns out that the second term can be transformed to the analogous form too (remember $`c_2=0`$) $$\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}(c_1)_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_2}(0)\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(c_1)=_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}(c_1)\text{str}_V_{\stackrel{}{\mathrm{}}}1+const$$ (5.2.4) For the proof we start from the Yang-Baxter equation $$\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}(u)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}}(u+v)\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(v)=\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(v)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}}(u+v)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}(u),$$ differentiate this equation with respect to $`u`$ and then put $`v=u`$: $$\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}^{}(u)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(u)+\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}(u)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}}^{}(0)\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(u)=$$ $$=\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(u)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}}^{}(0)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}(u)+\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(u)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}^{}(u),$$ then multiply both sides of the obtained equation by the permutation $`\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}}`$ from the left ($`\stackrel{}{\mathrm{}}_2=\stackrel{}{\mathrm{}}`$): $$\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}^{}(u)\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(u)+\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}(u)\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_2}^{}(0)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(u)=$$ $$=\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}(u)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}}^{}(0)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}(u)+\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}(u)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}^{}(u),$$ and calculate $`\text{str}_V_{\stackrel{}{\mathrm{}}}`$ using the equalities: $$\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}^{}(u)\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(u)=const;\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}}^{}(0)=const.$$ After all we obtain $$\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}(u)\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_2}^{}(0)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(u)=\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}(u)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}^{}(u)\text{str}_V_{\stackrel{}{\mathrm{}}}1+const$$ and using the evident identity (consequence of eq. (5.2.1)): $$\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}(u)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}^{}(u)=\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}^{}(u)\text{}_{\stackrel{}{\mathrm{}}_2,\stackrel{}{\mathrm{}}_1}(u)+const$$ we arrive at the general formula: $$\text{str}_V_{\stackrel{}{\mathrm{}}}\text{}_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_1}(u)_{\stackrel{}{\mathrm{}},\stackrel{}{\mathrm{}}_2}(0)\text{}_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}}(u)=_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}(u)\text{str}_V_{\stackrel{}{\mathrm{}}}1+const$$ Finally we obtain the following representation for $`𝒯_{\stackrel{}{\mathrm{}}}^{}(0)`$: $$𝒯_{\stackrel{}{\mathrm{}}}^1(0)𝒯_{\stackrel{}{\mathrm{}}}^{}(0)=_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}(c_1)+\underset{i=3}{\overset{N1}{}}_{\stackrel{}{\mathrm{}}_{i1},\stackrel{}{\mathrm{}}_i}(0)+_{\stackrel{}{\mathrm{}}_{N1},\stackrel{}{\mathrm{}}_N}(c_N)+const$$ This operator, commuting with all integrals of motions, is a sum of two-particle operators and can be considered as the Hamiltonian. The two-particle Hamiltonians entering the sum can be easily calculated from the universal R-matrix by (5.2.3). The expression for the $`_{\stackrel{}{\mathrm{}}_{i1},\stackrel{}{\mathrm{}}_i}(0)`$ coincides with the two-particle Hamiltonian for the periodic chain (4.2.4) up to overall coefficient two but the expression for the $`_{\stackrel{}{\mathrm{}}_1,\stackrel{}{\mathrm{}}_2}(c)`$ is rather lengthy to be presented here. ## 6 Conclusions In this paper we have obtained the general solution of the Yang-Baxter equation acting on the tensor product of arbitrary representations of the superalgebra $`s\mathrm{}(2\|1)`$. We have represented the lowest weight module of $`s\mathrm{}(2\|1)`$ by polynomials in one even ($`z`$) and two odd ($`\theta ,\overline{\theta }`$) variables. Therefore the general R-matrices is an operator acting on two-point functions being polynomials in the two sets of variables ($`z_1,\theta _1,\overline{\theta }_1`$) and ($`z_2,\theta _2,\overline{\theta }_2`$). Instead of calculating this operator explicitly we have obtained its matrix elements on the space of lowest weights. The eigenfunctions and eigenvalues can be easily calculated from the obtained matrix. From the general R-matrix for two isomorphic representations we have calculated the nearest neighbour interaction Hamiltonians for an homogeneous closed chain. The result applies both for the finite and infinite-dimensional representations on the sites. Using the general R-matrix an integrable open chain can be constructed with arbitrary isomorphic representations on the inner sites and other arbitrary representations at the end points. The nearest neighbour interaction Hamiltonian has been calculated. ## 7 Acknowledgments We thank L.N.Lipatov and A.Manashov for the stimulating discussion and critical remarks. This work has been supported by Deutsche Forschungsgemeinschaft, grant No KI 623/1-2 and by INTAS,grant No 96-524. One of us (D.K.) is grateful to Saxonian Ministry of Science and Arts for supporting his visit at Leipzig University. ## 8 Appendix A In this Appendix we discuss briefly the derivation of the expression (3.3.5) for the general R-matrix. In matrix form the defining relation for the R-matrix reads as follows: $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)K_{AB}=\overline{K}_{AB}\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u);A,B=1,2,3.$$ (8.0.5) where: $$K=\left(\begin{array}{ccc}K_{11}& K_{12}& K_{13}\\ K_{21}& K_{22}& K_{23}\\ K_{31}& K_{32}& K_{33}\end{array}\right)\frac{u}{2\eta }(\text{}_2\text{}_1)+\text{}_1\text{}_2;\overline{K}\frac{u}{2\eta }(\text{}_2\text{}_1)+\text{}_2\text{}_1.$$ The matrix element $`\overline{K}_{AB}`$ can be obtained from $`K_{AB}`$ by formal substitution $`12`$ and $`uu`$: $$\mathrm{}_1,b_1,Z_1\mathrm{}_2,b_2,Z_2;uu,$$ where $`Z(z,\theta ,\overline{\theta })`$. The operators $`K_{AB}`$ and lowest weights transform as follows: $$K_{AB}\overline{K}_{AB};\mathrm{\Phi }_n^\pm (1)^n\mathrm{\Phi }_n^{};\mathrm{\Psi }_n^\pm (1)^{n+1}\mathrm{\Psi }_n^\pm .$$ There are nine equations and we start from the simplest one. ### 8.1 Equation $`\text{}K_{13}=\overline{K}_{13}\text{}`$ The operator $`K_{13}`$ commutes with the lowering generators $`V^{},W^{},S^{}`$ and the covariant derivatives $`D_1^\pm `$. Therefore operator $`K_{13}`$ maps lowest weight vectors with definite chirality to lowest weight vector with the same chirality and decrease its power by one: $$K_{13}\mathrm{\Phi }_n^\pm =\alpha _n^\pm (\text{u})\mathrm{\Phi }_{n1}^\pm ;K_{13}\mathrm{\Psi }_n^\pm =\beta _n^\pm (\text{u})\mathrm{\Psi }_{n1}^\pm $$ The explicit calculation gives: $$\alpha _n^+(\text{u})=n(\text{u}+\mathrm{}_n);\alpha _n^{}(\text{u})=n(\text{u}+\mathrm{}_n1);\beta _n^\pm (\text{u})=n(\text{u}+\mathrm{}_n),$$ where: $$\mathrm{}_nn+\mathrm{}_1+\mathrm{}_2;\text{u}\frac{u}{\eta }+b_1b_2.$$ The action of operator $`\overline{K}_{13}`$ on lowest weights vectors can be obtained from formulae for $`K_{13}`$ by formal substitution $`\mathrm{}_1,b_1,Z_1\mathrm{}_2,b_2,Z_2`$ and $`uu`$: $$\overline{K}_{13}\mathrm{\Phi }_n^\pm =\alpha _n^{}(\text{u})\mathrm{\Phi }_{n1}^\pm ;K_{13}\mathrm{\Psi }_n^\pm =\beta _n^\pm (\text{u})\mathrm{\Psi }_{n1}^\pm .$$ We project the operator equation $`\text{}K_{13}=\overline{K}_{13}\text{}`$ onto the lowest weight vectors $`\mathrm{\Phi }_n^\pm `$: $$\text{}K_{13}\mathrm{\Phi }_n^+=\overline{K}_{13}\text{}\mathrm{\Phi }_n^+$$ $$\alpha _n^+(\text{u})[A_{n1}\mathrm{\Phi }_{n1}^++B_{n1}\mathrm{\Phi }_{n1}^{}]=A_n\alpha _n^{}(\text{u})\mathrm{\Phi }_{n1}^+B_n\alpha _n^+(\text{u})\mathrm{\Phi }_{n1}^{}$$ $$\text{}K_{13}\mathrm{\Phi }_n^{}=\overline{K}_{13}\text{}\mathrm{\Phi }_n^{}$$ $$\alpha _n^{}(\text{u})[C_{n1}\mathrm{\Phi }_{n1}^++D_{n1}\mathrm{\Phi }_{n1}^{}]=C_n\alpha _n^{}(\text{u})\mathrm{\Phi }_{n1}^+D_n\alpha _n^+(\text{u})\mathrm{\Phi }_{n1}^{}$$ which results in the recurrent relations: $$\alpha _n^+(\text{u})A_{n1}=\alpha _n^{}(\text{u})A_n;\alpha _n^+(\text{u})B_{n1}=\alpha _n^+(\text{u})B_n$$ $$\alpha _n^{}(\text{u})C_{n1}=\alpha _n^{}(\text{u})C_n;\alpha _n^{}(\text{u})D_{n1}=\alpha _n^+(\text{u})D_n$$ with the following general solution: $$A_n(u)=A(1)^{n+1}\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)};B_n(u)=B(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}$$ $$C_n(u)=C(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)};D_n(u)=D(1)^{n+1}\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}$$ The projection onto the odd lowest weight vectors $`\mathrm{\Psi }_n^\pm `$ leads to analogous recurrent relations $$\beta _n(\text{u})F_{n1}=\beta _n(\text{u})F_n;\beta _n(\text{u})E_{n1}=\beta _n(\text{u})E_n$$ with general solution: $$E_n(u)=E(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)};F_n(u)=F(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}.$$ We see that equation $`\text{}K_{13}=\overline{K}_{13}\text{}`$ fixes the n-dependence of the matrix elements of -matrix. The remaining equations fix the coefficients $`A,B,\mathrm{}`$ up to overall normalization. ### 8.2 Equation $`\text{}K_{12}=\overline{K}_{12}\text{}`$ Due to commutativity of the operator $`K_{12}`$ with lowering generators $`W^{},S^{}`$ we can write down the general formulae for the action of operator $`K_{12}`$ on even lowest weight vectors: $$K_{12}\mathrm{\Phi }_n^\pm =a^\pm W^+\mathrm{\Phi }_{n1}^++b^\pm W^+\mathrm{\Phi }_{n1}^{}+c^\pm \mathrm{\Psi }_{n1}^{}.$$ The coefficients $`a^\pm `$ and $`b^\pm `$ can be calculated using the following commutation relation: $$\{V^{},K_{12}\}=K_{13}.$$ Remembering the known results for operator $`K_{13}`$ $$\{V^{},K_{12}\}=K_{13}V^{}K_{12}\mathrm{\Phi }_n^\pm =K_{13}\mathrm{\Phi }_n^\pm =\alpha _n^\pm (\text{u})\mathrm{\Phi }_{n1}^\pm $$ and the simple formula ($`b=b_1+b_2`$): $$V^{}W^+\mathrm{\Phi }_n^\pm =(\mathrm{}_n+b)\mathrm{\Phi }_n^\pm V^{}K_{12}\mathrm{\Phi }_n^\pm =a^\pm (\mathrm{}_{n1}+b)\mathrm{\Phi }_{n1}^+b^\pm (\mathrm{}_{n1}+b)\mathrm{\Phi }_{n1}^{},$$ we obtain: $$(\mathrm{}_{n1}+b)a^+=\alpha _n^+(\text{u}),b^+=0;a^{}=0,(\mathrm{}_{n1}+b)b^{}=\alpha _n^{}(\text{u}).$$ The coefficient $`c^+`$ can be calculated using the commutativity of $`K_{12}`$ and $`D_1^+`$: $$\{D_1^+,K_{12}\}=0D_1^+K_{12}\mathrm{\Phi }_n^+=a^+D_1^+W^+\mathrm{\Phi }_{n1}^++c^+D_1^+\mathrm{\Psi }_{n1}^{}=0$$ and the simple formulae: $$D_1^+W^+\mathrm{\Phi }_{n1}^+=(\mathrm{}_1+b_1)\mathrm{\Phi }_{n1}^+;D_1^+\mathrm{\Psi }_{n1}^{}=\mathrm{\Phi }_{n1}^+.$$ We obtain $$(\mathrm{}_1+b_1)a^+=c^+.$$ The coefficient $`c^{}`$ can be calculated using the commutation relation with the covariant derivative $`D_2^+`$: $$\{D_2^+,K_{12}\}=D_2^+W_1^{}+(\mathrm{}_2+b_2)S_1^{}$$ $$D_2^+K_{12}\mathrm{\Phi }_n^{}=b^{}D_2^+W^+\mathrm{\Phi }_{n1}^{}+c^{}D_2^+\mathrm{\Psi }_{n1}^{}=(\mathrm{}_2+b_2)S_1^{}\mathrm{\Phi }_n^{}$$ and the formulae: $$D_2^+W^+\mathrm{\Phi }_{n1}^{}=(\mathrm{}_2+b_2)\mathrm{\Phi }_{n1}^{};D_2^+\mathrm{\Psi }_{n1}^{}=\mathrm{\Phi }_{n1}^+;S_1^{}\mathrm{\Phi }_n^{}=n\mathrm{\Phi }_{n1}^{}.$$ We obtain $$(\mathrm{}_2+b_2)b^{}+c^{}=n(\mathrm{}_2+b_2).$$ Finally we have ($`b=b_1+b_2`$): $$K_{12}\mathrm{\Phi }_n^+=\frac{\alpha _n^+(\text{u})}{\mathrm{}_{n1}+b}W^+\mathrm{\Phi }_{n1}^+\frac{\alpha _n^+(\text{u})(\mathrm{}_1+b_1)}{\mathrm{}_{n1}+b}\mathrm{\Psi }_{n1}^{}$$ $$K_{12}\mathrm{\Phi }_n^{}=\frac{\alpha _n^{}(\text{u})}{\mathrm{}_{n1}+b}W^+\mathrm{\Phi }_{n1}^{}+\frac{n(\text{u}b_1b_2)(\mathrm{}_2+b_2)}{\mathrm{}_{n1}+b}\mathrm{\Psi }_{n1}^{}.$$ The analogous calculations for odd lowest weight vectors give: $$K_{12}\mathrm{\Psi }_n^+=\frac{\beta _n(\text{u})}{\mathrm{}_n+b}W^+\mathrm{\Psi }_{n1}^++\frac{(\text{u}b)(\mathrm{}_2+b_2)}{\mathrm{}_n+b}\mathrm{\Phi }_n^++\frac{\beta _n(\text{u})(\mathrm{}_1+b_1)}{n(\mathrm{}_n+b)}\mathrm{\Phi }_n^{};K_{12}\mathrm{\Psi }_n^{}=\frac{\beta _n(\text{u})}{\mathrm{}_{n1}+b}W^+\mathrm{\Psi }_{n1}^{}$$ The expression for the action of the operator $`\overline{K}_{12}`$ is obtained by symmetry: $$\overline{K}_{12}\mathrm{\Phi }_n^{}=\frac{\alpha _n^+(\text{u})}{\mathrm{}_{n1}+b}W^+\mathrm{\Phi }_{n1}^{}\frac{\alpha _n^+(\text{u})(\mathrm{}_2+b_2)}{\mathrm{}_{n1}+b}\mathrm{\Psi }_{n1}^{}$$ $$\overline{K}_{12}\mathrm{\Phi }_n^+=\frac{\alpha _n^{}(\text{u})}{\mathrm{}_{n1}+b}W^+\mathrm{\Phi }_{n1}^+\frac{n(\text{u}+b)(\mathrm{}_1+b_1)}{\mathrm{}_{n1}+b}\mathrm{\Psi }_{n1}^{}$$ $$\overline{K}_{12}\mathrm{\Psi }_n^+=\frac{\beta _n(\text{u})}{\mathrm{}_n+b}W^+\mathrm{\Psi }_{n1}^++\frac{(\text{u}+b)(\mathrm{}_1+b_1)}{\mathrm{}_n+b}\mathrm{\Phi }_n^{}\frac{\beta _n(\text{u})(\mathrm{}_2+b_2)}{n(\mathrm{}_n+b)}\mathrm{\Phi }_n^{};\overline{K}_{12}\mathrm{\Psi }_n^{}=\frac{\beta _n(\text{u})}{\mathrm{}_{n1}+b}W^+\mathrm{\Psi }_{n1}^{}.$$ The projection of the operator equation $`\text{}K_{12}=\overline{K}_{12}\text{}`$ onto lowest weight vectors leads to the following new recurrent relations ($`b=b_1+b_2`$): $$\text{}K_{12}\mathrm{\Phi }_n^+=\overline{K}_{12}\text{}\mathrm{\Phi }_n^+$$ (8.2.1) $$(\mathrm{}_1+b_1)\alpha _n^+(\text{u})E_{n1}=n(\mathrm{}_1+b_1)(\text{u}+b)A_n+(\mathrm{}_2+b_2)\alpha _n^+(\text{u})B_n$$ $$\text{}K_{12}\mathrm{\Phi }_n^{}=\overline{K}_{12}\text{}\mathrm{\Phi }_n^{}$$ $$n(\mathrm{}_2+b_2)(\text{u}b)E_{n1}=n(\mathrm{}_1+b_1)(\text{u}+b_1+b_2)C_n+(\mathrm{}_2+b_2)\alpha _n^+(\text{u})D_n$$ $$\text{}K_{13}\mathrm{\Phi }_n^+=\overline{K}_{13}\text{}\mathrm{\Phi }_n^+$$ $$n(\mathrm{}_2+b_2)(\text{u}b)A_n+(\mathrm{}_1+b_1)\beta _n(\text{u})C_n=(\mathrm{}_2+b_2)\beta _n(\text{u})F_n$$ $$n(\mathrm{}_2+b_2)(\text{u}b)B_n+(\mathrm{}_1+b_1)\beta _n(\text{u})D_n=n(\mathrm{}_1+b_1)(\text{u}+b_1+b_2)F_n$$ ### 8.3 Equation $`\text{}K_{23}=\overline{K}_{23}\text{}`$ There exists an automorphism of the algebra $`s\mathrm{}(2\|1)`$: $$W^\pm V^\pm ;BB$$ which has the following form in our representation: $$\theta \overline{\theta };bb.$$ Due to this automorphism the matrix $`K_{AB}`$ has definite symmetry properties with respect to the transformation: $$\mathrm{}_1,z_1\mathrm{}_2,z_2;\theta _1,b_1\overline{\theta }_2,b_2$$ The operators $`K_{12}`$,$`K_{23}`$ and lowest weights transform as follows: $$K_{12}K_{23};\mathrm{\Phi }_n^\pm (1)^n\mathrm{\Phi }_n^\pm ;\mathrm{\Psi }_n^\pm (1)^{n+1}\mathrm{\Psi }_n^{}.$$ This symmetry allows to use the results of previous section and write down the formulae for the action of the operator $`K_{23}`$ on lowest weight vectors ($`b=b_1+b_2`$): $$K_{23}\mathrm{\Phi }_n^+=\frac{\alpha _n^+(\text{u})}{\mathrm{}_{n1}b}V^+\mathrm{\Phi }_{n1}^+\frac{\alpha _n^+(\text{u})(\mathrm{}_2b_2)}{\mathrm{}_{n1}b}\mathrm{\Psi }_{n1}^+,$$ $$K_{23}\mathrm{\Phi }_n^{}=\frac{\alpha _n^{}(\text{u})}{\mathrm{}_{n1}b}V^+\mathrm{\Phi }_{n1}^{}+\frac{n(\text{u}+b)(\mathrm{}_1b_1)}{\mathrm{}_{n1}b}\mathrm{\Psi }_{n1}^+,$$ $$K_{23}\mathrm{\Psi }_n^{}=\frac{\beta _n(\text{u})}{\mathrm{}_nb}V^+\mathrm{\Psi }_{n1}^{}\frac{(\text{u}+b)(\mathrm{}_1b_1)}{\mathrm{}_nb}\mathrm{\Phi }_n^+\frac{\beta _n(\text{u})(\mathrm{}_2b_2)}{n(\mathrm{}_nb)}\mathrm{\Phi }_n^{};K_{23}\mathrm{\Psi }_n^+=\frac{\beta _n(\text{u})}{\mathrm{}_{n1}b}V^+\mathrm{\Psi }_{n1}^+,$$ $$\overline{K}_{23}\mathrm{\Phi }_n^{}=\frac{\alpha _n^+(\text{u})}{\mathrm{}_{n1}b}V^+\mathrm{\Phi }_{n1}^{}\frac{\alpha _n^+(\text{u})(\mathrm{}_1b_1)}{\mathrm{}_{n1}b}\mathrm{\Psi }_{n1}^+,$$ $$\overline{K}_{23}\mathrm{\Phi }_n^+=\frac{\alpha _n^{}(\text{u})}{\mathrm{}_{n1}b}V^+\mathrm{\Phi }_{n1}^+\frac{n(\text{u}b)(\mathrm{}_2b_2)}{\mathrm{}_{n1}b}\mathrm{\Psi }_{n1}^+,$$ $$\overline{K}_{23}\mathrm{\Psi }_n^{}=\frac{\beta _n(\text{u})}{\mathrm{}_nb}V^+\mathrm{\Psi }_{n1}^{}\frac{(\text{u}b)(\mathrm{}_2b_2)}{\mathrm{}_nb}\mathrm{\Phi }_n^{}+\frac{\beta _n(\text{u})(\mathrm{}_1b_1)}{n(\mathrm{}_nb)}\mathrm{\Phi }_n^+;\overline{K}_{23}\mathrm{\Psi }_n^+=\frac{\beta _n(\text{u})}{\mathrm{}_{n1}b}V^+\mathrm{\Psi }_{n1}^+.$$ The projection of the operator equation $`\text{}K_{23}=\overline{K}_{23}\text{}`$ onto lowest weight vectors leads to the following new recurrent relations ($`b=b_1+b_2`$): $$\text{}K_{23}\mathrm{\Phi }_n^+=\overline{K}_{23}\text{}\mathrm{\Phi }_n^+$$ (8.3.1) $$(\mathrm{}_2b_2)\alpha _n^+(\text{u})F_{n1}=n(\mathrm{}_2b_2)(\text{u}b)A_n+(\mathrm{}_1b_1)\alpha _n^+(\text{u})B_n$$ $$\text{}K_{23}\mathrm{\Phi }_n^{}=\overline{K}_{23}\text{}\mathrm{\Phi }_n^{}$$ $$n(\mathrm{}_1b_1)(\text{u}+b)F_{n1}=n(\mathrm{}_2b_2)(\text{u}b)C_n+(\mathrm{}_1b_1)\alpha _n^+(\text{u})D_n$$ $$\text{}K_{23}\mathrm{\Phi }_n^+=\overline{K}_{23}\text{}\mathrm{\Phi }_n^+$$ $$n(\mathrm{}_1b_1)(\text{u}+b)A_n+(\mathrm{}_2b_2)\beta _n(\text{u})C_n=(\mathrm{}_1b_1)\beta _n(\text{u})E_n$$ $$n(\mathrm{}_1b_1)(\text{u}+b)B_n+(\mathrm{}_2b_2)\beta _n(\text{u})D_n=n(\mathrm{}_2b_2)(\text{u}b)E_n$$ The systems of equations (8.2.1), (8.2.1) fix the coefficients $`A,B,C\mathrm{}`$ up two arbitrary constants. The next operator equation $`\text{}K_{11}=\overline{K}_{11}\text{}`$ fixes the remaining ambiguity but we avoid presenting the rather lengthy formulae here. Alternatively the missing equation can be obtained as follows. The even lowest weight vectors $`\mathrm{\Phi }_n^\pm `$ coincide for $`n=0`$. Therefore one obtains: $$\text{}\mathrm{\Phi }_0^+=\text{}\mathrm{\Phi }_0^{}A(\text{u}+\mathrm{}_0)+B=\frac{C(\text{u}+\mathrm{}_0)}{\text{u}+\mathrm{}_0}\frac{D}{\text{u}+\mathrm{}_0}$$ Finally the systems (8.2.1), (8.3.1) and this last equation fix the solution completely: $$A=\text{u}+b_1b_2;B=(\mathrm{}_1+b_1)(\mathrm{}_2b_2);C=(\mathrm{}_1b_1)(\mathrm{}_2+b_2)$$ $$D=(\mathrm{}_2b_2)(\mathrm{}_2+b_2)\left(\text{u}b_1b_2\right)\left(\text{u}+b_1+b_2\right)\left(\text{u}b_2\mathrm{}_1\right)\left(\text{u}b_2+\mathrm{}_1\right)$$ $$E=(\text{u}+b_1\mathrm{}_2)(\text{u}+b_1+\mathrm{}_2);F=(\text{u}b_2\mathrm{}_1)(\text{u}b_2+\mathrm{}_1).$$ We have checked that obtained -matrix really is the solution also of the remaining seven equations but the involved formulae become rather lengthy starting from the equation $`\text{}K_{11}=\overline{K}_{11}\text{}`$. ## 9 Appendix B The R-matrix acting in the tensor product of chiral modules can be obtained by simple reduction from the general R-matrix (3.3.4). First of all we have to change the overall normalization multiplying all matrix elements by the factor $`(\mathrm{}_1b_1)(\mathrm{}_2+b_2)`$. Let us consider all possible special cases: 1. chiral at 1, generic at 2 $$V_{\mathrm{}_1,\mathrm{}_1}V_{\mathrm{}_2,b_2}=\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n,b}+\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n+\frac{1}{2},b\frac{1}{2}};\mathrm{}_2\pm b_2;\mathrm{\Phi }_n^+V_{\mathrm{}+n,b};\mathrm{\Psi }_n^+V_{\mathrm{}+n+\frac{1}{2},b\frac{1}{2}}$$ $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)\left(\begin{array}{c}\mathrm{\Phi }_n^+\\ \mathrm{\Psi }_n^+\end{array}\right)=R\left(\begin{array}{cc}A_n(u)& 0\\ 0& F_n(u)\end{array}\right)\left(\begin{array}{c}\mathrm{\Phi }_n^+\\ \mathrm{\Psi }_n^+\end{array}\right)$$ $$A_n(u)=(1)^{n+1}\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)};F_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}(\text{u}+\mathrm{}_1b_2)$$ 2. antichiral at 1, generic at 2 $$V_{\mathrm{}_1,\mathrm{}_1}V_{\mathrm{}_2,b_2}=\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n,b}+\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n+\frac{1}{2},b+\frac{1}{2}};\mathrm{}_2\pm b_2;\mathrm{\Phi }_n^{}V_{\mathrm{}+n,b};\mathrm{\Psi }_n^{}V_{\mathrm{}+n+\frac{1}{2},b+\frac{1}{2}}$$ $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)\left(\begin{array}{c}\mathrm{\Phi }_n^{}\\ \mathrm{\Psi }_n^{}\end{array}\right)=R\left(\begin{array}{cc}D_n(u)& 0\\ 0& E_n(u)\end{array}\right)\left(\begin{array}{c}\mathrm{\Phi }_n^{}\\ \mathrm{\Psi }_n^{}\end{array}\right)$$ $$D_n(u)=(1)^{n+1}\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}(\text{u}\mathrm{}_1b_2);E_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}$$ 3. chiral at 1, antichiral at 2 $$V_{\mathrm{}_1,\mathrm{}_1}V_{\mathrm{}_2,\mathrm{}_2}=\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n,b};\mathrm{\Phi }_n^+V_{\mathrm{}+n,b}$$ $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)\mathrm{\Phi }_n^+=RA_n(u)\mathrm{\Phi }_n^+;A_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)}$$ 4. antichiral at 1, chiral at 2 $$V_{\mathrm{}_1,\mathrm{}_1}V_{\mathrm{}_2,\mathrm{}_2}=\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n,b};\mathrm{\Phi }_n^{}V_{\mathrm{}+n,b}$$ $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)\mathrm{\Phi }_n^{}=RD_n(u)\mathrm{\Phi }_n^{};D_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}$$ 5. antichiral at 1 and 2 $$V_{\mathrm{}_1,\mathrm{}_1}V_{\mathrm{}_2,\mathrm{}_2}=\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n+\frac{1}{2},b+\frac{1}{2}};\mathrm{\Psi }_n^{}V_{\mathrm{}+n+\frac{1}{2},b+\frac{1}{2}}$$ $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)\mathrm{\Psi }_n^{}=RE_n(u)\mathrm{\Psi }_n^{};E_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}$$ 6. chiral at 1 and 2 $$V_{\mathrm{}_1,\mathrm{}_1}V_{\mathrm{}_2,\mathrm{}_2}=\underset{n=0}{\overset{\mathrm{}}{}}V_{\mathrm{}+n+\frac{1}{2},b\frac{1}{2}};\mathrm{\Psi }_n^+V_{\mathrm{}+n+\frac{1}{2},b\frac{1}{2}}$$ $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)\mathrm{\Psi }_n^+=RF_n(u)\mathrm{\Psi }_n^+;F_n(u)=(1)^n\frac{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}{\mathrm{\Gamma }\left(\text{u}+\mathrm{}_n+1\right)}$$ Note that we do not fix the overall normalization of the obtained R-matrix. ## 10 Appendix C In this Appendix we discuss shortly the case of finite-dimensional representations and show that the obtained general R-matrix reduces to the known ones for the tensor product of modules with minimal dimensions. The tensor product of two finite-dimensional $`s\mathrm{}(2\|1)`$-modules has the following direct sum decomposition : $$V_{\mathrm{}_1,b_1}V_{\mathrm{}_2,b_2}=V_{\mathrm{},b}+V_{\mathrm{}+N,b}+2\underset{n=1}{\overset{N1}{}}V_{\mathrm{}+n,b}+\underset{n=0}{\overset{N1}{}}V_{\mathrm{}+n+\frac{1}{2},b\frac{1}{2}}+\underset{n=0}{\overset{N1}{}}V_{\mathrm{}+n+\frac{1}{2},b+\frac{1}{2}};\mathrm{}_i\pm b_i$$ (10.0.2) $$\mathrm{}_1=\frac{n_1}{2},n_12;\mathrm{}_2=\frac{n_2}{2},n_22;Nmin(n_1,n_2);\mathrm{}=\frac{n_1+n_2}{2},b=b_1+b_2$$ Note that this formula is applicable in the generic situation $`\mathrm{}_i\pm b_i,\mathrm{}_i1/2`$. The direct sum decomposition reduces for the module $`V_{\frac{1}{2},b}`$: $$V_{\mathrm{}_1,b_1}V_{\frac{1}{2},b_2}=V_{\mathrm{},b}+V_{\mathrm{}+1,b}+V_{\mathrm{}+\frac{1}{2},b\frac{1}{2}}+V_{\mathrm{}+\frac{1}{2},b+\frac{1}{2}};\mathrm{}=\frac{n_1+1}{2},b=b_1+b_2;n_12,\mathrm{}_2=\frac{1}{2}$$ $$V_{\frac{1}{2},b_1}V_{\frac{1}{2},b_2}=V_{1,b}+V_{\frac{1}{2},b\frac{1}{2}}+V_{\frac{1}{2},b+\frac{1}{2}};b=b_1+b_2,\mathrm{}_1=\mathrm{}_2=\frac{1}{2}.$$ The origin of modifications for the tensor product involving the $`V_{\frac{1}{2},b}`$-module is very simple. The module $`V_{\frac{1}{2},b}`$ is the four-dimensional vector space with the following basis: $$\mathrm{\Phi }_0=1;\mathrm{\Phi }_1=z+b\theta \overline{\theta };\mathrm{\Psi }_0^+=\overline{\theta },\mathrm{\Psi }_0^{}=\theta .$$ It is evident that we are able to construct the following two-point lowest weight vectors only: $$V_{\mathrm{}_1,b_1}V_{\frac{1}{2},b_2}\mathrm{\Phi }_0^+=\mathrm{\Phi }_0^{}=1;(b_2\frac{1}{2})\mathrm{\Phi }_0^+(b_2+\frac{1}{2})\mathrm{\Phi }_0^{};\mathrm{\Psi }_0^+=\overline{\theta }_{12};\mathrm{\Psi }_0^{}=\theta _{12}.$$ Note that $$(b_2\frac{1}{2})\mathrm{\Phi }_0^+(b_2+\frac{1}{2})\mathrm{\Phi }_0^{}=z_1+b_2\theta _1\overline{\theta }_1+z_2+b_2\theta _2\overline{\theta }_2+(b_2\frac{1}{2})\overline{\theta }_1\theta _2(b_2+\frac{1}{2})\theta _1\overline{\theta }_2,$$ $$V_{\frac{1}{2},b_1}V_{\frac{1}{2},b_2}\mathrm{\Phi }_0^+=\mathrm{\Phi }_0^{}=1;\mathrm{\Psi }_0^+=\overline{\theta }_{12};\mathrm{\Psi }_0^{}=\theta _{12}.$$ Let us consider the general R-matrix (3.3.4) acting in the tensor product $`V_{\frac{1}{2},b_1}V_{\frac{1}{2},b_2}`$. In this case we have: $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)=(A_0(u)+B_0(u))\text{}_1+E_0(u)\text{}_2+F_0(u)\text{}_3$$ where $`\text{}_i`$ are projectors on the modules in the direct sum decomposition: $$V_{\frac{1}{2},b_1}V_{\frac{1}{2},b_2}=V_{1,b}+V_{\frac{1}{2},b\frac{1}{2}}+V_{\frac{1}{2},b+\frac{1}{2}};b=b_1+b_2,\mathrm{}_1=\mathrm{}_2=\frac{1}{2}$$ $$\text{}_1V_{1,b};\text{}_2V_{\frac{1}{2},b\frac{1}{2}};\text{}_3V_{\frac{1}{2},b+\frac{1}{2}}.$$ After a simple calculation one obtains: $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)\text{}_1+\frac{2\text{u}1+2b_1}{2\text{u}+12b_2}\text{}_2+\frac{2\text{u}12b_2}{2\text{u}+1+2b_1}\text{}_3.$$ This result coincides with the expression for the R-matrix given in : $$\text{}(\mu )=\text{}_1+\frac{4\mu 1+b_1+b_2}{4\mu +1b_1b_2}\text{}_2+\frac{4\mu 1b_1b_2}{4\mu +1+b_1+b_2}\text{}_3$$ up to the overall normalization and a redefinition of the spectral parameter: $$\text{u}=2\mu \frac{b_1b_2}{2}.$$ The direct sum decomposition for the tensor product of the chiral modules (atypical representations) is well known . Now we need the simplest ones: $$V_{\frac{1}{2},\frac{1}{2}}V_{\frac{1}{2},b_2}=V_{1,b}+V_{\frac{1}{2},b\frac{1}{2}};b=b_1+b_2,\mathrm{}_1=\mathrm{}_2=\frac{1}{2},b_1=\frac{1}{2}$$ Let us consider the R-matrix (3.3.4) acting in the tensor product $`V_{\frac{1}{2},\frac{1}{2}}V_{\frac{1}{2},b_2}`$. In this case we have (see Appendix B): $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)=A_0(u)\text{}_1+F_0(u)\text{}_2$$ where $`\text{}_i`$ are projectors on the modules in the direct sum decomposition: $$\text{}_1V_{1,b};\text{}_2V_{\frac{1}{2},b\frac{1}{2}}.$$ After a simple calculation one obtains: $$\text{}_{\stackrel{}{\mathrm{}}_1\stackrel{}{\mathrm{}}_2}(u)\text{}_1+\frac{2\text{u}12b_2}{2\text{u}+2}\text{}_2.$$ This result also coincides with the expression for such a -matrix given in .
warning/0003/astro-ph0003444.html	ar5iv	text	# 1 Status of Observational Cosmology ## 1 Status of Observational Cosmology The ten most significant parameters to be determined in observational cosmology are: 1. Age of the universe: $`t_0`$ 2. The Hubble constant $`H_0`$ from the Hubble relation: $`v=H_0d`$ 3. Density parameter : $`\mathrm{\Omega }_m=\frac{\rho _m}{3H_0^2/8\pi G}`$ 4. Deceleration parameter: $`q_0=\frac{\ddot{R}(t_0)R(t_0)}{\dot{R}^2(t_0)}`$ 5. The baryon density $`\mathrm{\Omega }_B`$ and the vacuum density $`\mathrm{\Omega }_\mathrm{\Lambda }`$ 6. The parameters associated with microwave backgound fluctuations: $`n,\sigma _8,T/S,N_T`$ The strongest lower limit for $`t_0`$ is determined from studies of the stellar populations of globular clusters. The main error in the globular clusters age estimate comes from the uncertain distance to the globular clusters. A 0.25 magnitude error in the distance translates into a $`22\%`$ error in the cluster age . Independent age limits come from the cooling of white dwarfs. The best estimates give $`t_013`$ Gyr, with a lower limit of $`11`$ Gyr. For $`t_0>13`$ Gyr, we have $`h0.50`$ (where $`h`$ is defined by the Hubble parameter: $`H_0=100`$ h km $`s^1\mathrm{Mpc}^1`$) for matter density $`\mathrm{\Omega }_m=1`$, and for $`h0.73`$ we have $`\mathrm{\Omega }_m0.3`$ in spatially flat cosmologies with $`\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }=1`$. The Hubble parameter is now better determined (it used to be known to within a factor of two). Most measurements are now consistent with a value: $`h=0.65\pm 0.08`$ . It is remarkable that data obtained from several different methods for determing $`H_0`$ lead to similar results, which gives hope for an ultimate convergence of measurements. For $`\mathrm{\Omega }_m=0.4`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.6`$ and $`h=0.65\pm 0.08`$, the age of the universe would be $`t_0=13\pm 2`$ Gyr, in agreement with globular cluster age estimates. This result is one of the strong arguments for a low matter density $`\mathrm{\Omega }_m0.3`$ and a non-zero cosmological constant $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$. A most promising new way of measuring $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ on cosmological scales is to use small-angle anisotropies in the CMB radiation and high-redshift Type Ia supernovae (SNe Ia). The Supernovae Cosmology Project (CSP) and the High-Z Supernovae team have found a significant number of Type Ia supernovae . The more recent larger SCP data set of 42 high redshift data gives for the flat case $`\mathrm{\Omega }_m=0.28_{0.080.04}^{+0.09+0.05}`$ . The High-Z Supernovae group has also measured $`\mathrm{\Omega }_m`$ giving in the flat case $`\mathrm{\Omega }_m=0.4\pm 0.3`$. Two possible sources of problems are the dimming by dust and the assumption made that evolution for nearby and far supernovae is uniform. In CMB anisotropy studies, the location of the first acoustic Doppler peak at angular wave number $`l250`$ is a strong indication of a flat universe $`\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }=1`$. The MAXIMA and BOOMERANG balloon flights seem to confirm this result, and the existence of a second and possible third peak would appear to be consistent with the predictions of simple inflation models. New data from the NASDA Microwave Anisotropy Probe satellite will hopefully strengthen these results. We can summarise the main observational results: 1. Age of universe$`t_0=916\mathrm{Gyr}`$ (from globular clusters) $`=917\mathrm{Gyr}^{}`$ 2. Hubble parameter$`H_0=100hs^1Mpc^1,h=0.65\pm 0.08`$. 3. Baryon density$`\mathrm{\Omega }_bh^2=0.019\pm 0.001`$ (from $`D/H`$) $`>0.015`$ from Ly$`\alpha `$ forest opacity\* 4. Matter density$`\mathrm{\Omega }_m=0.4\pm 0.2`$ (from cluster baryons) $`=0.34\pm 0.1`$ from Ly$`\alpha `$ forest $`P(k)^{}(P(k)=Ak^n\mathrm{with}n=1`$ for the Harrison-Zel’dovich spectrum) $`=0.4\pm 0.2`$ from cluster evolution\* $`>\frac{3}{4}\mathrm{\Omega }_\mathrm{\Lambda }\frac{1}{4}\pm \frac{1}{8}`$ from SN Ia $`>0.3(2.4\sigma `$ from flows) 5. Total density$`\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }=1\pm 0.3`$ (from CMB peak location) 6. Dark vacuum energy density$`\mathrm{\Omega }_\mathrm{\Lambda }=0.8\pm 0.3`$ (from last two lines) 7. Neutrino density$`\mathrm{\Omega }_\nu 0.001`$ (from Superkamiokande) $`0.1^{}`$ Here, the cosmological parameters with * are obtained by assuming $`\mathrm{\Lambda }CDM`$ i.e cold dark matter models with non-zero cosmological constant. The Type Ia reshift measurements have indicated the remarkable result that the cosmic expansion is presently undergoing an acceleration. ## 2 Bimetric Gravity Theory A new kind of vector-tensor and scalar-tensor theory of gravity, which exhibits a bimetric structure and contains two or more light cones \[6, Clayton+Moffat:1999a, 7\], has been introduced, recently. This type of model has attracted some attention , and similar effects have been noted elsewhere . The motivation for considering these models is derived form earlier work , which provided a scenario in which some of the outstanding issues in cosmology can be resolved. These models provide a fundamental dynamical mechanism for varying speed of light theories and generate a new mechanism for an inflationary epoch that could solve the initial value problems of early universe cosmology. In the following, we shall review some of the main features of this new kind of gravity theory and its application to cosmology. In these models matter that satisfies the strong energy condition can nevertheless contribute to the cosmic acceleration. Our cosmological model can be mapped to a model with varying fundamental constants , albeit not uniquely and requiring some care in the interpretation of the varying constants that appear. It is hoped that the models can shed some light on the new observational data that suggests the expansion of the universe at present is undergoing an acceleration . Although there has been some success in understanding the latter problem by the inclusion of a class of very particular scalar field potentials , it is fair to say that not all issues have been resolved. Using the scalar field version of the model, we expect that not only will we be able to generate sufficient inflation, but that a quintessence-like solution should be achieveable. We shall be considering models wth an action of the form $$S=\overline{S}_{\mathrm{gr}}[\overline{g}]+S[g,\psi ]+\widehat{S}[\widehat{g},\widehat{\varphi }^I].$$ (1) The first term is the usual Einstein-Hilbert action for general relativity constructed from a metric $`\overline{g}_{\mu \nu }`$, and the final term is the contribution from the non-gravitational (matter) fields in spacetime $`\widehat{\varphi }^I`$, and is built from a different but related metric $`\widehat{g}_{\mu \nu }`$. The contribution $`S[g,\psi ]`$ is constructed from a metric $`g_{\mu \nu }`$ and includes kinetic terms for a field or fields (unspecified as yet) $`\psi `$ that may be considered to be part of the gravitational sector, modifying the reaction of spacetime to the presence of the matter fields in $`\widehat{S}[\widehat{g},\widehat{\varphi }^I]`$. The manner in which $`\psi `$ accomplishes this is by modifying the metric that appears in each of the actions. For example, in $`\psi `$ was a vector field, $`\overline{g}_{\mu \nu }=g_{\mu \nu }`$ and $`\widehat{g}_{\mu \nu }=g_{\mu \nu }+b\psi _\mu \psi _\nu `$, whereas in \[Clayton+Moffat:1999a\] $`\psi `$ was a scalar field, $`\overline{g}_{\mu \nu }=g_{\mu \nu }`$ and $`\widehat{g}_{\mu \nu }=g_{\mu \nu }+b_\mu \psi _\nu \psi `$. These relations imply that matter and gravitational fields propagate at different velocities if $`\psi `$ is non-vanishing. Since the matter action $`\widehat{S}`$ is built using only $`\widehat{g}_{\mu \nu }`$, it is the null surfaces of $`\widehat{g}_{\mu \nu }`$ along which matter fields propagate. If we assume that other than the presence of a “composite” metric the matter action is otherwise a conventional form (perfect fluid, scalar field, Maxwell, etc.), then variation of the matter action yields the matter energy-momentum tensor $`\widehat{T}^{\mu \nu }`$, which will be conserved $$\widehat{}_\nu \widehat{T}^{\mu \nu }=0,$$ (2) by virtue of the matter field equations $`\widehat{F}_I=0`$. Throughout we will write, for example, $`\widehat{}_\nu `$ for the covariant derivative constructed from the Levi-Civita connection of $`\widehat{g}_{\mu \nu }`$. Since we also assume that the matter fields satisfy the dominant energy condition, we therefore know (assuming appropriate smoothness of $`\widehat{g}_{\mu \nu }`$) that matter fields cannot travel faster than the speed of light as determined by $`\widehat{g}_{\mu \nu }`$. The gravitational action is written $$\overline{S}_{\mathrm{gr}}[\overline{g}]=\frac{1}{\kappa }𝑑\overline{\mu }\overline{R},$$ (3) where we use a metric with ($`+`$$``$$``$$``$) signature and have defined $`\kappa =16\pi G/c^4`$. We will denote the metric densities by, e.g., $`\overline{\mu }=\sqrt{det(\overline{g}_{\mu \nu })}`$ and in addition write $`d\overline{\mu }=\overline{\mu }dtd^3x`$. We will not consider a cosmological constant, since it can easily be included later. We can identify the metric $`\overline{g}_{\mu \nu }`$ as providing the light cone for the gravitational system. We consider a Proca model with arbitrary potential $$S[g,\psi ]=\frac{1}{\kappa }𝑑\mu \left(\frac{1}{4}B^2V(X)\right),$$ (4) where we will use the definition $$X=\frac{1}{2}\psi ^2,$$ (5) and $`V^{}(X)=V(X)/X`$. We will also use $`B_{\mu \nu }=_\mu \psi _\nu _\nu \psi _\mu `$, $`\psi ^2=g^{\mu \nu }\psi _\mu \psi _\nu `$ and $`B^2=g^{\alpha \mu }g^{\beta \nu }B_{\alpha \beta }B_{\mu \nu }`$. We will assume that as $`\psi _\mu 0`$ we have $`V(X)m^2X`$ and therefore the linearized (in $`\psi _\mu `$) limit of our model is identical to Einstein-Proca field equations coupled to matter. The standard energy-momentum tensor for the vector field is $$T^{\mu \nu }=B^{\mu \alpha }B_{}^{\nu }{}_{\alpha }{}^{}+\frac{1}{4}g^{\mu \nu }B^2+V^{}\psi ^\mu \psi ^\nu Vg^{\mu \nu }.$$ (6) Although there exists a more general class of models, we will limit ourselves to the choice $$\widehat{g}_{\mu \nu }=g_{\mu \nu }+b\psi _\mu \psi _\nu ,\overline{g}_{\mu \nu }=g_{\mu \nu }+g\psi _\mu \psi _\nu ,$$ (7) where $`b`$ and $`g`$ are constants, so that the variations of $`\widehat{g}_{\mu \nu }`$ and $`\overline{g}_{\mu \nu }`$ are related to those of $`g_{\mu \nu }`$ and $`\psi _\mu `$. The field equations are given by $$_\mu B^{\mu \nu }+V^{}\psi ^\nu +gT^{\mu \nu }\psi _\mu +\kappa \frac{\widehat{\mu }}{\mu }(gb)\widehat{T}^{\nu \mu }\psi _\mu =0,$$ (8) $$\overline{\mu }\overline{G}^{\mu \nu }=\frac{1}{2}\mu T^{\mu \nu }+\frac{1}{2}\kappa \widehat{\mu }\widehat{T}^{\mu \nu }.$$ (9) It is clear that $`\widehat{g}_{\mu \nu }`$ and $`\overline{g}_{\mu \nu }`$ provide the characteristic surfaces for matter and gravitational fields, respectively. We can prove that any matter model that conserves energy-momentum with respect to $`\widehat{g}_{\mu \nu }`$ is consistent with the gravitational structure that we have introduced . The “most physical” metric is clearly $`\widehat{g}_{\mu \nu }`$, since it describes the geometry on which matter propagates and interacts. Because all matter fields are coupled to the same metric $`\widehat{g}_{\mu \nu }`$ in exactly the same way, the weak equivalence principle is satisfied. Furthermore, because one can work in a local Lorentz frame of $`\widehat{g}_{\mu \nu }`$, in which non-gravitational physics takes on its special relativistic form, the Einstein equivalence principle is also satisfied. However, because $`\widehat{g}_{\mu \nu }`$ does not couple to matter in the same way as in general relativity unless $`\psi _\mu =0`$, the strong equivalence principle will be violated. The main motivation for considering these theories is that they should have something to say about the horizon problem in the early universe. If $`\psi _\mu 0`$, then if we choose $`b>g`$, matter fields will propagate outside the light cone of the gravitational field. As $`\psi _\mu 0`$ the matter light cone will ‘contract’ and matter and gravitational disturbances will eventually propagate at the same velocity. If one considers a frame in which gravitational waves propagate at a constant speed, then as the light cone of matter contracts, the universe will appear to material observers to expand acausally. ## 3 Cosmology Implicit in the idea of a varying light speed is that the speed of light is changing with respect to some fixed frame of reference. If one introduces a fundamental frame for this, then it is perhaps sensible to introduce a function $`c(t,x)`$ to describe this variability . The models that we are considering are based on the idea that the speed of light can be changing with respect to the speed of gravitational disturbances, and therefore any indication of the speed of light as a function of spacetime is frame-dependent. In particular, we will see that a frame in which the speed of light is constant and the speed of gravitational disturbances is changing is connected via a diffeomorphism to a frame where the speed of gravitational disturbances is constant, and the speed of light is changing. Quantities of interest such as the local light cone, horizons, etc. are derived directly from the relevant metric, thereby avoiding any guesswork as to which ‘speed of light’ to use—the gravitational or electromagnetic . The constant $`c`$ is fixed in the present universe by making measurements of the electromagnetic field. In a homogeneous and isotropic (FRW) universe, the vector field $`\psi _\mu `$ has components $`\psi _\mu =(c\psi _0(\tau ),0,0,0)`$. We will begin with the metric $`g_{\mu \nu }`$ in comoving form $$g_{\mu \nu }dx^\mu dx^\nu =c^2d\tau d\tau R^2(\tau )\gamma _{ij}dx^idx^j,$$ (10) and therefore $$\widehat{g}_{\mu \nu }dx^\mu dx^\nu =\widehat{\mathrm{\Theta }}^2(\tau )c^2d\tau d\tau R^2(\tau )\gamma _{ij}dx^idx^j,$$ $$\overline{g}_{\mu \nu }dx^\mu dx^\nu =\overline{\mathrm{\Theta }}^2(\tau )c^2d\tau d\tau R^2(\tau )\gamma _{ij}dx^idx^j.$$ (11) The spatial metric in spherical coordinates has the standard form $$\gamma _{ij}=\mathrm{diag}(1/(1\mathrm{kr}^2),\mathrm{r}^2,\mathrm{r}^2\mathrm{sin}^2\theta ),$$ (12) and we have defined $$\widehat{\mathrm{\Theta }}=\sqrt{1+2bX},\mathrm{and}\overline{\mathrm{\Theta }}=\sqrt{1+2gX},$$ (13) where from (5) we have $`X=\frac{1}{2}\psi _0^2`$. Although we begin with the choice (10), once we have derived the field equations, we will make a coordinate transformation in order to put $`\widehat{g}_{\mu \nu }`$ in comoving form and thereby make a comparison with the standard cosmological results a simpler matter. Note that we are reversing the definitions of $`t`$ and $`\tau `$ as used in our previous article . The matter energy-momentum tensor will have a perfect fluid form $$\widehat{T}^{\mu \nu }=\left(\rho +\frac{p}{c^2}\right)\widehat{u}^\mu \widehat{u}^\nu p\widehat{g}^{\mu \nu },$$ (14) where we have written the velocity field as $`\widehat{u}^\mu `$ to emphasize that it is normalized using the metric $`\widehat{g}_{\mu \nu }`$, so that $$\widehat{g}_{\mu \nu }\widehat{u}^\mu \widehat{u}^\nu =c^2.$$ (15) The matter conservation laws (2) lead to the usual relation $$_\tau \rho +3\frac{_\tau R}{R}\left(\rho +\frac{p}{c^2}\right)=0.$$ (16) The Friedmann equations take the form $$\left(\frac{_\tau R}{R}\right)^2+\frac{kc^2\overline{\mathrm{\Theta }}^2}{R^2}=\frac{\kappa c^4}{6}\overline{\mathrm{\Theta }}^3\left[\frac{\rho }{\widehat{\mathrm{\Theta }}}+\frac{1}{\kappa c^2}(2XV^{}V)\right],$$ (17) $$2\frac{_\tau ^2R}{R}+\left(\frac{_\tau R}{R}\right)^2+\frac{kc^2\overline{\mathrm{\Theta }}^2}{R^2}2\frac{_\tau R}{R}\frac{_\tau \overline{\mathrm{\Theta }}}{\overline{\mathrm{\Theta }}}=\frac{\kappa c^2}{2}\overline{\mathrm{\Theta }}\left[\widehat{\mathrm{\Theta }}p+\frac{1}{\kappa }V\right].$$ (18) The single remaining Proca field equation from (8) is $$\frac{1}{c\widehat{\mathrm{\Theta }}}\psi _0\left[\widehat{\mathrm{\Theta }}\left(\overline{\mathrm{\Theta }}^2V^{}gV\right)\kappa (bg)c^2\rho \right]=0.$$ (19) We now perform the coordinate transformation $$dt=\widehat{\mathrm{\Theta }}d\tau ,$$ (20) and defining $$\eta =\frac{\overline{\mathrm{\Theta }}}{\widehat{\mathrm{\Theta }}}=\sqrt{\frac{1+2gX}{1+2bX}},$$ (21) we see that the metric $`\widehat{g}_{\mu \nu }`$ is put into comoving form $$\widehat{g}_{\mu \nu }dx^\mu dx^\nu =c^2dtdtR^2(t)\gamma _{ij}dx^idx^j,$$ $$\overline{g}_{\mu \nu }dx^\mu dx^\nu =\eta ^2(t)c^2dtdtR^2(t)\gamma _{ij}dx^idx^j.$$ (22) We have $$\left(\frac{\dot{R}}{R}\right)^2+\frac{kc^2\eta ^2}{R^2}=\eta ^2\frac{\kappa c^4}{6}\rho _{\mathrm{eff}},$$ (23) $$2\frac{\ddot{R}}{R}+\left(\frac{\dot{R}}{R}\right)^2+\frac{kc^2\eta ^2}{R^2}2\frac{\dot{R}}{R}\frac{\dot{\eta }}{\eta }=\eta ^2\frac{\kappa c^2}{2}p_{\mathrm{eff}},$$ (24) where $`\dot{\rho }=_t\rho `$. In (23), we have defined the effective energy and pressure densities as $$\rho _{\mathrm{eff}}=\eta \left(\rho +\frac{1}{\kappa c^2}\widehat{\mathrm{\Theta }}(2XV^{}V)\right),p_{\mathrm{eff}}=\frac{1}{\eta }\left(p+\frac{1}{\kappa \widehat{\mathrm{\Theta }}}V\right).$$ (25) The reason for making these definitions is that (23) has exactly the form of the Friedmann equations for the metric $`\overline{g}_{\mu \nu }`$, and therefore these effective energy and momentum densities will also satisfy the conservation laws (16). The function $`R(t)`$ is written in comoving coordinates and, therefore, the speed of light is constant. This emphasizes that having a ‘varying speed of light’ is a frame-dependent statement. In a frame where the speed of matter propagation (including electromagnetic fields) is constant, the speed of gravitational waves will be changing. In a frame where the speed of gravitational waves is constant, the speed of matter propagation will be changing. This, of course, is as it should be, since we have not introduced any nondynamical preferred frame into our model. In the following we will specialize to a model where the vector field potential is a simple mass term: $$V=m^2X,V^{}=m^2.$$ (26) In this case (25) becomes $$\rho _{\mathrm{eff}}=\eta \left(\rho +(bg)\rho _{pt}\widehat{\mathrm{\Theta }}X\right),$$ $$p_{\mathrm{eff}}=\frac{1}{\eta }\left(p+c^2(bg)\rho _{pt}\frac{X}{\widehat{\mathrm{\Theta }}}\right).$$ (27) The nontrivial solution ($`\psi _00`$) of the field equation (19) leads to $$\rho =\rho _{pt}\widehat{\mathrm{\Theta }}(1+gX),$$ (28) where $$\rho _{pt}=\frac{m^2}{\kappa c^2(bg)},H_{pt}=\sqrt{\frac{c^2m^2}{6(bg)}},$$ (29) are the density at which $`\psi _0^2=0`$ is reached, and the inverse Hubble time at which this occurs (assuming that $`k=0`$). We can now write the acceleration parameter as observed by material observers from (23) as $$\widehat{q}=\frac{\ddot{R}}{H^2R}=\frac{\kappa c^4}{12}\frac{\eta ^2}{H^2}\left(\rho _{\mathrm{eff}}+\frac{3}{c^2}p_{\mathrm{eff}}\right)\frac{\dot{\eta }}{H\eta },$$ (30) where we have defined the Hubble function $`H=\dot{R}/R`$. We have $$\frac{\dot{\eta }}{H\eta }=3\frac{(bg)}{\rho _{pt}\overline{\mathrm{\Theta }}^2\widehat{\mathrm{\Theta }}(g+b+3bgX)}\left(\rho +\frac{p}{c^2}\right).$$ (31) ## 4 The Very Early Universe For very short times following the initial singularity, we expect that $`\psi _0`$ is large, and if we assume that $`gX1`$ and $`bX1`$, then from (28) we find that $$\rho =\rho _{pt}\sqrt{2b}gX^{3/2}.$$ (32) This results in the Friedmann equation $$\left(\frac{\dot{R}}{R}\right)^2+\frac{k\overline{c}^2}{R^2}=\frac{\overline{\kappa }\overline{c}^4}{6}\rho ,$$ (33) where $$\overline{c}=c\sqrt{\frac{g}{b}},\overline{G}=G\sqrt{\frac{g}{b}},\overline{\kappa }=\frac{16\pi \overline{G}}{\overline{c}^4}.$$ (34) Although the behaviour of the solutions are well-known, it is worth pointing out that the ‘effective’ constants $`\overline{c}`$ and $`\overline{G}`$ are not interpretable as the effective speed of light and gravitational constant, rather they are effective constants that dictate how the gravitational field reacts to the presence of matter. Matter fields continue to propagate with speed $`c`$ consistent with (3). It is the gravitational field perturbations that propagate with speed $`\overline{c}`$, which is the justification for the notation. During this phase there is clearly no inflation, but the horizon scales of the gravitational field and matter fields are related by $$\overline{d}_H(t)=\frac{\overline{c}}{c}\widehat{d}_H(t)=\sqrt{\frac{g}{b}}\widehat{d}_H(t),\mathrm{where}\widehat{d}_H(t)=cR(t)_0^t\frac{ds}{R(s)},$$ (35) with a similar definition for $`\overline{d}_H(t)`$ using the metric $`\overline{g}_{\mu \nu }`$. Because we have $`g<b`$ we expect that not only is the speed of gravitational disturbances slower than that of matter, but also that the coupling between matter and the gravitational sector is also lessened. What we have here is very close to what was originally envisaged by one of us in . This is part of the motivation for including the $`g0`$ possibility, the other is that the approach to the initial singularity in this phase follows the same path as in ordinary GR+matter models, with a re-interpretation of the parameters. In this case we have a model that interpolates between this initial period where $`\overline{c}>c`$ and the later universe where $`\overline{c}=c`$. ## 5 Inflation and Light Cone Contraction As $`\psi _0`$ decreases towards the point where $`gX1`$ the solution will no longer be a good approximation. If we now consider the solution when $`gX1`$, from (28) we have $$\widehat{\mathrm{\Theta }}=\frac{\rho }{\rho _{pt}},\mathrm{or}X=\frac{1}{2b}\left[\left(\frac{\rho }{\rho _{pt}}\right)^21\right],$$ (36) and the Friedmann equation (23) becomes $$\left(\frac{\dot{R}}{R}\right)^2+\frac{kc^2\eta ^2}{R^2}\left(\frac{\rho _{pt}}{\rho }\right)^2=\frac{\kappa c^4}{12}\rho _{pt}\left[1+\left(\frac{\rho _{pt}}{\rho }\right)^2\right].$$ (37) In this limit $$\rho _{\mathrm{eff}}+\frac{3}{c^2}p_{\mathrm{eff}}=\frac{1}{\rho _{pt}}\left[\rho \left(\rho +\frac{3}{c^2}p\right)+\rho ^2\rho _{pt}^2\right],$$ (38) which is greater than zero if the strong energy condition is satisfied, since $`\rho \rho _{pt}`$, and (30) reduces to $$\widehat{q}=\frac{\kappa c^4\rho _{pt}}{12H^2}\left[\frac{1}{\rho }\left(\rho +\frac{3}{c^2}p\right)+1\left(\frac{\rho _{pt}}{\rho }\right)^2\right]\frac{3}{\rho }\left(\rho +\frac{p}{c^2}\right).$$ (39) Since we expect that $`H^2`$ is large in the early universe (we can arrange that $`\rho _{pt}\rho _c`$ where $`\rho _c=12H^2/(\kappa c^4)`$, it is clear from (39) that even if matter satisfies the strong energy condition, the final term will dominate and $`\widehat{q}<0`$ (unless, perhaps, the weak energy condition is also violated). This is the expansion of the universe as seen by material observers. The acceleration of the gravitational geometry $`\overline{q}`$ would lack the final term and therefore $`\overline{q}>0`$. That we get inflation was demonstrated previously , where an exact solution for $`k=0`$ and $`g=0`$ was found. Although we discovered that we could not get enough expansion to solve the horizon problem with pure radiation, a slowly rolling scalar field could provide the necessary negative pressure. The role that the extra structure of our model plays is that the fine-tuning that is required in a simple scalar field, potential-driven model is alleviated. The flatness problem requires a bit more explanation. Dividing (23) by $`H^2`$ and defining $$\overline{ϵ}=\frac{kc^2\eta ^2}{(\dot{R})^2},$$ (40) we find a differential equation that $`\overline{ϵ}`$ satisfies by taking a derivative and using (23) to give $$\dot{\overline{ϵ}}=\frac{\kappa c^4\eta ^2}{6H}\overline{ϵ}\left(\rho _{\mathrm{eff}}+\frac{3}{c^2}p_{\mathrm{eff}}\right).$$ (41) Therefore, since $`\overline{ϵ}>0`$ and $`H>0`$ in the early universe, the only way for $`\overline{ϵ}=0`$ to be an attractor for (23) is for $`\rho _{\mathrm{eff}}+\frac{3}{c^2}p_{\mathrm{eff}}<0`$ at least for part of the history of the universe. What is not so obvious is whether the quantity $`\overline{ϵ}`$ as defined in (40) is of physical relevance. The quantity of geometrical importance is the $`3`$-curvature of the spacelike slices, $`6k/R^2`$, which suggests that the physically meaningful quantity to examine would be $$\widehat{ϵ}=\frac{kc^2}{(\dot{R})^2},$$ (42) which has the equation of motion $$\dot{\widehat{ϵ}}=2\widehat{ϵ}\widehat{q}.$$ (43) Another way of stating this is that the curvature radius defines $`\mathrm{\Omega }`$ through $$R_{\mathrm{curv}}=\frac{R}{\|k\|^{1/2}}=\frac{c}{H(\|\mathrm{\Omega }1\|)^{1/2}},$$ (44) and so $`\widehat{ϵ}=\|\mathrm{\Omega }1\|`$. Since we found from (39) that $`\widehat{q}<0`$ in the early universe, clearly $`\widehat{ϵ}=0`$ is an attractor for (23), and since it is most-likely the quantity of physical importance for matter physics, we can also claim to have solved the flatness problem once the horizon problem is solved. ## 6 Conclusions In our bimetric model, the universe generically accelerates ($`\widehat{q}<0`$) during some period in the early universe, and in the same period the physical importance of spatial curvature diminishes ($`\|\mathrm{\Omega }1\|`$ is decreasing). This can occur even when the matter fields satisfy the strong energy condition. The model that we have considered generalizes that which appeared in in a way that more closely follows the scenario discussed in . In the very early universe, matter and gravitational fields propagate with different and approximately constant velocities. During a period in which the matter light cone, originally much larger than the light cone of gravity, contracts, material observers will see an acausal expansion of the universe similar to inflation. Because the light cone of gravity does not undergo the same contraction, we expect there to be an observable difference in the scalar versus tensor contributions to the cosmic microwave background anisotropies. ## 7 Acknowledgements This work was supported by the Natural Sciences and Engineering Research Council of Canada.
warning/0003/astro-ph0003328.html	ar5iv	text	# The dynamical evolution of the fragmented, bipolar dust shell around the carbon star IRC +10 216 based on observations performed with the 6 m telecope at the Special Astrophysical Observatory, Russia ## 1 Introduction IRC +10 216 (CW Leo) is the nearest and best–studied carbon star and one of the brightest infrared sources. It experiences a strong mass loss at a rate of $`\dot{M}25\times 10^5`$M$`_{}`$yr<sup>-1</sup> (see e.g. Loup et al. LoupForveilleEtAl93 (1993)). The central star of IRC +10 216 is a long–period variable star (LPV) with a period of approximately 649 days (Le Bertre LeBertre92 (1992)). Recent distance estimates of 110 pc to 135 pc (Groenewegen Groenewegen97 (1997)) and 150 pc (Crosas & Menten CrosasMenten97 (1997)) were reported. IRC +10 216’s initial mass can be expected to be close to 4 M (Guelin et al. GuelinForestiniEtAl95 (1995), Weigelt et al. WeigeltBalegaEtAl98 (1998)). The bipolar appearance of the nebula around this object was already reported by Christou et al. (ChristouRidgwayEtAl90 (1990)) and Kastner & Weintraub (KastnerWeintraub94 (1994)). The non-spherical structure is consistent with the conjecture that IRC +10 216 is in a phase immediately before entering the protoplanetary nebula stage. High–resolution observations of this object and its circumstellar dust shell were reported by McCarthy et al. (McCarthyMcLeodEtAl90 (1990)), Christou et al. (ChristouRidgwayEtAl90 (1990)), Dyck et al. (DyckBensonEtAl91 (1991)), Danchi et al. (DanchiBesterEtAl94 (1994)), Osterbart et al. (OsterbartBalegaEtAl97 (1997)), Weigelt et al. (WeigeltBalegaEtAl97 (1997), WeigeltBalegaEtAl98 (1998) WeigeltBloeckerEtAl99 (1999)), Skinner et al. (SkinnerMeixnerEtAl98 (1998)), and Haniff & Buscher (HaniffBuscher98 (1998)). The results of Dyck et al. (DyckBensonEtAl91 (1991)) and Haniff & Buscher (HaniffBuscher98 (1998)) showed that the structure of the dust shell of IRC +10 216 has been changing for some years. Detailed radiative transfer calculations for IRC +10 216 were recently performed by Ivezić & Elitzur (IvezicElitzur96 (1996)), Crosas & Menten (CrosasMenten97 (1997)), and Groenewegen (Groenewegen97 (1997)) using a large amount of spectroscopic and visibility data. The aim of this paper is to discuss the properties of the inner dust shell of IRC +10 216 on the basis of a series of high–resolution observations. In a second paper (Men’shchikov et al., in prep.) we will present a detailed two–dimensional radiative transfer model for this object. ## 2 Observational results The IRC +10 216 speckle interferograms were obtained with the 6 m telescope at the Special Astrophysical Observatory in Russia and our Nicmos 3 camera at four epochs and with our Hawaii–array camera at one epoch (Nov. 1998). Table 1 lists the observational parameters. ### 2.1 $`J`$-, $`H`$-, and $`K`$-band reconstructions Figs. 1 and 2 show the $`K`$ and $`H`$ images, respectively, of the central region of IRC +10 216 for all epochs. The high–resolution images were reconstructed from the speckle interferograms using the bispectrum speckle interferometry method (Weigelt Weigelt77 (1977), Lohmann et al. LohmannWeigeltEtAl83 (1983), Weigelt Weigelt91 (1991)). The achieved resolution of the images depends slightly on the data quality (seeing and number of recorded interferograms) and is indicated in the figure captions. The object power spectra were determined with the speckle interferometry method (Labeyrie Labeyrie70 (1970)). The speckle transfer functions were derived from speckle interferograms of the unresolved stars mentioned in Table 1. We denote the resolved components in the center of the nebula as A, B, C, and D (see Figs. 1g and 2b) in the order of decreasing peak intensity (based on the $`K`$ band results from 1996). In addition, Fig. 2b shows three fainter components denoted with E, F, and G. In Fig. 2c, cuts through the centers of components A and B of the $`H`$ and $`K`$ images from 1997 are shown to illustrate the differences in the relative intensities of A and B for different wavelengths. Fig. 3 shows images of the faint nebula around IRC +10 216, which seems to have a bipolar structure. The resolution of these images was reduced to 149 mas ($`J`$) and 95 mas ($`H`$ and $`K`$) to increase the signal–to–noise ratio in the outer parts of the nebula. The faint granular structure of the images ($`0.5\%`$ of the peak intensity) is partly caused by speckle noise. Note that the faint extended feature at position angle PA $``$340° in the $`J`$ image corresponds quite well to a very faint component (denoted with E in Figs. 1g and 2b) visible in all images in Figs. 1 and 2 (assuming that the brightest component in the $`J`$ image is roughly coinciding with component A in the $`H`$ and $`K`$ images). Two other faint features at PA $``$20° and PA $``$50° (F and G in Fig. 2b) are only visible in some of the images. ### 2.2 Polarimetry Fig. 4 shows the results of polarimetric observations with the HST Nicmos camera at a wavelength of 1.1 $`\mu `$m (raw data retrieved from the Hubble Data Archive, STScI). The data were obtained on April 30, 1997 at a photometric phase of $`\mathrm{\Phi }=0.76`$. Data for three polarization angles were available (0°, 120°, and 240°). It has been carefully checked that the three images were correctly registered. The images have been obtained without new pointing of the telescope so that no correction was necessary (see Voit Voit97 (1997)). From the data the total intensity (contours in Fig. 4a and b), the polarized intensity (Fig. 4c), the degree and the position angle of the polarization have been derived (see e.g. Voit Voit97 (1997), Fischer et al. FischerHenningEtAl94 (1994)). Fig. 4b shows the derived polarization map superimposed on the contours of the total intensity. In Fig. 4a the two images were superimposed by centering the brightest peaks onto each other. Looking at the southern tails of A and at the northern arms this centering seems to be appropriate. The structure of the 1.1 $`\mu `$m HST image at the position of B is affected by the diffraction pattern of A. Nevertheless, the contours indicate that a structure corresponding to B might be present in the 1.1 $`\mu `$m image. ### 2.3 The core structure and the faint nebula The multi–component structure of the IRC +10 216 $`K`$–band image (of the innermost 300 mas $`\times `$ 300 mas) has already been reported by Weigelt et al. (WeigeltBalegaEtAl98 (1998)) and Haniff & Buscher (HaniffBuscher98 (1998)). This bright inner region is surrounded by a larger faint nebula (with $`1\%`$ of the peak brightness of A; see Fig. 3b and c). Three arms of the nebula can be seen at position angles of $``$30° (NE), 340° (NW), and 210° (SW) with respect to component A. At 160° (counterside of the NW-arm) the nebula is much fainter. The direction from component A to component B (PA$``$20°) can be taken as the direction of the main axis (see also Kastner & Weintraub KastnerWeintraub94 (1994)). This direction fits well to the main axis of the H<sup>13</sup>CN($`J=10`$) emission (Dayal & Bieging DayalBieging95 (1995)) which is weakly elongated on a scale of about 10<sup>′′</sup>. On larger scales asymmetries are not observed (Dayal & Bieging DayalBieging95 (1995), Groenewegen et al. GroenewegenVanDerVeenEtAl97 (1997), Mauron & Huggins MauronHuggins99 (1999)). ### 2.4 The bipolar structure in $`J`$ and at shorter wavelengths The $`J`$–band image (Fig. 3a) and the 0.79 $`\mu `$m and 1.06 $`\mu `$m HST images (Haniff & Buscher HaniffBuscher98 (1998), Skinner et al. SkinnerMeixnerEtAl98 (1998), see also Fig. 4a for the total intensity of the HST 1.1 $`\mu `$m polarization data) show a bipolar shape of the nebula. The southern lobe has a cometary or fan–shaped structure whereas the northern area of the images shows two arms reminiscent to (but weaker than) the northern X–shaped structure of the Red Rectangle (see, e.g., Men’shchikov et al. MenshchikovBalegaEtAl98 (1998)). However, the fact that even in the polarized intensity the nebula is very faint on the southeastern side suggests the main axis to be at PA$``$20° to 30°. ### 2.5 $`HK`$ color image For the January 1997 data, the integral intensities in the full fields of view of our camera (5$`\stackrel{}{.}`$1 in $`H`$ and 7$`\stackrel{}{.}`$8 in $`K`$) were compared to those of the photometric standard stars HIP 71284 (BS 5447) and HD 106965 (cf. Elias et al. EliasFrogelEtAl82 (1982)). The resulting IRC +10 216 magnitudes are $`K=2.3`$ and $`H=5.4`$. According to an extrapolation of the fitted light curve of Le Bertre (LeBertre92 (1992)), IRC +10 216 was very close to its light minimum in January 1997. Our magnitudes are in good agreement with this prediction. In a square aperture of 1$`\stackrel{}{.}`$6 the magnitudes can be determined to be $`K=2.5\pm 0.1`$ and $`H=5.7\pm 0.1`$. The integral color in this field is thus $`HK3.2`$. The resolved two–dimensional $`HK`$ color image is shown in Fig. 5 together with the contours of the $`H`$–band image. The $`H`$ and $`K`$ images used for the determination of the $`HK`$ color were reconstructed with a common resolution of 95 mas which is also the resolution of the color image. The dependence of the $`HK`$ color image on relative shifts between the two images was investigated because it is not a priori known whether the positions of the intensity maxima in $`H`$ and $`K`$ coincide. The relative position of the components A and B is very similar in the $`H`$ and $`K`$ images from 1997 so that a solution could be found where both components A and B are almost coinciding (within a few mas) for the two images. We found that the color image is not very sensitive to relative shifts within the realistic uncertainties of some milli–arcseconds. ### 2.6 Separation of components A and B Fig. 6 shows the separation of the components A and B in the $`K`$–band images for different epochs. Phase 0 in Fig. 6 corresponds to JD=2449430. This date of the photometric maximum was derived from the results of Le Bertre (LeBertre92 (1992)). The separations are: 191 mas, 201 mas, 214 mas, 245 mas, and 265 mas, for the 5 epochs from 1995 to 1998 shown in Fig. 1. The linear regression fit gives a value of 23 mas/yr for the average increase in the apparent separation of the components. Interpreting this increase as a real motion would lead to 14 km/s within the plane of the sky (for a distance of $`D=130`$ pc). The apparent relative motion of the nebula components is obviously not simply related to the stellar variability which has a period of approximately 649 days (Le Bertre LeBertre92 (1992)). It may thus be related to either an overall expansion of an inhomogeneous circumstellar dust medium or a variability of the dust shell with a period significantly larger than the stellar pulsation period (cf. Fleischer et al. FleischerGaugerEtAl92 (1992), Winters et al. WintersFleischerEtAl94 (1994), WintersFleischerEtAl95 (1995)). Table 2 lists the apparent relative velocities within the plane of the sky of the components A to E with respect to either A or B or the center positions defined by A and B or by A, B, C, and D. The values were determined from linear regression fits. ### 2.7 Structural changes within the nebula. The images by Haniff & Buscher (HaniffBuscher98 (1998)) taken in 1989 and 1997, as well as the one–dimensional data collected by Dyck et al. (DyckBensonEtAl91 (1991)) covering a larger range of epochs, already show that the structure of the envelope changes on time scales of some years. Our $`K`$–band observations now allow us to study the changes within approximately 2 stellar pulsation cycles in more detail. Besides the apparent motion of the components, Fig. 1 shows that these components change their shapes and relative fluxes. The brightest component A appears to be rather symmetric in 1995. At the later epochs it approximately keeps its size perpendicular to the axis A–B ($``$20°) but becomes narrower along this axis. The peak–to–peak intensity ratio of B and A is approximately constant from 1995 to 1997. Thereafter the component B started fading. At the same time the other components become brighter and detached from A. Note that the photometric phases of IRC +10 216 in January 1997 and November 1998 are almost identical. In fact, the integral $`K`$ magnitudes were the same ($`K=2.3`$). Again we find that the time scale for the changes seen in our images is significantly different from the period of the stellar pulsation. ## 3 Discussion The observations presented in the previous section as well as other high–resolution observations (e.g. Haniff & Buscher HaniffBuscher98 (1998)) provide detailed information on the variable structure in IRC +10 216. The question of where, behind all the dust, the central star is located, is of specific interest to understand the physical properties of the nebula. At short wavelengths the nebula shows a bipolar structure. A comparison of the observed structure with other bipolar nebulae like the Red Rectangle (Men’shchikov et al. MenshchikovBalegaEtAl98 (1998)) suggests that the X–like arms originate mainly from scattering of stellar light on the surfaces of cavities. The star is then at least partially obscured by an optically thick dust shell or torus. Haniff & Buscher (HaniffBuscher98 (1998)) argued convincingly that the main axis of the object is tilted with its southern side towards the observer. ### 3.1 Is component A the star? Motion of the components. From molecular line observations the terminal velocity of the expansion of the circumstellar envelope was found to be $`14`$ to 15 km/s (Dayal & Bieging DayalBieging95 (1995), Gensheimer & Snyder GensheimerSnyder97 (1997)). The velocities of B and D (14 km/s and 13 km/s, resp.) with respect to A can, therefore, only be understood if the direction of motion is approximately within the plane of the sky. However, the structure of the bipolar nebula at short wavelengths shows that we are looking neither pole–on nor edge–on at the nebula but at an intermediate viewing angle (e.g. 50° to 60°). The assumption that A is the star is thus not satisfactorily consistent with our observation. Polarimetry. The synthetic polarization maps of Fischer et al. (FischerHenningEtAl96 (1996)) show that a significant polarization at the position of the star is only present for nearly edge–on configurations. The high degree of polarization of A ($`14\%`$) is thus not in agreement with A being the star. At larger separations from A the polarization pattern is centrosymmetric. The center of such a pattern is thought to be at the position of the illuminating source (Fischer et al. FischerHenningEtAl94 (1994), FischerHenningEtAl96 (1996)). The fact that in Fig. 4b this center does not coincide with A but is located significantly north of it, supports our interpretation that A is not the star. ### 3.2 Is the star at or near the position of B? In the following (a to d) we will argue that the assumption that the star is at or near the position of B is consistent with the observations. (a) The cometary shapes of A in the $`H`$ and the $`J`$ images and the 0.79 $`\mu `$m and 1.06 $`\mu `$m HST images (Haniff & Buscher HaniffBuscher98 (1998)), as well as the polarization data strongly suggest that A is part of a scattering lobe within a bipolar structure. Consistently, A and its southern tails are relatively blue ($`HK`$ ranging from 2 to 3.2 in Fig. 5) compared to the integral color. (b) The northern components B, C, and D are, on the other hand, rather red ($`HK4.2`$) in comparison with the integral color. This suggests that these structures are strongly obscured and reddened by circumstellar dust. (c) The brightest northern component in the $`J`$ image (Fig. 3a) can hardly be seen in the $`H`$ image and is thus very blue. It has a separation of $``$500 mas from A at a position angle of $``$27°. In the $`K`$–band image from the same epoch (Fig. 1b) component B is at a separation of $``$200 mas and at a position angle of 21°. This means that B is almost in the center between the northern X–shaped $`J`$-band component and the southern cometary component A (see Sect. 2.4). This morphology suggests that the latter components are the opposite lobes of a bipolar structure around a central star approximately located at the position of B. (d) The polarization map (Fig. 4b) fits well with the picture that the star is at B. The region between the two $`J`$–band lobes is only weakly polarized with polarization null points in the east and northwest. Such null points indicate a transition from an optically thick region with multiple scattering to a region with predominantly single scattering of photons (Piirola et al. PiirolaScaltritiEtAl92 (1994), Fischer et al. FischerHenningEtAl94 (1994), FischerHenningEtAl96 (1996)). The center of the centrosymmetric polarization pattern is located between the two lobes. Because of the scatter in the direction of the polarization vectors it is not possible to determine very precisely if the center coincides with B, but the polarization pattern is consistent with this assumption. We note that at the position of B the polarized intensity (Fig. 4c) may be slightly contaminated by the wings of the HST diffraction pattern associated with the dominant component A. The polarimetry of A itself, however, is not affected. Since the polarimetric features discussed above and in Sec. 3.1 are only slightly influenced by this contamination, the conclusions drawn so far do not change. Changes in the mass loss rate. The change of the shape of component A and the fading of B can be attributed to an increasing mass loss which is accompanied by a gradual increase of the optical depth of the dust shell. This is most obvious for the later observation epochs, suggesting an enhanced mass loss since 1997. A strongly variable mass loss has, in fact, been predicted by theoretical models treating the dust formation mechanism in the envelopes of long–period variable carbon stars (Fleischer et al. FleischerGaugerEtAl92 (1992), Winters et al. WintersFleischerEtAl94 (1994), WintersFleischerEtAl95 (1995)). Periods of this mechanism may be significantly longer than the stellar pulsation period (Winters et al. WintersFleischerEtAl95 (1995)). An increasing optical depth of the inner dust shell also results in an increasing apparent separation of the dust shell structures. The apparent motion of the components is thus not solely determined by the velocity of the dust particles but also by the changes of the optical properties of the circumstellar dust. Alternative model: the star between A and B. From the present observational data it is not possible to exclude the possibility that the star may be located between A and B, close to B, or in the center of A, B, C, and D. In particular, the precision of the polarization map is not sufficient to conclude whether the star is at the position of B or only close to it. Radiative transfer modeling. An answer to the question of where the star is located requires radiative transfer calculations. In a second paper (Men’shchikov et al., in preparation), we will present the results of our two–dimensional radiative transfer modeling showing that the shapes of A and B cannot be reproduced when assuming a position of the star between A and B. On the other hand, it was possible to reproduce these shapes and the intensity ratio of A and B in the case where the star is assumed to be at the position of B. Clear preference is thus given to the latter model. ### 3.3 Stellar evolution and bipolar structure IRC +10 216 is without doubt in a very advanced stage of its AGB evolution due to its long pulsational period, high mass-loss rate, and carbon-rich dust-shell chemistry indicating that a significant number of thermal pulses have already taken place. The star’s initial mass can be estimated to be $`4`$M$`{}_{}{}^{}\pm 1`$M due to the observed isotopic ratios of C, N and O in the dust shell (Guelin et al. GuelinForestiniEtAl95 (1995)) and the luminosity of the central star (Weigelt et al. WeigeltBalegaEtAl98 (1998)). Accordingly, the core mass should be $`0.7`$ to $`0.8`$M with corresponding thermal-pulse cycle times of $`1\mathrm{3\hspace{0.17em}10}^4`$yr (Blöcker Bloecker95 (1995)). Introducing the mean observed mass-loss rate to these thermal-pulse periods shows that the present stellar wind leads to a very effective erosion of the envelope per thermal pulse cycle, possibly as high as $`1`$M/cycle. Consequently, the whole envelope may be lost during the next few thermal pulses leading to the termination of the AGB evolution. Thus, it is not unlikely to assume that IRC +10 216 has entered a phase immediately before moving off the AGB. This is strongly supported by the non-spherical appearence of its dust shell showing even bipolar structures. Unlike AGB stars, post-AGB objects as protoplanetary nebulae often expose prominent features of asphericities, in particular in axisymmetric geometry (e.g. Olofsson Oloff93 (1993), Oloff96 (1996)). Accordingly, IRC +10 216 can be thought to be an object in transition. It is noteworthy that the establishment of bipolar structures, i.e. the metamorphosis into a protoplanetary nebulae, obviously already begins during the (very end of) AGB evolution. The clumpiness within the bipolar shape is probably due to small scale fluctuations of the dust condensation radius which, in turn, might be influenced by, e.g., giant surface convection cells (Schwarzschild Schwarzschild75 (1975)). The formation of giant convection cells can be assumed to be a common phenomenon in red giants. The shaping of planetary nebulae can successfully be described by interacting stellar wind models (Kwok et al. Kwoketal78 (1978), Kahn & West KahnWest85 (1985)). Within this scenario a fast (spherical) wind from the central star interacts with the slow wind of the preceding AGB evolution. The slow AGB wind is asssumed to be non-spherical (axisymmetric) which leads to the observed morphology of planetary nebulae (Mellema Mel96 (1996)). The cause of an aspherical AGB mass loss is still a matter of debate. Different mechanisms to provide the required equatorial density enhancements are discussed (cf. Livio Liv93 (1993)). Among these, binarity is one channel including common envelope evolution and spin up of the AGB star due to the interaction with its companion (Morris Mor81 (1981), Bond & Livio BondLiv90 (1990)). Not only stellar companions are found to be able to spin up the AGB star but also substellar ones as brown dwarfs and planets, most effectively by evaporation in the AGB star’s envelope (Harpaz & Soker SokHar92 (1992), Soker Soker97 (1997)). Currently there is no observational evidence for a possible binary nature of IRC +10 216. The fact that the polarization pattern in the southeastern part of the nebula at 1.1 $`\mu `$m has a different orientation than in the rest of the nebula might be an indication of a second illuminating source. Mechanisms inherent to the star include rotation (Dorfi & Höfner DorHoef96 (1996), Garcia-Segura et al. GarBer99 (1999)), non-radial pulsations (Soker & Harpaz SokHar92 (1992)) and magnetic fields (Pacoli et al. Pacetal92 (1992), Garcia-Segura et al. GarBer99 (1999)). Both non-radial $`p`$-modes and magnetic fields appear to only be important for significant rotation rates. Often spin-up agents due to binarity are assumed. For instance, Groenewegen (Groenewegen96 (1996)) favours non-radial pulsation or an as yet unidentified companion which spun up the central star as the most likely explanation for the non-spherical shape of the dust shell of IRC +10 216. AGB stars are known to be slow rotators. Stars with initial masses below $`1.3`$M can be expected to lose almost their entire angular momentum during the main sequence phase due to magnetic braking operating in their convective envelopes. Consequently, they are not believed to devolop non-spherical mass-loss due to rotation. Stars with larger initial masses are spun down due to expansion and mass loss in the course of evolution, but may achieve sufficiently high rotation rates at the end of their AGB stage (Garcia-Segura et al. GarBer99 (1999)). Even small rotation rates influence dust-driven winds considerably, yielding a mass loss preferentially driven in the equatorial plane (Dorfi & Höfner DorHoef96 (1996)). Furthermore, for supergiants leaving the Hayashi line, Heger & Langer (HegLan98 (1998)) found that significant spin up of the surface layers may take place. Thus, on second glance, inherent rotation might be able to support axisymmetric mass loss during the transition to the proto-planetary nebula phase for more massive AGB stars as IRC +10 216. ## 4 Conclusions We have presented high-resolution $`J`$–, $`H`$–, and $`K`$–band observations of IRC +10 216 with the highest resolution so far at $`H`$ of 70 mas. A series of $`K`$–band images from five epochs between October 1995 and November 1998 shows that the inner nebula is non-stationary. The separations of the four dominant resolved components increased within the 3 years by up to $`35\%`$. For the two brightest components a relative velocity within the plane of the sky of about 23 mas/yr or 14 km/s was found. Within these 3 years the rather faint components C and D became brighter whereas component B faded. The general geometry of the nebula seems to be bipolar. We find that the most promising model to explain the structures and changes in the inner nebula is to assume that the star is at or near the position of component B. The star is then strongly but not totally obscured at $`H`$ and $`K`$. Consistently component B is very red in the $`HK`$ color while A and the northern $`J`$–band components are relatively blue. The polarization pattern with strong polarization in the northern arms and also a significant polarization in the peak supports this model. The inner nebula and the apparent motions seem to be rather symmetric around this position and the observed changes are consistent with the assumption of an enhanced mass loss becoming apparent at least in 1997. IRC +10 216 is without doubt in a very advanced stage of its AGB evolution. The observed bipolarity of its dust shell even reveals that it has possibly entered the phase of transformation into a protoplanetary nebula. ###### Acknowledgements. We thank the referee for valuable comments. This research has made use of the SIMBAD database, operated at CDS, Strasbourg, France. The HST images for the polarization analysis have been retrieved from the Hubble Data Archive operated at the STScI, Baltimore, USA.
warning/0003/cond-mat0003217.html	ar5iv	text	# Thermal diffusion of sine-Gordon solitons ## 1 Introduction As a key subject within nonlinear science, the dynamics of emergent, coherent structures (solitons, vortices, etc) has been a research topic that has attracted very much attention in the past quarter century scott . One question, extensively investigated in the literature bass ; kiv ; oldanx ; kvbook ; yuripr ; anx is the following: Is, and if so, how is the motion and the shape of those excitations modified by the presence of small perturbations? Indeed, when applied to physical situations of interest, nonlinear models must incorporate additional terms, such as damping, constant or periodic external forces, or noise, to name a few. Among those, stochastic perturbations are very much of interest in view of their highly non trivial effects on nonlinear systems gsbook , and a great deal of work has been devoted to them bass ; oldanx ; kvbook . In particular, of the very many nonlinear models applied to physical problems, the sine-Gordon (sG) equation has been considered in much detail in this context, as it applies to, e.g., propagation of ultra-short optical pulses in resonant laser media lamb , a unitary theory of elementary particles skyrme1 ; skyrme2 ; enz ; raja , propagation of magnetic flux in Josephson junctions Barone , transmission of ferromagnetic waves feld , epitaxial growth of thin films CW ; krug ; us , motion of dislocations in crystals frenkel ; Nabarro , flux-line unlocking in type II superconductors nuevo , or DNA dynamics eng ; DNA ; yaku , situations in which noise (of different origins) can play, and often does, a crucial role. As an example, let us mention the recent work on long Josephson junctions reported in exp , where the authors calculated the escape rate from the zero-voltage state induced by thermal fluctuations, obtaining very satisfactory results compared with the experimental ones. Specifically, this work is devoted to the study of the diffusive dynamics of sG kink solitons subjected to a thermal bath, as given by the stochastically perturbed, damped sG equation $$\varphi _{tt}\varphi _{xx}+\mathrm{sin}(\varphi )=\alpha \varphi _t+f(x,t,\varphi _,\mathrm{}),$$ (1) where $`\alpha \varphi _t`$ is a damping term with a dissipation coefficient $`\alpha `$, and $`f(x,t,\varphi _,\mathrm{})`$ is a thermal (gaussian) noise term fulfilling $$\begin{array}{c}f(x,t,\varphi _,\mathrm{})=\sqrt{D}\eta (x,t),\eta (x,t)=0,\hfill \\ \\ \eta (x,t)\eta (x^{},t^{})=\delta (xx^{})\delta (tt^{}),\hfill \end{array}$$ (2) where $`\sqrt{D}`$ is related to temperature through the fluctuation-dissipation theorem $`D=2\alpha k_bT`$, $`k_b`$ being the Boltzmann constant and $`T`$ the temperature. To our knowledge, the first results on problems directly related to the one we deal with here were obtained by Joergensen et al. joer , who performed experiments on Josephson junctions and presented a derivation of the diffusion constant for kinks. Subsequently, Kaup and Osman kaup studied, in a more rigorous way, the motion of damped sG kinks, driven by a constant force, in the presence of thermal fluctuations by using a singular perturbation expansion. They analyzed the temperature effect on the mean velocity of the kink and also the changes in the shape of the kink. In addition, they calculated the diffusion coefficient of the kink up to first-order in temperature and the energy values corresponding to the translational ($`E_T=k_bT/2`$) and radiational ($`E_R=k_bT`$) modes. These values of the energy have been also obtained by Marchesoni march , who applied the McLaughlin and Scott approach McL to investigate kink motion under thermal fluctuations (see bass ; kiv ; kvbook for reviews). For the sake of completeness, let us mention work done along a different line, namely that devoted to the diffusive motion of the kink in equilibrium with phonons in isolated sG systems (possibly perturbed) yuripr ; Wsch ; March2 ; March . In this case, the kink diffusive motion is characterized by two diffusion coefficients. The first one of them is proportional to $`T^2`$ and is related to the anomalous diffusion, that arises from the phase shifts of kinks colliding with phonons and takes place on a short time scale in which the collision among kink and wave packet is elastic; the kink retains the same velocity after the collision (non-dissipative diffusion) and suffers only a spatial shift. However, for large times and in slightly perturbed sG systems, this interaction is nonlinear and becomes inelastic and the velocity of the kink changes after the collisions Wada . This diffusive regime is called viscous and has a diffusion coefficient proportional to $`T^1`$. The diffusion of the kink when the low energy excitations are represented by breathers has also been studied, and in Niko it has been demonstrated that both descriptions (breathers or phonons) are equivalent and give rise to the same diffusion coefficient in the anomalous regime. In any event, we want to stress that, although in this type of diffusion problem there are many open questions Wada , we will concern ourselves with the other kind of diffusion problem, in which the phonons appear as a consequence of an external heat bath, represented by white noise correlated in space and in time and the damping is included explicitly à la Langevin. The main aim of this work is to extend a previous study of ours about the overdamped limit of sine-Gordon kink diffusion nos to the more physical and general case of the underdamped dynamics (i.e., with finite dissipation coefficient). As we will see below, the general perturbative approach raja we resorted to in the overdamped case can also be applied, albeit with more difficulties, to the underdamped problem. The corresponding theoretical analysis is presented in Sec. 2, where we obtain explicit expressions for the long-time diffusive dynamics of kinks up to second-order in temperature, thus going beyond the currently available knowledge. The accuracy and importance of the new terms is assessed by numerical simulations in Sec. 3: we will see there that the quadratic corrections are in good agreement with the simulations and, most importantly, that they must be taken into account even for not so large temperatures. Finally, in the conclusions we summarize the main results of this work, comparing the underdamped and overdamped dynamics of the sG equation and discussing other related questions. ## 2 Analytical results We begin by briefly reviewing the basic results we need for our analytical approach. We will concern ourselves with the perturbation effect on the kink solutions of the unperturbed sG equation, whose static form is $$\varphi _0(x,t)=4\text{arctan}[\mathrm{exp}(x)].$$ (3) Small perturbations over this equation can be treated by calculating the spectrum of linear excitations around the kink solution bis : To this end, we write $$\varphi (x,t)=\varphi _0(x)+\psi (x,t),\psi (x,t)<<\varphi _0(x),$$ (4) substitute in (1) (with $`\alpha =D=0`$) and linearize around $`\varphi _0(x)`$, arriving at the following equation for $`\psi (x,t)`$: $$\psi _{tt}=\psi _{xx}\left[1\frac{2}{\mathrm{cosh}^2(x)}\right]\psi .$$ (5) Then, assuming that the solution of (5) has the form $$\psi (x,t)=f(x)\mathrm{exp}\left(i\omega t\right)$$ (6) we find the eigenvalue problem for $`f(x)`$, $$\frac{^2f}{x^2}+\left[1\frac{2}{\mathrm{cosh}^2(x)}\right]f=\omega ^2f.$$ (7) This equation admits the following eigenfunctions with their respective eigenvalues $`f_T(x)={\displaystyle \frac{2}{\mathrm{cosh}(x)}},\omega _T^2`$ $`=`$ $`0,`$ (8) $`f_k(x)={\displaystyle \frac{\mathrm{exp}(ikx)[k+i\text{tanh}(x)]}{\sqrt{2\pi }\omega _k}},\omega _k^2`$ $`=`$ $`1+k^2,`$ (9) which represent, respectively, the translation (Goldstone) mode and the radiation modes. Importantly, the functions $`f_T(x)`$ and $`f_k(x)`$ form a complete set with the orthogonality relations $`{\displaystyle _{\mathrm{}}^+\mathrm{}}f_T^2(x)𝑑x`$ $`=`$ $`8,`$ (10) $`{\displaystyle _{\mathrm{}}^+\mathrm{}}f_T(x)f_k(x)𝑑x`$ $`=`$ $`0,`$ (11) $`{\displaystyle _{\mathrm{}}^+\mathrm{}}f_k(x)f_k^{}^{}(x)𝑑x`$ $`=`$ $`\delta (kk^{}).`$ (12) We can now proceed with our problem: In order to tackle Eq. (1), with noise as given in (2), we use the same Ansatz proposed for the overdamped case in nos (or for the general Klein-Gordon system in raja ): We assume that the solution of Eq. (1) is $$\varphi (x,t)=\varphi _0[xX(t)]+_{\mathrm{}}^+\mathrm{}𝑑kA_k(t)f_k[xX(t)],$$ (13) where $`X(t)`$ is the kink position. We now insert (13) in (1) and use the orthogonality of $`f_k`$ and $`f_T`$ bis , obtaining the following system of differential equation for $`X(t)`$ and $`A_k(t)`$: $`\ddot{X}(t)+\alpha \dot{X}(t)`$ $`=`$ $`{\displaystyle \frac{\alpha }{8}}\dot{X}(t){\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑kA_k(t)I_1(k){\displaystyle \frac{1}{16}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k^{}A_k(t)A_k^{}(t)R_3(k,k^{})+`$ (14) $`+`$ $`{\displaystyle \frac{\sqrt{D}}{8}}{\displaystyle _{\mathrm{}}^+\mathrm{}}f_T[xX(t)]\eta (x,t)𝑑x`$ $``$ $`{\displaystyle \frac{1}{48}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k_1{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k_2A_k(t)A_{k_1}(t)A_{k_2}(t)R_6(k,k_1,k_2)`$ $``$ $`{\displaystyle \frac{\dot{X}(t)}{4}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle \frac{A_k}{t}}I_1(k){\displaystyle \frac{\ddot{X}(t)}{8}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑kA_k(t)I_1(k)+`$ $`+`$ $`{\displaystyle \frac{\dot{X}^2(t)}{8}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle \frac{A_k}{t}}I_2(k),`$ $`{\displaystyle \frac{^2A_k}{t^2}}+\alpha {\displaystyle \frac{A_k}{t}}+\omega _k^2A_k(t)`$ $`=`$ $`\alpha \dot{X}(t){\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k^{}A_k^{}(t)I_3(k^{},k)+`$ (15) $`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k^{}A_k(t)A_k^{}(t)R_4(k,k^{})`$ $``$ $`\sqrt{D}{\displaystyle _{\mathrm{}}^+\mathrm{}}f_k^{}[xX(t)]\eta (x,t)𝑑x+`$ $`+`$ $`{\displaystyle \frac{1}{6}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k^{}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k_1{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k_2A_k^{}(t)A_{k_1}(t)A_{k_2}(t)R_7(k^{},k,k_1,k_2)+`$ $`+`$ $`2\dot{X}(t){\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k^{}{\displaystyle \frac{A_k^{}}{t}}I_3(k^{},k)+`$ $`+`$ $`\ddot{X}(t){\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k^{}A_k^{}(t)I_3(k^{},k)+\dot{X}^2(t)I_1(k),`$ where $`I_1(k)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{f_k}{\theta }}f_T(\theta )𝑑\theta ={\displaystyle \frac{i\pi \omega _k}{\sqrt{2\pi }\mathrm{cosh}\left({\displaystyle \frac{\pi k}{2}}\right)}},`$ $`I_2(k)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{^2f_k}{\theta ^2}}f_T(\theta )𝑑\theta ,`$ $`R_3(k,k^{})`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}f_T(\theta ){\displaystyle \frac{f_T}{\theta }}f_k(\theta )f_k^{}^{}(\theta )𝑑\theta ={\displaystyle \frac{i(\omega _k^2\omega _k^{}^2)^2}{4\omega _k\omega _k^{}\mathrm{sinh}\left({\displaystyle \frac{\pi \mathrm{\Delta }k}{2}}\right)}},\mathrm{\Delta }k=k^{}k,`$ $`I_3(k,k^{})`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{f_k}{\theta }}f_k^{}^{}(\theta )𝑑\theta ,`$ $`R_4(k,k^{})`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}[f_k^{}^{}(\theta )]^2{\displaystyle \frac{f_T}{\theta }}f_k(\theta )𝑑\theta ,R_4(k,k)={\displaystyle \frac{3i\omega _k}{8\sqrt{2\pi }\mathrm{cosh}\left({\displaystyle \frac{\pi k}{2}}\right)}},`$ $`R_6(k,k_1,k_2)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{^2f_T}{\theta ^2}}f_k(\theta )f_{k_1}^{}(\theta )f_{k_2}(\theta )𝑑\theta ,`$ $`R_7(k,k^{},k_1,k_2)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}\mathrm{cos}(\varphi _0)f_k^{}^{}(\theta )f_k(\theta )f_{k_1}^{}(\theta )f_{k_2}(\theta )𝑑\theta .`$ (16) It goes without saying that these equations can not be solved. Therefore, in order to extract information from them, we resort to a perturbative approach assuming the noise term is small, or equivalently, that the temperature and the dissipation are not too large (this is not a serious restriction since our single-kink approach does not apply to high temperatures, when kink-antikink pairs are thermally generated butt ). We then expand $`X(t)`$ and $`A_k(t)`$ in powers of $`\sqrt{D}`$, i.e., $`X(t)=_{n=1}^{\mathrm{}}(\sqrt{D})^nX_n(t)`$ and $`A_k(t)=_{n=1}^{\mathrm{}}(\sqrt{D})^nA_k^n(t)`$, since when $`\sqrt{D}=0`$ and $`\alpha =0`$ we recover the static kink solution (in this case initially centered at the origin) of the sG equation. By substituting these expansions in (14) and (15) we find a set of linear equations for the coefficients of these series. The first members of this hierarchy correspond to order $`O(\sqrt{D})`$: $$\ddot{X}_1(t)+\alpha \dot{X}_1(t)=\frac{1}{8}_{\mathrm{}}^+\mathrm{}f_T[xX(t)]\eta (x,t)𝑑xϵ_1(t),$$ (17) from where we obtain the statistical properties of $`ϵ_1(t)`$, $`ϵ_1(t)=0,ϵ_1(t)ϵ_1(t^{})={\displaystyle \frac{1}{8}}\delta (tt^{}),`$ (18) and $$\frac{^2A_k^1}{t^2}(t)+\alpha \frac{A_k^1}{t}(t)+\omega _k^2A_k^1(t)=_{\mathrm{}}^+\mathrm{}f_k^{}[xX(t)]\eta (x,t)𝑑x\xi _k(t),$$ (19) which in turn leads to $`\xi _k(t)=0,\xi _k(t)\xi _k^{}(t^{})=\delta (tt^{})\delta (kk^{}).`$ (20) Equations (17)–(20) have been obtained in sal by using a similar, but more restrictive perturbative approach. By integrating these two equations we obtain the first-order terms, $`X_1(t)`$ and $`A_k^1`$: $$\begin{array}{c}X_1(t)=_0^te^{\alpha t^{}}_0^t^{}e^{\alpha \tau }ϵ_1(\tau )𝑑\tau 𝑑t^{},A_k^1(t)=e^{{\scriptscriptstyle \frac{\alpha t}{2}}}\left\{C_1(t)\mathrm{sin}\omega t+C_2(t)\mathrm{cos}\omega t\right\},\hfill \\ \\ C_1(t)=\frac{1}{\omega }_0^t\xi _k(\tau )e^{{\scriptscriptstyle \frac{\alpha \tau }{2}}}\mathrm{cos}\omega \tau d\tau ,C_2(t)=\frac{1}{\omega }_0^t\xi _k(\tau )e^{{\scriptscriptstyle \frac{\alpha \tau }{2}}}\mathrm{sin}\omega \tau d\tau ,\hfill \end{array}$$ (21) where $`\omega ^2=\omega _k^2(\alpha ^2/4)`$. From these relations we can calculate the mean values and correlation functions up to first order in $`\sqrt{D}`$: $`X_1(t)=0,X(t)X(t^{})`$ $`=`$ $`DX_1(t)X_1(t^{})={\displaystyle \frac{D}{16\alpha ^3}}[e^{\alpha M}e^{\alpha \|\mathrm{\Delta }t\|}+e^{\alpha M\alpha \|\mathrm{\Delta }t\|}`$ (22) $``$ $`e^{\alpha (t+t^{})}+e^{\alpha t}+e^{\alpha t^{}}+2(\alpha M1)],`$ $`\dot{X}_1(t)`$ $`=`$ $`0,\dot{X}(t)\dot{X}(t^{})=D\dot{X}_1(t)\dot{X}_1(t^{})={\displaystyle \frac{D}{16\alpha }}\left[e^{\alpha \|\mathrm{\Delta }t\|}e^{\alpha (t+t^{})}\right],`$ (23) $`A_k^1(t)=0,A_k(t)A_k(t^{})`$ $`=`$ $`DA_k^1(t)A_k^1(t^{})={\displaystyle \frac{D}{\omega ^2}}e^{\alpha (t+t^{})/2}[{\displaystyle \frac{e^{\alpha M}1}{2\alpha }}\mathrm{cos}\omega \mathrm{\Delta }t`$ (24) $``$ $`{\displaystyle \frac{\alpha e^{\alpha M}}{8\omega _k^2}}\mathrm{cos}\omega \mathrm{\Delta }t{\displaystyle \frac{\omega e^{\alpha M}}{4\omega _k^2}}\mathrm{sin}\omega \|\mathrm{\Delta }t\|+{\displaystyle \frac{\alpha }{8\omega _k^2}}\mathrm{cos}\omega (t+t^{})`$ $``$ $`{\displaystyle \frac{\omega }{4\omega _k^2}}\mathrm{sin}\omega (t+t^{})],`$ where $`\mathrm{\Delta }t=tt^{}`$, and $`M=\mathrm{min}(t,t^{})`$. Of course, for $`t^{}=t`$ in Eq. (22) we recover the result in kaup for $`[X(t)]^2`$. We now turn to the main point of our work, namely obtaining the next-order corrections for the position and the velocity of the center of the kink. This requires the calculation of the next two contributions to $`X(t)`$ as well as the second order in the radiation terms, which are: $`O(D)`$ $`\ddot{X}_2(t)+\alpha \dot{X}_2(t)`$ $`=`$ $`ϵ_2(t),`$ (25) $`ϵ_2(t)`$ $``$ $`{\displaystyle \frac{ϵ_1(t)}{8}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑kA_k^1(t)I_1(k){\displaystyle \frac{\dot{X}_1(t)}{4}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle \frac{A_k^1}{t}}(t)I_1(k)`$ (26) $``$ $`{\displaystyle \frac{1}{16}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k^{}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑kA_k^1(t)A_k^{}^1(t)R_3(k,k^{}),`$ $`{\displaystyle \frac{^2A_k^2}{t^2}}(t)+\alpha {\displaystyle \frac{A_k^2}{t}}(t)+\omega _k^2A_k^2(t)`$ $`=`$ $`ϵ_1(t){\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k^{}A_k^{}^1(t)I_3(k^{},k)+`$ (27) $`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k^{}A_k^1(t)A_k^{}^1(t)R_4(k,k^{})`$ $``$ $`\dot{X}_1^2(t)I_1(k)+2\dot{X}_1(t){\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k^{}{\displaystyle \frac{A_k^{}^1}{t}}(t)I_3(k^{},k);`$ $`O([\sqrt{D}]^3)`$ $`\ddot{X}_3(t)+\alpha \dot{X}_3(t)`$ $`=`$ $`ϵ_3(t),`$ (28) $`ϵ_3(t)`$ $``$ $`{\displaystyle \frac{ϵ_1(t)}{8}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑kA_k^2(t)I_1(k){\displaystyle \frac{ϵ_2(t)}{8}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑kA_k^1(t)I_1(k)`$ (29) $``$ $`{\displaystyle \frac{1}{16}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k^{}A_k^2(t)A_k^{}^1(t)R_3(k,k^{})`$ $``$ $`{\displaystyle \frac{1}{16}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k^{}A_k^1(t)A_k^{}^2(t)R_3(k,k^{})`$ $``$ $`{\displaystyle \frac{1}{48}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k_1{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k_2A_k^1(t)A_{k_1}^1(t)A_{k_2}^1(t)R_6(k,k_1,k_2)+`$ $`+`$ $`{\displaystyle \frac{\dot{X}_1^2(t)}{8}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle \frac{A_k^1}{t}}(t)I_2(k){\displaystyle \frac{\dot{X}_1(t)}{4}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle \frac{A_k^2}{t}}(t)I_1(k)`$ $``$ $`{\displaystyle \frac{\dot{X}_2(t)}{4}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle \frac{A_k^1}{t}}(t)I_1(k).`$ Analogously to what we have done for the first-order corrections, from the solutions of Eqs. (25) and (28) we find that $`X_2(t)`$ $`=`$ $`0,\dot{X}_2(t)=0,`$ (30) $`X_3(t)`$ $`=`$ $`0,\dot{X}_3(t)=0.`$ (31) As for higher moments, taking into account that the cross-correlation function of $`X_1(t)`$ and $`X_3(t^{})`$ is of the same order as $`X_2(t)X_2(t^{})`$, and also that $`X_1(t)X_2(t^{})=0`$ we obtain that $`[X(t)]^2`$ $`=`$ $`D[X_1(t)]^2+`$ (32) $`+`$ $`D^2([X_2(t)]^2+2X_1(t)X_3(t))+\mathrm{},`$ $`[\dot{X}(t)]^2`$ $`=`$ $`D[\dot{X}_1(t)]^2+`$ (33) $`+`$ $`D^2([\dot{X}_2(t)]^2+2\dot{X}_1(t)\dot{X}_3(t))+\mathrm{}`$ The expressions for the functions $`[X_2(t)]^2`$, $`[\dot{X}_2(t)]^2`$, $`X_1(t)X_3(t)`$, and $`\dot{X}_1(t)\dot{X}_3(t)`$ can be obtained after a lengthy calculation, and are very cumbersome indeed. We therefore do not include them here. However, for large time ($`t>>1/\alpha `$) these relations can be simplified, yielding, as $`t\mathrm{}`$ $`[X(t)]^2`$ $`=`$ $`{\displaystyle \frac{k_bTt}{4\alpha }}\left\{1+{\displaystyle \frac{k_bT}{32}}\left(1+{\displaystyle \frac{9\sigma ^2}{4}}\right)\right\},`$ (34) $`[\dot{X}(t)]^2`$ $`=`$ $`{\displaystyle \frac{k_bT}{8}}\left\{1+{\displaystyle \frac{3k_bT}{128}}\left(12+\sigma ^2\right)\right\},`$ (35) with $$\sigma =_{\mathrm{}}^+\mathrm{}\frac{dk}{\omega _k\mathrm{cosh}\left({\displaystyle \frac{\pi k}{2}}\right)}=1.62386.$$ (36) To complete the characterization of the kink diffusion, we can now compute in a straightforward way the average value of the wave function $`\varphi (x,t)`$, defined as $`\varphi (x,t)`$ $`=`$ $`\varphi _0[x\sqrt{D}X_1(t)]+O(D)=`$ (37) $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑X_1p(X_1)\varphi _0[x\sqrt{D}X_1(t)],`$ where $`p(X_1)`$ is the probability distribution function for $`X_1`$. To find explicitly this function we note that, if we rewrite Eq. (17) as a system of two differential equations, $`\dot{X}_1`$ $`=`$ $`V,`$ $`\dot{V}`$ $`=`$ $`\alpha V+ϵ_1(t),`$ (38) the last equation represents an Ornstein-Uhlenbeck process for the velocity, and its distribution function is given by $$p(V)=\sqrt{\frac{1}{2\pi V^2}}\mathrm{exp}\left(\frac{V^2}{2V^2}\right),$$ (39) (see noise ). Subsequently, by integrating the first equation of (38), we obtain that $`X_1=_0^tV(\tau )𝑑\tau `$. Since $`V`$ has a Gaussian distribution function, $`X_1`$ has also a Gaussian distribution function, given by (recall that $`X_1(t)=0`$) $$p(X_1)=\sqrt{\frac{1}{2\pi [X(t)]^2}}\mathrm{exp}\left(\frac{1}{2}\frac{X_1^2}{[X(t)]^2}\right),$$ (40) where the first and second moments of $`X_1`$ were obtained before, see Eq. (22). With this result, the integral (37), can be evaluated numerically taking into account Eqs. (40) and (22). In the next section we will compare this result with the mean value of the wave function as obtained from simulations of the full partial differential equation (1). ## 3 Numerical simulations In order to test the approximate theory developed in the previous section, we have simulated numerically Eq. (1) by using the Heun method maxra . In our simulations we begin with a kink, initially at rest, with free boundary conditions. For the damping coefficient we choose $`\alpha =0.1`$, which is not too small because from (34) we can see that $`[X(t)]^2`$ is proportional to $`1/\alpha `$. This means that if $`\alpha `$ is too small the kink can move in a much larger region, forcing us to increase the length of our simulated system in the simulations, already quite time consuming. Furthermore, the relation (34) is only valid for large times ($`t>>1/\alpha `$). Again, for too small $`\alpha `$ we would need to simulate our equation for very long times and, as $`[X(t)]^2`$ increases linearly with time \[see Eq. (34)\], the system length would once more have to be large. The other parameters are $`\mathrm{\Delta }x=0.2`$, $`\mathrm{\Delta }t=0.001`$ and the length of the system $`L=200`$. We have calculated all average values over 1000 realizations up to a final time $`400`$. It is important to point out that, this system being inertial, the accuracy of the averages is considerably less than for overdamped problems, this being the reason why we have to use such large ensembles of trajectories to obtain reasonably good results. An important, nontrivial issue is the question as to how can we find the center of the kink. We solve this problem by finding all the discrete lattice points $`x_i`$ and $`x_{i+1}`$ such that $`\varphi _i\pi `$ and $`\varphi _{i+1}\pi `$ or vice versa, and then estimating the corresponding points $`\stackrel{~}{x}_i`$ where $`\varphi =\pi `$ by linear interpolation. Afterwards, among the $`n`$ such points $`\stackrel{~}{x}_n`$, we choose to be the center of the kink the value $`\stackrel{~}{x}=\stackrel{~}{x}_n`$, which minimizes $`_{i=1}^{L/\mathrm{\Delta }x}[\varphi _i(t)\varphi _0(x\stackrel{~}{x}_n,t)]^2`$, i.e., the discrete version of the integral of the square of the difference of $`\varphi `$ and $`\varphi _0`$. It has to be realized that this involves an assumption, namely that individual realizations of the kink have a shape similar to that of the unperturbed kink. As can be seen from Fig. 1, where the individual realizations are compared with the initial condition (an exact kink), this is indeed the case and our procedure is truly sensible. Therefore, we are sure that this method to compute the kink center avoids the false centers, which can appear for higher temperatures due to fluctuations introducing a systematic difference between the numerical and the theoretical results (see nos ). With the procedure we have just summarized, that works even for relatively large temperatures, we believe we find a very accurate approximation to the actual center of the kink. We will come back to Fig. 1 below. As an example of the comparison of the numerical simulations of Eq. (1) with the theoretical results obtained in the previous section and valid for large times, Fig. 2 shows the numerically computed variance of the center of the kink, $`[X(t)]^2X(t)^2`$, as well as the first- and second-order analytical expressions. The plot clearly evidences that the numerical variance asymptotically coincides with the second-order expression: Note that to compare the different curves one has to look at the slopes at times $`t>>1/\alpha `$ (in this case, $`t100`$, for instance, as $`\alpha =0.1`$); the theoretical result is not valid at early times and therefore there is a bias between analytics and numerics coming from that. The small, irregular oscillations in the numerical curve arise from the difficulty in accurately computing averages in an underdamped system like this mentioned above; however, we believe that the present accuracy is enough to confirm the validity of our approach. We have observed the same agreement for other values of temperatures ($`k_bT=0.2,0.6,0.8`$, not shown). In all cases, we have computed the diffusion coefficient for large times as the slope of the variance of $`X(t)`$ again for $`t100`$, the regime in which we expect our analytical approximation to be valid. Summarizing our results, these numerical values of the diffusion coefficient are plotted in Fig. 3 together with the theoretical results. It is clear that for large temperatures the quadratic behavior in $`k_bT`$ of the diffusion coefficient becomes important. For higher values of the temperature, such as $`k_bT=0.8`$, the numerical value of the diffusion coefficient is not so close to the predicted one. This effect arises because of the large diffusivity of the kink in that range: Indeed, for this and higher temperatures the kink performs very long excursions away from the center, reaching the boundaries of the numerical integration interval; it is clear that when this occurs, the diffusion of the kink is not in free space anymore and hence those realizations spoil the quality of the averages. The way to solve this problem would be to resort to much larger numerical systems, but within our present computing capabilities this would necessitate a simultaneous decrease in the number of realizations in the average, leading again to poorer results. However, it is important to realize that this boundary effect leads to an underestimation of the diffusion coefficient (as the boundary prevents the kink from travelling as far as it should) and therefore the point in Fig. 3 for $`k_bT=0.8`$ is a lower bound for the diffusion coefficient, with the actual one lying even closer to our second order prediction. Finally, there is one last question that deserves discussion, namely that of the physical significance of the mean value of the wave function $`\varphi `$. In Fig. 1 we can clearly see that, whereas individual realizations of kinks look very similar to the unperturbed ones, the mean value of $`\varphi `$ is a much wider excitation, not even close to the original kink. Figure 1 clearly shows that this does not mean that the width of individual kinks increases; indeed, much as we discussed regarding the overdamped problem nos , we have verified numerically that the mean wave function $`\varphi (x,t_{fix})`$ increases due to the variance of the kink center of individual realizations, and hence it should not be interpreted as the typical deformation of the shape of kinks: Indeed, the widening of the mean value of $`\varphi `$ arises from the contributions of the stochastically moving, but mostly undistorted kinks whose center positions have the distribution of a rigid, diffusing particle. To further check this interpretation, we can look at Fig. 4, where we have represented the mean value of the wave function for two fixed times $`t_{fix}=100,300`$, obtained from the numerical simulation of the full partial differential equation (1), for $`k_bT=0.6`$ and $`\alpha =0.1`$. The overimposed points, computed by using the Gaussian distribution function $`p(X_1)`$ \[Eq. (40)\] of the kink center $`X(t)=\sqrt{D}X_1`$ found in the last section, show the excellent agreement between our theory and the simulation. Of course, there is a small discrepancy that is likely to disappear if one would go to a next order calculation, but for the present purpose of understanding the mean wave function $`\varphi `$ the first order calculation is enough. In addition, we have plotted the initial kink (at rest) in order to see that the mean value of the wave function increases with time. ## 4 Discussion and conclusions To summarize, in this work we have studied the diffusive dynamics of sine-Gordon kinks subjected to thermal fluctuations. We have analytically calculated expressions, valid up to second order in temperature, for the average position and variance of the kink center, as well as for the mean shape of the kink. We have numerically checked the validity of these results up to temperatures of the order of $`k_bT=0.8`$ (in dimensionless units, equivalent to about a 10% of the kink rest mass), already close to the temperature at which kink-antikink nucleation becomes a likely event. Therefore, our first conclusion is that the second-order theory developed here is the proper one, meaning it is accurate and higher order terms are negligible, to describe the thermal diffusion of sine-Gordon kinks in the single kink propagation regime. Interestingly, our calculation pinpoints the fact that the second-order correction in $`k_bT`$ comes from the interaction between kink and phonons. This implies that the physics behind this contribution is basically the same as for the case of anomalous diffusion in an isolated chain mentioned in the introduction yuripr ; Wsch ; March2 ; March . Note that we do not expect $`T^1`$ contributions in our analytical calculations, as they are carried out in a continuum sG equation yuripr and, in any case, they would show up in simulations only for very low temperatures. Apart from that, it is also interesting to note that, according to Eq. (35), the second order term implies an increase of the energy carried by the kink beyond the $`k_bT/2`$ predicted by statistical mechanics (recall that the kink mass is 8 in our units). This can be interpreted in the following way: The kink is dressed by phonons which increase its mass. Thus, the kink energy would be $`M(T)\dot{X}^2/2`$, with a temperature dependent mass $`M(T)`$ whose expression can be easily found from Eq. (35). In order to confirm this interpretation, one could compute the energy carried by the phonons which dress the kink, but we believe that it is not necessary because, on the one hand, it would be a rather involved calculation (far beyond the scope of this work) and, on the other hand, we do not think that there is any other possible interpretation of this result. A second relevant point of this study relates to the numerical simulation and center location procedures. As this is an underdamped (inertial) system, the thermal mobility of kinks is quite large, the larger the higher the temperature. Because of this, we have not been able to obtain very precise numerical averages at the top of the temperature range studied, since the lengths of the systems and the number of realizations required are very large and consequently time consuming. However, we believe that the results presented here are enough to verify our theory. This is reinforced by the very good agreement between analytics and numerics regarding the mean shape of the field, even for temperatures as large as $`k_bT=0.6`$ (see Fig. 3), which shows that our approach indeed captures the physics of the diffusion process. In addition, we want to emphasize that, to our knowledge, we have designed a new algorithm to detect the kink center which gives very good results even for the highest temperatures studied, where previous researches, such as nos , had found problems arising from the many false centers detected. Another important issue is the comparison of the present analysis to that in nos for the overdamped problem. We have found that the diffusion coefficient given by (34) for the present case practically coincides with that obtained in nos for the overdamped limit of the equation: the difference in the second order is approximately $`0.06k_bT`$, i.e., very small compared to the magnitude of the quadratic contribution itself. Furthermore, the width of the mean value of the wave function increases with time for the overdamped case nos in the same manner as that reported here. Therefore, we can conclude that for large times the dynamic of the underdamped sG kinks is very similar to the overdamped case. This is an important point, because in principle one can expect similar results for other kink-bearing systems such as the $`\varphi ^4`$ equation, for instance, whose overdamped diffusive dynamics is known (see dzi for the $`\varphi ^4`$ case), thus avoiding the much more involved calculation of the underdamped case. Finally, we want to mention the relevance of this work to experimental systems, such as long Josephson junctions. As has been shown in exp , the thermal sG equation (1) is a good description of the physics of in-line Josephson junctions (although different boundary conditions are needed in that case). The work in exp compared the predictions from the sG model to experimentally measured escape rates from the zero voltage state. Therefore, it should be possible to design similar experiments in order to test our results and, specifically, the increased (quadratic) diffusivity of kinks at higher temperatures vs the linear behavior at lower ones. We hope that our theoretical work stimulates further experimental research in that direction. ## Acknowledgement We are grateful to A. R. Bishop for his comments and Esteban Moro for his suggestions. Work at GISC (Leganés) has been supported by CICyT (Spain) grant MAT95-0325 and DGES (Spain) grant PB96-0119. Travel between Bayreuth and Madrid is supported by “Acciones Integradas Hispano-Alemanas”, a joint program of DAAD (Az. 314-AI) and DGES. This research is part of a project supported by NATO grant CRG 971090.
warning/0003/cond-mat0003410.html	ar5iv	text	# Anderson-type model for a molecule adsorbed on a metal surface ## 1 Introduction In the last decade, a lot of progress has been made in the theoretical and experimental investigation of nanostructures like quantum dots and quantum wires . These structures are produced by lithography. Using this very flexible technology, a large variety of different structures with different geometries has been made. Simple dots, multiple dots, dot arrays and the corresponding leads have been realized both using metals and semiconductors. A disadvantage of lithographic methods is the presence of fluctuations in shape and size of these structures; they are unavoidable because of errors in writing the mask and because of the stochastic nature of the etching process. The idea to use molecules or supramolecular structures as quantum wires and quantum dots has been around for quite a while, see Ref. for an early suggestion of a molecular rectifier and Refs. for recent reviews. Powerful chemical synthesis methods are available that allow the production of atomic and molecular clusters with linear dimensions of up to 5 nm, i.e., they are approaching the lowest linear dimensions of nanostructures produced by lithography. It is a fascinating idea that such clusters which would have custom-made electronic properties could be included in nanostructures and replace, say, a quantum dot. Different copies of these systems would be identical in structure and properties (e.g., there is no difference in the composition of different C<sub>60</sub> molecules). First steps in the direction of this ‘bottom-up approach’ are under way in a number of labs around the world. Scanning probe techniques have been used to arrange atoms on a substrate in arbitrary ways . A particularly impressive example is the ‘quantum corral’ , a circularly shaped structure of Fe atoms on a Cu(111) surface that may be considered a precursor of a device. Carbon nanotubes have been contacted and proven to exhibit Coulomb-blockade behavior . A new kind of transistor whose central element is a C<sub>60</sub> molecule has been proposed . The theory of conduction through molecular wires has been studied by a number of groups (see, e.g., Refs. ). Also, the thermal conduction through molecular wires has been investigated . Recently, the coupling of a molecule to a metal has been studied both experimentally and theoretically . A central problem facing the use of molecules as electronic devices is the question of how to contact them in a reproducible way. An interesting answer to this question was given in a recent experiment in which a custom-designed “lander” molecule was adsorbed on a $`(111)`$ copper surface. The molecules can be thought of as consisting of an aromatic platform and four spacer “legs”. The spacer legs keep the main board at such a distance from the metal surface that the aromatic part is electronically decoupled from the surface. The molecules self-assembled perpendicularly to a step on the surface, thereby forming a contact between the metal and the aromatic board (the length of the legs was chosen to be comparable to the step height). A scanning-tunneling microscope (STM) was used to scan along the surface and the molecule, and it was found that the tunneling conductance decayed in an exponential way along the molecule. In this paper, we would like to study a (relatively) simple model for the molecule at the metal step. In contrast to the quantum-chemical calculations provided in that take into account the precise structure of the lander molecule, we propose to use a modified Anderson model to describe this situation. The ‘impurity’ of the new model is spatially extended, i.e., contains internal degrees of freedom not present in the standard Anderson model. The aim will be to calculate the local density of states of an extended structure (the molecule) using a Green’s function method. The molecule will be approximated by a (finite) one- or two-dimensional tight-binding lattice with a Hubbard-like interaction. Our goal is to gain a qualitative understanding of the physical problem and its phenomenology (as opposed to quantum-chemical calculations for specific molecules) to be able to propose and analyze more complicated (e.g., multiply-connected) structures in the future. ## 2 Model The system we have in mind is shown in Fig. 1. Our model is defined by a Hamiltonian that consists of three parts, $`H=H_c+H_d+H_T`$. The metal to which the molecular wire is attached is described by non-interacting electrons in a conduction band $$H_c=\underset{𝐤\sigma }{}ϵ_𝐤c_{𝐤\sigma }^{}c_{𝐤\sigma },$$ (1) where $`𝐤`$ and $`\sigma =,`$ denote electron wave vectors and spins, respectively. The single-electron energy $`ϵ_𝐤`$ is measured with respect to the Fermi energy ($`ϵ_F=0`$). The molecular wire is assumed to be a one-dimensional lattice of $`L`$ atomic sites and is described within the Hubbard model $$H_d=\underset{i,j=1}{\overset{L}{}}t_{ij}d_{i\sigma }^{}d_{j\sigma }+U\underset{j=1}{\overset{L}{}}n_jn_j.$$ (2) The parameters that enter here are the on-site electron-electron interaction $`U`$ and the matrix elements $`t_{ij}=ϵ_d\delta _{ij}t\delta _{i,j\pm 1}`$, where $`ϵ_d`$ is the on-site energy of the single-particle level at each site and $`t`$ is the hopping element between nearest neighbors. The coupling of the wire to the metal is assumed to be small so that we can describe it by the tunneling amplitude $`V_𝐤`$ of the electron in the state $`(𝐤,\sigma )`$ to the first site $`j=1`$: $$H_T=\underset{𝐤\sigma }{}(V_𝐤c_{𝐤\sigma }^{}d_{1\sigma }+h.c.).$$ (3) The model defined by $`H`$ is a generalization of the Anderson impurity model: for $`L=1`$, the molecule reduces to a localized level, and (3) is the hybridization term of the Anderson impurity model. The qualitative nature of the spectrum of the system is shown in the right part of Fig. 1. We characterize the metal by a model density of states $`N_c(E)`$ to be discussed below and shown schematically in (a). If the molecule is not coupled to the metal, $`V=0`$, its spectrum consists of discrete levels as shown in (b). Coupling the molecule to the metal broadens the spectral lines as shown in (c). If we assume a filling of one electron per site, the presence of strong on-site interactions leads to the formation of a Mott-Hubbard gap as shown in (d). We will use the imaginary-time path-integral formalism at finite temperature $`k_BT1/\beta `$ (see, e.g., Ref. ). In this formalism, the Euclidean action is given by $$S^E=_0^\beta 𝑑\tau \left(\underset{𝐤\sigma }{}c_{𝐤\sigma }^{}_\tau c_{𝐤\sigma }+\underset{j\sigma }{}d_{j\sigma }^{}_\tau d_{j\sigma }+H\right),$$ (4) and we want to calculate the thermal Green’s function $$𝒢(j\sigma ,k\sigma ^{};\tau )=T_\tau d_{j\sigma }^{}(0)d_{k\sigma ^{}}(\tau )$$ (5) and its Fourier transform defined by $`𝒢(j\sigma ,k\sigma ^{};i\omega _n)=_0^\beta 𝑑\tau e^{i\omega _n\tau }𝒢(j\sigma ,k\sigma ^{};\tau )`$, where $`\omega _n(2n+1)\pi /\beta `$ with $`n`$ integer are the Matsubara frequencies. First, we consider the non-interacting case, $`U=0`$. In this case, the Euclidean action (4) is quadratic both in $`c_{𝐤\sigma }(c_{𝐤\sigma }^{})`$ and $`d_{j\sigma }(d_{j\sigma }^{})`$, and it follows that $`𝒢(j\sigma ,k\sigma ^{};i\omega _n)=\delta _{\sigma \sigma ^{}}𝒢(j,k;i\omega _n)`$, $$𝒢^1(j,k;i\omega _n)=𝒢_0^1(j,k;i\omega _n)\mathrm{\Sigma }_c(j,k;i\omega _n),$$ (6) or more explicitly $$𝒢(i,j)=𝒢_0(i,j)+\frac{𝒢_0(i,0)\mathrm{\Sigma }_c(0,0)𝒢_0(0,j)}{1𝒢_0(0,0)\mathrm{\Sigma }_c(0,0)}.$$ (7) The unperturbed (non-interacting, isolated) Green’s function is given by $$𝒢_0^1(j,k;i\omega _n)=i\omega _nt_{jk}.$$ (8) In Eqs. (6) and (7), the coupling to the metal manifests its effect on the wire through the self-energy $$\mathrm{\Sigma }_c(j,k;i\omega _n)=\delta _{j1}\delta _{k1}_{\mathrm{}}^{\mathrm{}}\frac{dϵ}{2\pi }\frac{\mathrm{\Gamma }(ϵ)}{i\omega _nϵ}.$$ (9) The effect of the metal on the molecular wire is given by the function $$\mathrm{\Gamma }(E)2\underset{𝐤}{}\|V_𝐤\|^2\mathrm{Im}G_c^R(𝐤,E),$$ (10) where $`G_c^R(𝐤,E)`$ is the retarded Green’s function for the (unperturbed) metal. In the simplest case in which the energy-dependence of the tunneling amplitude $`V_𝐤`$ is not too large, $`V_𝐤V`$, (10) can be further reduced to an expression directly proportional to the density of states in the metal $`N_c(E)`$: $`\mathrm{\Gamma }(E)=2\pi V^2N_c(E)`$. The local density of states (LDOS) at the $`j`$th site is then given by the analytic continuation of the thermal Green’s function. $$\rho (j,E)=\frac{2}{\pi }\mathrm{Im}𝒢(j,j;i\omega _nE+i0^+).$$ (11) The factor 2 accounts for the contributions from the two spin components. We now turn to the opposite limit in which the electron-electron interaction on the molecule is very strong; $`Ut`$ . In this case, exact solutions are only available in special cases; the isolated molecule ($`V_𝐤=0`$), i.e, the usual Hubbard model and the single-site molecule ($`L=1`$), i.e., the usual Anderson impurity model . There are many approximation methods for the Hubbard model or the Anderson model (see Ref. for a recent review), and we will adopt a self-consistent mean-field approximation . We start with the strong-$`U`$ limit of $`H_d`$ in (2), i.e., the $`tJ`$ Hamiltonian $$H_d\underset{ij}{}t_{ij}d_{i\sigma }^{}d_{j\sigma }\frac{1}{U}\underset{ijk}{}t_{ij}(d_i^{}d_j^{}d_i^{}d_j^{})t_{jk}(d_jd_kd_jd_k)$$ within the reduced Hilbert space without doubly occupied sites. The second term in (2) can be made quadratic in $`d_{j\sigma }(d_{j\sigma }^{})`$ by means of the Hubbard-Stratonovich transformation introducing the auxiliary field $`\mathrm{\Delta }_j`$ to get the molecular part (corresponding to $`H_d`$) in the Euclidean action (4) $$S_d^E=S_\mathrm{\Delta }^E+_0^\beta 𝑑\tau \underset{ij}{}\widehat{d}_i^{}\left[\delta _{ij}_\tau +t_{ij}\widehat{\tau }_3\widehat{\mathrm{\Delta }}_{ij}\right]\widehat{d}_j,$$ (12) where $`\tau _3`$ is the Pauli matrix and $`S_\mathrm{\Delta }^E=_0^\beta 𝑑\tau U_j\|\mathrm{\Delta }_j\|^2`$. In Eq. (12), we have introduced a two-component spinor representation $$\widehat{d}_j\left[\begin{array}{c}d_j\\ d_j^{}\end{array}\right],\widehat{\mathrm{\Delta }}_{ij}t_{ij}\left[\begin{array}{cc}0& \mathrm{\Delta }_i+\mathrm{\Delta }_j\\ \mathrm{\Delta }_i^{}+\mathrm{\Delta }_j^{}& 0\end{array}\right].$$ (13) In this representation, the Green’s function in Eq. (5) also has a matrix form $`\widehat{𝒢}(i,j;\tau )=T_\tau \widehat{d}_i(0)\widehat{d}_j^{}(\tau )`$. After integrating out the fields $`c_{𝐤\sigma }`$ and $`c_{𝐤\sigma }^{}`$, the Euclidean action can be written as $$S^E=S_\mathrm{\Delta }^E_0^\beta 𝑑\tau 𝑑\tau ^{}\widehat{d}_i^{}(\tau )\widehat{𝒢}^1(i,j;\tau \tau ^{})\widehat{d}_j(\tau ^{})$$ (14) where the Green’s function $`𝒢`$ is given by the Dyson equation $$\widehat{𝒢}^1(j,k;i\omega _n)=\widehat{𝒢}_0^1(j,k;i\omega _n)\widehat{\mathrm{\Sigma }}_c(j,k;i\omega _n)\widehat{\mathrm{\Sigma }}(j,k;i\omega _n),$$ (15) with the new self-energy term $`\widehat{\mathrm{\Sigma }}(i,j;\tau )=\delta (\tau )\widehat{\mathrm{\Delta }}_{ij}(\tau )`$. Here $`𝒢_0`$ in (8) and $`\mathrm{\Sigma }_c`$ in Eq. (9) have been extended to matrix form in a trivial way: $`\widehat{𝒢}_0=𝒢_0\widehat{\tau }_0`$ and $`\widehat{\mathrm{\Sigma }}_c=\mathrm{\Sigma }_c\widehat{\tau }_0`$, where $`\tau _0`$ is the identity matrix. So far no approximation has been made and the formal expression for the Euclidean action in Eq. (14) is exact. The interaction effect is correctly incorporated through the self-energy term $`\widehat{\mathrm{\Sigma }}`$ (whereas the effect of the metal is again manifested by $`\widehat{\mathrm{\Sigma }}_c`$) as long as the integration over the field $`\mathrm{\Delta }_j`$ is performed properly. At this point, we make our main approximation and neglect the fluctuations in the field $`\mathrm{\Delta }_j`$. We first assume a particular realization of the field $`\mathrm{\Delta }_j`$ in the Dyson equation (15) and then determine $`\mathrm{\Delta }_j`$ within the stationary-phase approximation. Iterating this procedure allows us to determine $`\mathrm{\Delta }_j`$ and the Green’s function self-consistently. Moreover, the self-consistency equation $$\mathrm{\Delta }_i=\frac{1}{U}\underset{j}{}t_{ij}d_id_jd_id_j$$ (16) corresponding to this procedure, gives us a physical interpretation of the field $`\mathrm{\Delta }_j`$: $`\mathrm{\Delta }_j`$ measures the spin-singlet correlation between the $`j`$th site and its neighboring sites. It is clear that in the limit $`U\mathrm{}`$, $`\mathrm{\Delta }_j2t/U`$ unless $`j`$ is too close to the boundaries. Finally, the LDOS can be obtained from the diagonal components of the retarded Green’s function $`\widehat{G}^R(j,k;E)=\widehat{𝒢}(j,k;i\omega _nE+i0^+)`$: $$\rho (j,E)=\frac{1}{\pi }\mathrm{Im}\left[\widehat{G}_{11}^R(j,j;E)+\widehat{G}_{22}^R(j,j;E)\right].$$ (17) ## 3 Results We will now report specific results on the LDOS of the molecular wire and discuss their physical implications. For the density of states (DOS) in the metal, we adopt the model functional form $`N_c(E)=N_c(0)(1E^2/W_c^2)`$. The parameters $`N_c(0)`$ and $`W_c`$ specify the DOS at the Fermi energy and the width of the conduction band, respectively. The self-energy contribution (9) from the coupling to the metal is then given by $$\mathrm{\Sigma }_c(j,k;E)=\delta _{j1}\delta _{k1}V^2\left[\alpha _c(E)i\pi N_c(E)\right]$$ (18) with $`\alpha _c(E)`$ defined by $$\alpha _c(E)=N_c(0)\left[\frac{2E}{W_c}+\left(1\frac{E^2}{W_c^2}\right)\mathrm{ln}\left\|\frac{E/W_c+1}{E/W_c1}\right\|\right].$$ (19) All results shown in this work will be generated using the values $`ϵ_d=4t`$, $`V=0.1t`$, $`W_c=40t`$, and $`N_c(0)=10/t`$. The other parameters will be specified as needed. We also remark that we only investigate the lower band, i.e., the band below the Mott-Hubbard gap, in the interacting case. In Fig. 2 we plot typical LDOS curves at three different sites ($`j=1`$, $`3`$, and $`5`$) on a molecule with size $`L=10`$ and on-site interaction $`U=10t`$. The structures get sharper as $`U`$ increases because tunneling is blocked, see the discussion below. At a given site, the LDOS decays rapidly outside the band. Figure 3 shows the spatial dependence of the LDOS along molecules with sizes $`L=10`$ and $`20`$, respectively, both for the non-interacting (open circle) and interacting (solid circle) case. The LDOS has been calculated at the fixed energy $`E=2t`$, slightly above the upper edge of the band (remember that at the parameter values chosen below Eq. (19), $`E=2t`$ is the upper edge of the band for an infinite tight-binding chain, i.e., $`H_d`$ with $`L=\mathrm{}`$ and $`U=0`$ in Eq. (2)). The LDOS decays exponentially with distance from the metal-molecule contact, except for the region close to the boundaries. This reproduces the exponential decay of the STM tunneling conductance in Ref. . It also agrees with the expectation that the electron wave function induced by the metal to the molecule should decay in an exponential way. Another interesting conclusion that can be drawn from Fig. 3 is the influence of the interaction on the suppression of the LDOS. The characteristic decay length decreases with the interaction $`U`$. This can be interpreted by saying that the strong repulsion by the electrons sitting on (singly-occupied) sites tends to block tunneling events from the metal to the molecule. This effectively leads to the suppression of the “hybridization” term $`H_T`$, Eq. (3). The suppression of the LDOS by interaction effects does not depend on the choice $`E=2t`$, however, the spatial scale would be different for different positions of the molecular level. The influence of repulsive interactions was also studied in , where it is shown that the suppression is lower for repulsive interactions than for a dimerization leading to the same value of the gap. In conclusion, we have calculated the local density of states of a molecule adsorbed to a metal surface using a new Anderson-type model. The LDOS decays along the molecule in an exponential way (like in the experiment ) and is suppressed by a local interaction on the molecule. Our formalism provides a qualitative understanding of this and similar experiments. It is able to include local interaction effects and can be used to treat more complicated geometries that will be investigated in future experiments. \*** We would like to acknowledge discussions with T. Jung and A. Baratoff and thank J. Gimzewski for providing a preprint of Ref. .
warning/0003/astro-ph0003105.html	ar5iv	text	# Quintessence-like Dark Matter in Spiral Galaxies ## Abstract Through the geodesic analysis of a static and axially symmetric space time, we present conditions on the state equation of an isotropic perfect fluid $`p=\omega d`$, when it is considered as dark matter in spiral galaxies. The main conclusion is that it can be an exotic fluid ($`1<\omega <1/3`$) as it is found for Quintessence at cosmological scale. There is no doubt about the importance of the mystery concerning the nature of dark matter in the Universe and in particular in galaxies. The consequences of observations made on SNIa supernovae have posed challenges to the available theoretical machinery, and certain models explaining such phenomena have arised which propose exotic types of matter and therefore unusual equations of state, such as a Cosmological Constant, Cold Dark Matter models, Dilaton Fields and Quintessence . However, at the galactic level there are no models consistent with the cosmological ones, and which give some light in the understanding of the nature of dark matter. In order to be precise about the problem let us recall the situation of the galactic dark matter, for which we confine ourselves to the observations made by Rubin et al. who found that for a few sample of spiral galaxies the interstellar gas and stars lying far away from the center (in the equatorial plane) of the corresponding galaxy behaves in a non Kepplerian way, but their circular velocity seems to be independent of the radius starting from a certain distance to the galactic center, i.e. the rotation curve profile of a spiral galaxy is flat outside a central galactic region. It was then inferred a distribution $`1/r^2`$ of non luminous matter (dark matter) which should contribute to the flatness of the rotation curves. There exists certain controversy about the flatness of such curves , but in general it is accepted that rotation curves are flat up to the precision of the meassurementes made by the astronomers and that this behavior is reproduced even for large samples of spiral galaxies . The most accepted scenario for a spiral galaxy reads as follows: it is an object composed by a luminous disc whose density distributed in an exponentially decay which conspires with a dark halo whose density is distributed as $`1/r^2`$ . In this way it is found an explanation for the kinetic behavior of gas and stars composing a spiral galaxy, but how was this mixture formed and what inspired nature to conspire in this way and not another one? If the dark matter is baryonic such as MACHOS for instance, why does its density have a non exponential distribution as luminous matter density does? If it is non baryonic, what is it made of, or at least which is its equation of state? This last question is the one that occupies ourselves in the present work. In this letter we proceed in the following way: Assume that a spiral galaxy lies on a background axysimmetric static space time, which is characterized by the presence of a perfect fluid with an arbitrary equation of state, i.e. $`p=\omega d`$ being $`\omega `$ a free function, and then we find conditions over $`\omega `$ that permit flat rotation curves of test particles. Other types of candidates to dark matter are discused in . First of all it must be clear that our treatment is valid only in the dark matter dominated region, i.e. where the rotation curves are flat, and we do not consider the galactic core region. Observational data show that the galaxies must be composed by almost 90% of dark matter, distributed at the halo, in order to explain the observed dynamics of particles in the luminous sector of the galaxy. We can thus assume that luminous matter does not contribute in a very important way to the total energy density of the halo of the galaxy in the mentioned region. On the other hand, the exact symmetry of the halo is still unknown, but it is reasonable to suppose that the halo is symmetric with respect to the rotation axis of the galaxy, so we choose the space time to be axial symmetric. Furthermore, the rotation of the galaxy does not have a large effect on the motion of test particles around it; dragging effects in the halo of the galaxy can be considered too small to seriously affect the motion of tests particles (stars) traveling around the galaxy. The circular velocity of stars (like the Sun) is of the order of 230 Km/s, much less than the speed of light. Hence, in the region of interest we can suppose the space-time to be static as well. Therefore, we start by considering the background described by the following line element: $$ds^2=e^{2\psi }dt^2+e^{2\psi }[e^{2\gamma }(d\rho ^2+dz^2)+\mu ^2d\phi ^2],$$ (1) which corresponds to an static axially symmetric space-time; the coordinates are the usual cylindrical ones. We recall the reader that observations are made upon objects lying in the galactic equatorial plane, thus the Lagrangian for a test particle travelling on such slide of the space time described by (1) is $$2=e^{2\psi }\dot{t}^2+e^{2\psi }[e^{2\gamma }\dot{\rho }^2+\mu ^2\dot{\phi }^2],$$ (2) where dot means derivative with respect to the proper time $`\tau `$ of the test particle. The radial geodesic motion equation is then $$\dot{\rho }^2e^{2(\psi \gamma )}\left[E^2e^{2\psi }L^2\frac{e^{2\psi }}{\mu ^2}1\right]=0.$$ (3) where $`E`$, and $`L`$, are constants associated with this geodesic motion along the equatorial plane. We are interested in circular and stable motion of test particles, therefore the following conditions must be satisfied i) $`\dot{\rho }=0`$, circular trajectories ii)$`\frac{V(\rho )}{\rho }=0`$, extreme ones iii)$`\frac{^2V(\rho )}{\rho ^2}\|_{extr}>0`$, and stable. being $`V(\rho )=e^{2(\psi \gamma )}\left[E^2e^{2\psi }L^2e^{2\psi }/\mu ^21\right]`$. Recalling that $`E`$ and $`L`$ are constants of motion for each circular orbit, it is straightforward to obtain expressions for the energy $`E`$, angular momentum $`L`$, angular velocity $`\mathrm{\Omega }=d\phi /dt`$ and the tangential velocity $`v^{(\phi )}=e^{2\psi }\mu \mathrm{\Omega }`$ , corresponding to a circular, stable equatorial motion: $`E`$ $`=`$ $`e^\psi \sqrt{{\displaystyle \frac{\frac{\mu _{,\rho }}{\mu }\psi _{,\rho }}{\frac{\mu _{,\rho }}{\mu }2\psi _{,\rho }}}},`$ (4) $`L`$ $`=`$ $`\mu e^\psi \sqrt{{\displaystyle \frac{\psi _{,\rho }}{\frac{\mu _{,\rho }}{\mu }2\psi _{,\rho }}}},`$ (5) $`\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \frac{e^{2\psi }}{\mu }}\sqrt{{\displaystyle \frac{\psi _{,\rho }}{\frac{\mu _{,\rho }}{\mu }\psi _{,\rho }}}},`$ (6) $`v^{(\phi )}`$ $`=`$ $`\sqrt{{\displaystyle \frac{\psi _{,\rho }}{\frac{\mu _{,\rho }}{\mu }\psi _{,\rho }}}},`$ (7) and for the stability condition: $$V_{,\rho \rho }\|_{extr}=\frac{2e^{2(\psi \gamma )}}{\frac{\mu _{,\rho }}{\mu }2\psi _{,\rho }}\left(\frac{\mu _{,\rho }}{\mu }\psi _{,\rho \rho }\frac{\mu _{,\rho \rho }}{\mu }\psi _{,\rho }+4\psi _{,\rho }^{}{}_{}{}^{3}6\frac{\mu _{,\rho }}{\mu }\psi _{,\rho }^{}{}_{}{}^{2}+3\left(\frac{\mu _{,\rho }}{\mu }\right)^2\psi _{,\rho }\right)>0$$ (8) where a coma stands for partial derivative. Now, observe that if the functions $`\psi `$ and $`\mu `$ are related by $$e^\psi =(\frac{\mu }{\mu _0})^l.$$ (9) being $`l=const,`$ we obtain a necessary and sufficient condition for the velocity $`v_{c}^{}{}_{}{}^{(\phi )}`$ to be the same for two orbits at different radii, given $`l=(v_{c}^{}{}_{}{}^{(\phi )})^2/\left(1+(v_{c}^{}{}_{}{}^{(\phi )})^2\right),`$ and equation (8) tells us that this motion is stable. We call equation (9) together with such value of $`l`$ the flat curve condition. We now write the Einstein’s equations $`G_{\alpha \beta }=8\pi T_{\alpha \beta }`$ for an arbitrary energy momentum tensor for the line element (1): $`\mu D^2\psi +D\mu D\psi `$ $`=`$ $`4\pi \mu [e^{2(\psi \gamma )}(e^{2\psi }T_{tt}+{\displaystyle \frac{e^{2\psi }}{\mu ^2}}T_{\phi \phi })+T_{\rho \rho }+T_{zz}],`$ (10) $`D^2\mu `$ $`=`$ $`8\pi \mu [T_{\rho \rho }+T_{zz}]`$ (11) $`\gamma _\rho \mu _\rho \gamma _z\mu _z\mu (\psi _{\rho }^{}{}_{}{}^{2}\psi _{z}^{}{}_{}{}^{2})+\mu _{zz}`$ $`=`$ $`8\pi \mu T_{\rho \rho },`$ (12) $`\gamma _\rho \mu _z+\gamma _z\mu _\rho 2\mu \psi _\rho \psi _z\mu _{\rho z}`$ $`=`$ $`8\pi \mu T_{\rho z}.`$ (13) where we have introduced the operator $`D=(_\rho ,_z)`$, see reference . In order to have flat tangential curve velocities, it is introduced the flat curve condition (9). This condition is valid on the equatorial plane. Nevertheless, the halo is expected to be almost spherically symmetric, that means that if we know the functional dependence of the gravitational potential on the equatorial plane, this dependence should be the same one in almost the rest of the halo. In that case it is reasonable to suppose that the flat curve condition (9) is valid in a region around the equatorial plane. Thus, in this region we substitute the relations (9) into the left hand side of equation (10) obtaining $`\mu D^2\psi +D\mu D\psi =lD^2\mu `$ and with (11) we get a constrain equation amount the components of the stress energy tensor: $$\left(\frac{1(v_{c}^{}{}_{}{}^{(\phi )})^2}{1+(v_{c}^{}{}_{}{}^{(\phi )})^2}\right)(T_{\rho \rho }+T_{zz})=e^{2(\psi \gamma )}\left(e^{2\psi }T_{tt}+\frac{e^{2\psi }}{\mu ^2}T_{\phi \phi }\right)$$ (14) Notice that this relation must be satisfied by any stress energy tensor which, within the approximation made in the analysis, curves the space time in such a way that the motion of test particles corresponds to the observed one. Let us consider the case of a stress energy tensor corresponding to a perfect fluid, $`T_{\mu \nu }=(d+p)u_\mu u_\nu +g_{\mu \nu }p`$, with $`d`$ the density of the fluid and $`p`$ its pressure. In this case we are thinking on a “dark fluid”, which is not seen but it is thought that it could be there affecting the geometry in the way needed in order to have the observed behavior in the tangential velocities of the luminous matter, as just mentioned. Considering the dark fluid as static, the four velocity of such dark fluid is given by $`u^\alpha =(u^0,0,0,0)`$ which, for the line element (1) reads: $`u^0=Ee^{2\psi }`$, thus $`u_0=E`$ and from $`u^\mu u_\mu =1`$, we obtain that $`E=e^{2\psi }`$. Therefore, the stress energy tensor has the form: $`T_{tt}=e^{2\psi }d,`$ (15) $`T_{\rho \rho }=T_{zz}=e^{2(\psi \gamma )}p,`$ (16) $`T_{\phi \phi }=\mu ^2e^{2\psi }p.`$ (17) Substituting these expressions into (14), we obtain that in the equatorial plane, in order to satisfy the observed behavior on the tangential velocities, the “dark fluid” has to fulfill the relation: $$2\left(\frac{1(v_{c}^{}{}_{}{}^{(\phi )})^2}{1+(v_{c}^{}{}_{}{}^{(\phi )})^2}\right)p=(d+p)$$ (18) Let us see which are the permitted relations between pressure and density of the perfect fluid providing flat rotation curves, we thus obtain: $$p=\frac{1+(v_{c}^{}{}_{}{}^{(\phi )})^2}{3(v_{c}^{}{}_{}{}^{(\phi )})^2}d$$ (19) relation from which the $`d`$ coefficient is identified with the square velocity dispersion of the dark particles, that appears to be negative. We are now in a convenient position to strict the state equation. As the velocities of the gas and stars rotating in the flat region must be within $`0<v_c^{(\phi )2}<1`$, (the observed ones are of the order of $`v_c^{(\phi )}10^3`$ ), relation (19) implies $`1<\omega <1/3`$, being $`p=\omega d`$. This result coincides with the one obtained at cosmological scale for the respective equation of state in the Quintessence model . Therefore the analysis presented in this letter, gives support to the hypothesis that a Quintessence-like equation of state could be the solution for the dark matter problem at galactic scale. In both cases it turns out the need of an exotic equation of state, with $`\omega =0.33`$ at a galactic scale and $`\omega =0.64`$ for the cosmos . In any case, we have shown that galactic dark matter satisfying an exotic equation of state certainly can be used to explain the observed behavior on the rotational curves of spiral galaxies. We want to thank Daniel Sudarsky, and Alejandro Corichi for many helpful discussions. We also want to express our acknowledgment to the relativity group in Jena for its kind hospitality and partial support. This work is also partly supported by CONACyT México, 94890 (Guzmán) and by the DGAPA-UNAM IN121298 (Núñez, Ramírez).
warning/0003/quant-ph0003083.html	ar5iv	text	# Nonlinear gauge interactions - A solution to the “measurement problem” in quantum mechanics? ## Abstract We propose that the mechanism responsible for the “collapse of the wave function” (or “decoherence” in its broadest meaning) in quantum mechanics is the nonlinearities already present in the theory via nonabelian gauge interactions. Unlike all other models of spontaneous collapse, our proposal is, to the best of our knowledge, the only one which does not introduce any new elements into the theory. Indeed, unless the gauge interaction nonlinearities are not used for exactly this purpose, one must then explain why the violation of the superposition principle which they introduce does not destroy quantum mechanics. A possible experimental test of the model would be to compare the coherence lengths for, e.g., electrons and photons in a double-slit experiment. The electrons should have a finite coherence length, while photons should have a much longer (in principle infinite) coherence length. PACS numbers: 03.65.-w, 03.65.Bz, 11.15.-q We start by noting an apparent paradox in the presently most fundamental description of (experimentally tested) physical reality. * 1) The absolute backbone of quantum mechanics is the superposition principle (e.g., interfering amplitudes, summation of Feynman diagrams, etc). It is also well known that superposition requires linear equations (i.e., the sum of two different solutions to a nonlinear equation is generally not a solution, ruining superposition). The Hilbert space of quantum mechanics and the Fock space of quantum field theory are linear spaces (based on the superposition requirement), suitable for linear mappings or operators. * 2) Nonabelian gauge field theories describing the fundamental interactions obey nonlinear evolution equations in the gauge fields. This is in apparent contradiction to point 1). For convenience, we write down the evolution equations for pure Yang-Mills fields below. Although the fermion evolution obeys linear equations, they become “contaminated” by nonlinearities through the interaction. The nonabelian vector gauge fields are governed by a set of coupled, second order, nonlinear PDEs on Minkowski spacetime. (The general argument for gravity is the same, but involve tensor fields on a dynamical spacetime. We do not explicitly write down those equations.) For pure Yang-Mills the evolution equations are given by the following formula, $$(^\mu g[A^\mu ,])_a^b(_\mu A_\nu _\nu A_\mu g[A_\mu ,A_\nu ])_b=0,$$ (1) where $`g`$ is the coupling constant and $`a`$, $`b`$ are indices of the gauge group (i.e., $`a,b1,2,3`$ for $`SU(2)`$ and $`a,b1,\mathrm{},8`$ for $`SU(3)`$). Summation over repeated indices is implied. The operator (“covariant derivative”) at the left works according to $`(^\mu g[A^\mu ,])(anything)=^\mu (anything)g[A^\mu ,anything]`$. We see that we get highly nonlinear (quadratic and cubic) terms in the gauge fields, especially when the coupling constant, $`g`$, is large. The commutator terms (square brackets) vanish identically for abelian fields (e.g. photons) because the gauge fields then commute, leaving only the ordinary, linear Maxwell equations. In the Feynman path-integral formulation of quantum mechanics the nonlinearities can be “hidden” in the action functional, but as the Schrödinger, Heisenberg and path-integral formulations are equivalent, a problem in one of them must translate into a problem in all formulations. Instead of trying to reconcile the two apparently contradictory statements above (by, for instance, modifying the rules of quantum mechanics), introduced by nonabelian gauge fields, we instead propose to turn a vice into a virtue by postulating that it is the dynamical nonlinear interaction terms which break the superposition of different quantum states of a system, i.e. acting as the physical mechanism which reduces the state vector, or “collapses the wave function” in the less general Schrödinger setting. We thus get a self-induced collapse - “SIC”, into the ordinary world of chairs, tables, people and indeed also recorded elementary particle tracks in a photographic emulsion, a bubble chamber, or a modern multi-purpose computer-aided detector. That a quantum mechanical state must be able to “self-decohere” is imperative in quantum cosmology, the quantum mechanical treatment of the whole universe, where no “outside” observer exists. The self-induced collapse puts an end to the infinite regress of quantum superposition, where first the measuring apparatus obtains a quantum mechanical nature, then the observer, and so on, ad infinitum, until the whole universe consists of infinitely many superimposed quantum states, without any one of them actually being “realized”. The “many worlds” interpretation of Everett purports to solve this problem by assuming that we only see events which take place in one of these branching universes, but it seems that the fundamental question of when, and how, the universe actually branches is unanswered by that model (this being the equivalent of the “measurement problem” in the orthodox interpretation). We now turn to the actual implementation of our idea of self-induced collapse. For simplicity, we choose the following (non-covariant) expression for the (average) self-decoherence time. $$\tau =\frac{\mathrm{}}{E_{N.L.}},$$ (2) where $`E_{N.L.}`$ is the energy stored in the nonlinear field configuration of the nonabelian interaction (which in turn depends on the strength of the coupling). Observe that the relation is not an uncertainty relation, despite its identical form, as $`\tau `$ and $`E_{N.L.}`$ are not uncertainties. As we want the energy of the full nonlinear theory, and we cannot today explicitly calculate this inherently non-perturbative quantity, we take the energy to be a characteristic energy for the interaction. If, for instance, for QCD, we as a rough approximation take the energy to be $`E_{N.L.}\mathrm{\Lambda }_{QCD}0.2`$ GeV, we get as a rough “ballpark” figure $`\tau _{QCD}10^{23}`$s for the decoherence time for strong QCD (e.g., inside a non-disturbed hadron). Although the exact result probably will differ by many orders of magnitude, this may explain why (semi-)classical models work so well for strong QCD, as the stronger the interaction is, the more “classical” it behaves according to our mechanism. In QCD the energy stored in gauge fields decreases as the absolute energy of the interaction increases, due to asymptotic freedom. In our model, the fields are the fundamental entities which obey quantum mechanics, the particle aspect appearing each time a self-collapse takes place. That is, the quantum mechanical (linear, unitary) evolution is constantly punctuated by (possibly random) “hits” of self-collapse at an average frequency of $`\tau ^1`$. This is similar to the case in orthodox quantum mechanics where an observation (or the initial preparation of a state) suddenly “realizes” one of the potential outcomes, after which the unitary (linear) evolution of the state takes over until the next observation. A “macroscopic” piece of matter has such a high energy stored in nonlinear field configurations that $`\tau =\frac{\mathrm{}}{E_{N.L.}}0`$, approximating a continuously collapsing state, i.e., a classical state. The model thus forbids quantum mechanical effects to “invade” the macroscopic world, and hence resolves the “Schrödinger’s cat” paradox and related questions such as “Wigner’s friend” , etc. Note that any significant nonlinear interaction, whether as part of a “measurement” carried out by conscious beings or not, bring about the decoherence of interfering amplitudes into (semi-)classical states. Conscious observation is therefore only a special case of the more general nonlinearity, as all “measuring apparatuses” (including human beings!) consist of both weakly (all particles) and strongly (quarks) nonlinearly interacting constituents. Hence, there should be no need to introduce the mind into the interpretation of quantum mechanics at a fundamental level. For pure QED the nonlinear terms are absent, hence a hypothetical world built by QED alone would never be classical. It also explains why, e.g., atomic physics works so closely to orthodox quantum mechanics, as it is being “classicalized” only by (very) weak interaction effects. Were it not for the existence of the other interactions besides QED, we would indeed have quantum mechanical superpositions of whole universes, i.e., the “many worlds” interpretation of quantum mechanics by Everett . The difference between our proposal for self-induced collapse, and other models aiming at the same goal, is that, as far as we know, all other models postulate additional equations and/or variables, * Decohering histories : new fundamental principle of irreversible coarse graining + additional constraints to remove “too many” decoherent histories * Altered Schrödinger equation: obvious extra (non-unitary or nonlinear) term in Schrödinger equation * Bohm QM : additional (nonlinear) evolution equation for objective positions whereas we use only nonlinearities which are already present in the dynamics of the accepted standard model of particle physics. Another difference is that, to our knowledge, all other models for spontaneous collapse are non-relativistic, whereas our scheme is based on covariant theories. Also, it must be stressed that if the nonlinearities introduced by the nonabelian gauge fields are not used to explain the decoherence to (quasi)classical behaviour, it must instead be explained how they can be reconciled with the superposition principle of quantum mechanics. It is well known that a nonlinear mapping is non-reversible, as there in that case does not exist a unique inverse mapping (i.e., the mapping is not one-to-one). We therefore propose that the nonlinear gauge interaction is the physical “mechanism” of the “irreversible amplification” emphasized by Bohr as being necessary to produce classical, observable results from the quantum mechanical formalism. Even though Bohr himself denounced the need, or even the possibility, to give a physical description of this “mechanism” , we believe that the central problem for truly understanding quantum mechanics lies in the quantum measurement problem. For instance, it is only there, in the collapse of the wave function, that the undeterminacy of quantum mechanics enters. It may even be possible, if not entirely likely, that deterministic chaos in the nonlinear self-interaction can be responsible for the seemingly statistical character of quantum mechanics. Our model can be experimentally tested, at least in principle, as differently charged (electric, weak isospin, color,…) particles should have different coherence lengths. In a double-slit experiment, for instance, the photon should have a much longer (in principle infinite) coherence length than, e.g., electrons which ought to have a finite coherence length due to nonlinear weak interactions. As the full nonlinear calculations are very complicated, it is not possible to quantitatively predict the coherence lengths at the present time, but if it turns out that electrons experimentally have shorter coherence lengths than photons it would strengthen our hypothesis. Our model could also have importance for (the not yet existing theory of) quantum gravity. Weak gravity would have extremely long decoherence times, completely swamped by the other interactions. However, exactly where quantum gravity is expected to become important (i.e., at the Planck mass/energy scale), we see that it spontaneously decoheres. Hence a strong quantum gravity might not exist in this scheme. The collapse postulated in orthodox quantum mechanics is not relativistically covariant, as it is instantaneous (over all space), which is not a covariant concept. Only the (deterministic) unitary development of the state is taken into account by the relativistic Dirac equation and, more generally, by quantum field theory. As our scheme for collapse is based on covariant gauge-field theories, it might be possible to describe the collapse of the state in a covariant way, although our present attempt, which for simplicity singles out the energy accumulated in the nonlinear field configurations, is not covariant. On the other hand, it might be that the collapse should be described by an inherently non-local mechanism, as it seems that quantum mechanics at its very foundation is non-local, as given by the results of Aspect et al. , and more recent experiments on quantum non-separability. As a nonlinear theory also in some cases can be non-local, it would be interesting to investigate if this model of spontaneous collapse can account for such non-local effects. Work in this direction is in progress , together with a detailed investigation of the nonlinear terms in the nonabelian evolution equations, as a means to better understand the quantitative details of the proposed mechanism for self-collapse. In conclusion, we have emphasized that automatic dynamical collapse of the wave function in quantum mechanics may already be implicit in the existing dynamical theory of nonabelian (i.e., nonlinear) gauge fields. These include the weak interaction, QCD, gravity, and any other nonabelian fields which eventually may be found in the future. The nonlinear self-interaction terms break the fundamental superposition principle of quantum mechanics, hence making it plausible that they can be just the right physical mechanism for the purpose of collapse.
warning/0003/gr-qc0003038.html	ar5iv	text	# A quasilocal calculation of tidal heating ## I Introduction In many physical problems in gravitation, one is interested in the interaction of a nearly isolated gravitating system with an external universe. The interaction effects are computed in a “buffer zone” (see Sec. 20.6 of Ref. and Sec. IB of Ref. ) surrounding the gravitating system, in which the radius of curvature, scale of inhomogeneity, and rate of change of curvature are much smaller than the size of the body. The formalism of Thorne and Hartle and Zhang has been used recently by Purdue and Favata to compute the gauge-invariant heating of a body interacting with an external tidal field. Until now, calculations of the sort described in these references have made use of pseudotensors to compute energy and momentum fluxes. However, quasilocal methods should be equally applicable in situations with a reasonably well defined buffer zone—in this case, the quasilocal surface can be conveniently located in the buffer zone. While quasilocal methods are not fundamentally different than pseudotensor methods , an advantage of quasilocal method is that all quantities (e.g., energy fluxes) can be computed in terms of real tensors on the quasilocal surface. Gauge ambiguities in the total amount of energy and energy flux such as those reported in Ref. and discussed in Ref. still exist for the quasilocal methods, but now the ambiguities can be understood in terms of distortions of the quasilocal surface and so their geometric origin is identified. In this paper we present a quasilocal formalism for computing the work done on a gravitating system by an external universe. Our formalism is based on the quasilocal mass of Brown and York —the on-shell value of the gravitational Hamiltonian—which coincides with the Arnowitt-Deser-Misner energy at spatial infinity and the Trautman-Bondi-Sachs energy at null infinity . It is complementary to but independent of which studied how motion of the observers affects the Brown-York energy. We use our expression for energy flux to compute (i) the energy lost in gravitational radiation from a gravitational system and (ii) the heating of a body through interactions with an external tidal field. Problem (i) demonstrates that the formula for the work reproduces the known gravitational radiation flux formula when the quasilocal surface is located in the wave-zone. Problem (ii) reproduces the calculation of Purdue using quasilocal methods and shows how these methods are applicable for problems in which the quasilocal surface is located in a buffer zone. ## II Quasilocal energy flux In this section, we derive an expression for the energy flux through a closed two-surface surrounding a gravitating system. Our analysis closely follows Sec. V of Ref. , which derives a conserved measure of mass for stationary systems. We relax the requirement that the quasilocal two-surface time evolution vector be a Killing vector of the spacetime and thereby obtain an expression for the rate of change in the mass of the system. Consider a gravitating system separated from the external universe by a ($`2+1`$)-dimensional timelike boundary $`B`$. This boundary has an outward “radial” normal vector $`n^a`$, a metric $`\gamma _{ab}=g_{ab}n_an_b`$ induced by its embedding in the spacetime with metric $`g_{ab}`$, and an extrinsic curvature $`\mathrm{\Theta }_{ab}=\frac{1}{2}\text{£}_n\gamma _{ab}`$ (with trace $`\mathrm{\Theta }=\gamma ^{ab}\mathrm{\Theta }_{ab}`$). Let $`\mathrm{}_a`$ be the derivative operator compatible with the metric $`\gamma _{ab}`$. Foliate the boundary $`B`$ into closed two-surfaces $`\mathrm{\Omega }_t`$ of constant time $`t`$; then the time evolution vector $`t^a`$ on $`B`$ satisfies $`t^a\mathrm{}_at=1`$ and can be decomposed into a lapse function $`N`$ and a shift vector $`V^a`$ on $`\mathrm{\Omega }_t`$ via $`t^a=Nu^a+V^a`$, where $`u^a`$ is the timelike normal to $`\mathrm{\Omega }_t`$ embedded in $`B`$. The closed, spacelike, two-surface $`\mathrm{\Omega }_t`$ has an induced metric $`\sigma _{ab}=\gamma _{ab}+u_au_b`$ and, viewed as a two-surface embedded in a three-dimensional spacelike hypersurface $`\mathrm{\Sigma }`$ locally defined such that $`n^aT\mathrm{\Sigma }`$, the extrinsic curvature of $`\mathrm{\Omega }_t`$ is $`k_{ab}=\frac{1}{2}\text{£}_n\sigma _{ab}`$. A full discussion of the geometry of the boundary $`B`$ and its foliation (including a diagram) may be found in . The notation there is substantially the same as here though $`u^a`$ is written as $`\stackrel{~}{u}^a`$. The Codazzi identity, $$\mathrm{}_a\tau ^{ab}=\gamma ^{bc}n^dR_{cd}/8\pi ,$$ (1) where $`\tau ^{ab}=(\mathrm{\Theta }\gamma ^{ab}\mathrm{\Theta }^{ab})/8\pi `$, relates the extrinsic curvature of $`B`$ to the spacetime Ricci curvature $`R_{ab}`$. It then follows from the Einstein field equations that $$\mathrm{}_a(t_b\tau ^{ab})=t^an^bT_{ab}+\frac{1}{2}\tau ^{ab}\text{£}_t\gamma _{ab}.$$ (2) We restrict our attention to a vacuum spacetime in which the stress-energy tensor $`T_{ab}`$ vanishes. Then, if $`t^a`$ is a Killing vector field of the boundary metric $`\gamma _{ab}`$, Eq. (2) is a conservation equation and the quantity $$M=_{\mathrm{\Omega }_t}d^2x\sqrt{\sigma }u_at_b\tau ^{ab}$$ (3) is a conserved measure of the total mass contained within the boundary $`\mathrm{\Omega }_t`$. It is the “non-orthogonal” Brown-York mass , up to a subtraction term that is required for it to be bounded for large surfaces in asymptotically flat spacetimes (see, e.g., Lau or Mann ). When $`t^a`$ is not a Killing vector of the boundary, then Eq. (2) represents an energy flow from the system. Between two times $`t_1`$ and $`t_2`$ one can integrate to find that $`\mathrm{\Delta }M=\frac{1}{2}_Bd^3x\sqrt{\gamma }\tau ^{ab}\text{£}_t\gamma _{ab}`$ is the change in the mass contained by $`\mathrm{\Omega }_t`$. Subtraction terms from a reference spacetime do not need to be included here as it expresses the change in the mass of the system. The rate at which this work is done is $$\frac{dW}{dt}=\frac{1}{2}_{\mathrm{\Omega }_t}d^2x\sqrt{\gamma }\tau ^{ab}\text{£}_t\gamma _{ab}$$ (4) which describes the rate of change of the system’s mass due to the purely gravitational interaction between it and the surrounding environment. It is illustrative to decompose the expression for the work into terms involving projections of $`\text{£}_t\gamma _{ab}`$ normal to and into the spatial two-surfaces $`\mathrm{\Omega }_t`$. We find $$\frac{dW}{dt}=_{\mathrm{\Omega }_t}d^2x\sqrt{\sigma }\{\frac{1}{2}s^{ab}\text{£}_t\sigma _{ab}\epsilon \text{£}_tN+j_a\text{£}_tV^a\}$$ (5) where $`\epsilon =\sigma ^{ab}k_{ab}/8\pi `$, $`j_a=\sigma _{ab}u_c\mathrm{\Theta }^{bc}/8\pi `$, and $`s^{ab}=[k^{ab}+\sigma ^{ab}(n^cu^d\mathrm{}_du_c\sigma ^{cd}k_{cd})]/8\pi `$ are the quasilocal surface energy, momentum, and stress densities. The first two are potentials conjugate to changes in the lapse function and shift vector respectively while the surface stress density is a work potential conjugate to changes in the size and shape of the surface $`\mathrm{\Omega }_t`$. The stress density can be further decomposed as follows. A change in the two-metric, $`\delta \sigma _{ab}=\varsigma _{ab}\delta \sqrt{\sigma }+\sqrt{\sigma }\delta \varsigma _{ab}`$, is written as a change in the “size” $`\sqrt{\sigma }`$ of the surface plus a change in the conformally-invariant part of the metric (the “shape” of the surface) $`\varsigma _{ab}=\sigma _{ab}/\sqrt{\sigma }`$. Correspondingly, the surface stress density is decomposed into a surface tension $`s=s^{ab}\varsigma _{ab}`$ and a shear $`\eta ^{ab}=s^{ab}/\sqrt{\sigma }`$. Then we rewrite the work term as $`\frac{1}{2}s^{ab}\text{£}_t\sigma _{ab}=\frac{1}{2}(s\text{£}_t\sqrt{\sigma }+\eta ^{ab}\text{£}_t\varsigma _{ab})`$. The above has a particularly nice application in the physics of thin shells. Israel first showed that a thin shell of matter can be described in general relativity by matching two spacetimes along a timelike boundary $`B`$ such that even though they induce the same surface metric on $`B`$, the extrinsic curvature in each spacetime is different. If $`\mathrm{\Theta }_{ab}^+`$ and $`\mathrm{\Theta }_{ab}^{}`$ are those curvatures this (mild) singularity can be accounted for if there is a (distributional) stress energy tensor $`S_{ab}=\tau _{ab}^+\tau _{ab}^{}`$ over $`B`$. A set of observers dwelling on the surfaces $`\mathrm{\Omega }_t`$ (which foliate $`B`$) measures the shell to have matter-energy $`M^+M^{}`$ \[Eq. (3)\]. A more detailed discussion of this may be found in but here we note that the above analysis for the quasilocal energy also shows that the set of observers dwelling on $`\mathrm{\Omega }_t`$ measures the matter-energy to change with rate $`dW/dt`$ (Eq. 5). Then the quasilocal densities defined above are the energy, angular momentum, and stress tensor of the matter shell. A set of observers being evolved by $`t^a=u^a`$ see work being done on the shell at a rate equal to the integral of the stress tensor contracted with the time rate of change of the area—exactly as one would expect from classical physics. ## III Gravitational radiation Equation (4) purportedly measures the change in the mass of a system. In this section we apply our work formula to obtain the correct mass loss for a system radiating gravitational waves. For this we suppose that the quasilocal surface is in the wave-zone, far away from the radiating system. Although this is not a very interesting application of a quasilocal method (since an asymptotic method, such as the Bondi-Sachs mass loss formula could as well be used), it is useful to confirm that Eq. (4) does recover the correct result. Gravitational radiation far from the generating source can be described as a transverse-traceless perturbation to the flat-space metric. In spherical-polar coordinates, the metric is given by $`ds^2=dt^2+dr^2+(rd\theta )^2+(r\mathrm{sin}\theta d\varphi )^2+h_{\mu \nu }dx^\mu dx^\nu `$ where $`h_{\mu \nu }dx^\mu dx^\nu =h_+[(rd\theta )^2(r\mathrm{sin}\theta d\varphi )^2]+2h_\times (rd\theta )(r\mathrm{sin}\theta d\varphi )`$ is the transverse, trace free perturbation. The “plus,” $`h_+`$, and “cross,” $`h_\times `$, polarizations represent outgoing, spherical waves, and have the form $`h_+(t,r,\theta ,\varphi )=s_+(tr,\theta ,\varphi )/r`$ and $`h_\times (t,r,\theta ,\varphi )=s_\times (tr,\theta ,\varphi )/r`$. We then find the energy lost by the radiating system by inserting this metric into Eq. (4) while taking the boundary to be a sphere of constant $`r`$ in the wave-zone (very large $`r`$). The integrand of Eq. (4) is $$\frac{dE}{dtd\mathrm{\Omega }_t}=\frac{r^2}{16\pi }[(h_+/t)^2+(h_\times /t)^2]$$ (6) to leading order in the perturbation and in $`r`$. This is the standard expression for the flux of gravitational radiation—see, e.g., Eq. (10) of Ref. . By inspection of the form of the perturbation, it is clear that the energy loss arises due to the shearing of the bounding two-surface $`\mathrm{\Omega }_t`$ since, to leading order, the perturbation does not affect the volume element on that two-surface. Thus the entire energy loss (in the transverse, trace-free gauge) arises from the “$`\eta ^{ab}\text{£}_t\varsigma _{ab}`$” work term. As a simple example, consider two point-particles, each of mass $`m=M/2`$, orbiting each other in the $`xy`$-plane with angular frequency $`\omega `$ and constant separation $`a`$. The quadrupole moment tensor $`_{jk}`$ in Cartesian coordinates is $`_{xx}=_{yy}=\frac{1}{8}Ma^3\mathrm{cos}2\omega t`$ and $`_{xy}=\frac{1}{8}Ma^2\mathrm{sin}2\omega t`$ (constant terms omitted). The far-field metric perturbation is $`h_{jk}=2(^2_{jk}/t^2)/r`$, so $`h_{yy}=h_{xx}=(Ma^2\omega ^2/r)\mathrm{cos}2\omega (tr)`$ and $`h_{xy}=(Ma^2\omega ^2/r)\mathrm{sin}2\omega (tr)`$. Using Eqs. (4.3) and (4.4) of Ref. , we find $`h_+`$ $`=`$ $`{\displaystyle \frac{1}{2}}Ma^2\omega ^2r^1(1+\mathrm{cos}^2\theta )\mathrm{cos}2[\omega (tr)\varphi ]`$ (8) $`h_\times `$ $`=`$ $`Ma^2\omega ^2r^1\mathrm{cos}\theta \mathrm{sin}2[\omega (tr)\varphi ].`$ (9) We then integrate Eq. (6) over the sphere at large $`r`$ to obtain the loss of energy from the system: $$dE/dt=\frac{2}{5}M^2a^4\omega ^6=\frac{2}{5}(M/a)^5$$ (10) where we have used Kepler’s law $`a^3\omega ^2=M`$ for particles in a circular orbit. ## IV Tidal heating We now calculate the work done by an external gravitational field to deform a self-gravitating body. The canonical example of this effect in the solar system is the tidal heating of Io by Jupiter. In this instance, the gradient of Jupiter’s gravitational field distorts Io from being a perfect sphere and then tidally locks it in its orbit so that it always presents the same face to Jupiter. That orbit is strongly perturbed by the other Gallilean moons and so its radial distance from Jupiter varies with time. With this variation comes a corresponding one in the gradient of the field and so Io is gradually stretched and then allowed to relax. The energy transferred by this pumping is largely dispersed as heat and it is this heat that produces the volcanic activity on Io. The same type of process occurs for any two bodies in non-circular orbits about each other. First from a Newtonian perspective, we may mathematically describe the gravitational fields in this situation as follows. We assume that the self-gravitating body is far enough away from the source of the external field that that field is nearly uniform close to the body. Then in a rectangular coordinate system that orbits with the body with its origin at the center of mass, the Newtonian potential of the external field may be written as $`\mathrm{\Phi }_{\text{ext}}=\frac{1}{2}_{ij}x^ix^j`$ where $`_{ij}`$ is the (time-dependent but symmetric and trace-free) quadrupole moment of the field and $`x^i`$ is the position vector based at the body’s centre of mass. At the same time, to quadrupolar order the Newtonian potential of the body is $`\mathrm{\Phi }_\text{o}=M/r\frac{3}{2}r^3_{ij}n^in^j`$, where $`M`$ is the mass of the body, $`r`$ is the radial distance from the centre of mass, $`_{ij}`$ is its (time-dependent but symmetric and trace-free) quadrupole moment, and $`n^i`$ is the unit normal radial vector. With this in mind the techniques of Thorne and Hartle can be used to construct a metric that describes these situations in the slow moving, nearly Newtonian limit. First, define an annulus surrounding the body whose inner boundary is chosen so that the gravitational field of the body is weak throughout and whose outer boundary is chosen so that the external field is nearly uniform. This region is termed the buffer zone. The rectangular coordinate system is replaced with one that is chosen so that the metric is as close to Minkowskian as possible over the buffer zone . Then to first order in perturbations from Minkowski and first order in time derivatives the metric can be written as $`ds^2`$ $`=`$ $`(1+2\mathrm{\Phi })dt^2+2(A_j+_t\xi _j)dx^jdt`$ (12) $`+[(12\mathrm{\Phi })\delta _{ij}+_i\xi _j+_j\xi _i]dx^idx^j`$ where the indices run from one to three and $`\delta _{ij}`$ is the Cartesian metric $`\mathrm{diag}[1,1,1]`$ on a spacelike slice. The Newtonian potential is $`\mathrm{\Phi }=M/r\frac{1}{2}(3r^3_{ij}r^2_{ij})n^in^j`$ and $`A_j=2r^2n^kd_{jk}/dt\frac{2}{21}r^3(5n_jn^k2\delta _j^k)n^ld_{kl}/dt`$ is a vector potential that must be added so that the metric is a solution to the first order Einstein equations. Here, $`n^i`$ is the radial normal with respect to the flat spatial metric $`\delta _{ij}`$ and $`r^2=x^2+y^2+z^2`$. The diffeomorphism generating vector field $`\xi _j`$ represents the gauge ambiguity in setting up a nearly Minkowski coordinate system. In order that the metric be slowly evolving and nearly Minkowski, $`\xi _j`$ must be of the form $`\xi _j=\alpha r^2_{jk}n^k+\beta r^3_{jk}n^k+\gamma r^3_{kl}n^kn^ln_j`$, where $`\alpha `$, $`\beta `$, and $`\gamma `$ are free constants of order one. We set up a constant $`r`$ timelike quasilocal surface $`B`$ in the buffer zone and foliate with constant $`t`$ spacelike two-surfaces $`\mathrm{\Omega }_t`$. Then the time vector $`t^a`$ is $`/t`$. In calculating the rate of change of the mass contained within $`\mathrm{\Omega }_t`$ it is most convenient to switch to spherical coordinates. We make the standard transformation to spherical coordinates $`x^i=r[\mathrm{sin}\theta \mathrm{cos}\varphi ,\mathrm{sin}\theta \mathrm{sin}\varphi ,\mathrm{cos}\theta ]`$; in these coordinates, the metric is $`ds^2`$ $`=`$ $`(1+2\mathrm{\Phi })dt^2+(12\mathrm{\Phi })[dr^2+(rd\theta )^2`$ (17) $`+(r\mathrm{sin}\theta d\varphi )^2]+2\overline{A}_rdrdt+2\overline{A}_\theta (rd\theta )dt`$ $`+2\overline{A}_\varphi (r\mathrm{sin}\theta d\varphi )dt+H_{rr}dr^2+H_{\theta \theta }(rd\theta )^2`$ $`+H_{\varphi \varphi }(r\mathrm{sin}\theta d\varphi )^2+2H_{r\theta }dr(rd\theta )`$ $`+2H_{r\varphi }dr(r\mathrm{sin}\theta d\varphi )+2H_{\theta \varphi }(rd\theta )(r\mathrm{sin}\theta d\varphi )^2`$ where $`H_{rr}=4\alpha r^3_{rr}+6(\beta +\gamma )r^2_{rr}`$, $`H_{\theta \theta }=2\alpha r^3_{\theta \theta }+2\beta r^2_{\theta \theta }+2\gamma r^2_{rr}`$, $`H_{\varphi \varphi }=2\alpha r^3_{\varphi \varphi }+2\beta r^2_{\varphi \varphi }+2\gamma r^2_{rr}`$, $`H_{r\theta }=\alpha r^3_{r\theta }+(4\beta +2\gamma )r^2_{r\theta }`$, $`H_{r\varphi }=\alpha r^3_{r\varphi }+(4\beta +2\gamma )r^2_{r\varphi }`$, and $`H_{\theta \varphi }=2\alpha r^3_{\theta \varphi }+2\beta r^2_{\theta \varphi }`$. In these expressions $`_{rr}=_{ij}e_r^ie_r^j`$, $`_{r\theta }=_{ij}e_r^ie_\theta ^j`$, etc., with $`e_r^i=n^i`$, $`e_\theta ^i=_\theta e_r^i`$ and $`e_\varphi ^i=(1/\mathrm{sin}\theta )_\varphi e_r^i`$. Also, $`\overline{A}_r=(A_j+_t\xi _j)e_r^j`$, etc., but we don’t need their expanded forms since only time derivatives of them show up in later calculations and we are ignoring second order time derivatives. As might be expected, the subsequent calculations are quite involved and we did them partially with GRTensor . To lowest order $`{\displaystyle \frac{dW}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_t}}d^2x\sqrt{\gamma }\tau ^{ab}\text{£}_t\gamma _{ab}`$ (18) $`=`$ $`{\displaystyle \frac{1}{2}}_{ij}{\displaystyle \frac{d_{ij}}{dt}}+{\displaystyle \frac{1}{60}}{\displaystyle \frac{d}{dt}}[2(32\beta 2\beta ^2+4\gamma `$ (22) $`+4\gamma ^2+8\beta \gamma )r^5_{ij}_{ij}`$ $`+2(32\alpha +6\beta 12\gamma +8\alpha \gamma )_{ij}_{ij}`$ $`(9+12\alpha +4\alpha ^2)r^5_{ij}_{ij}].`$ The calculations used the identities $`_{\mathrm{\Omega }_t}𝑑\theta 𝑑\varphi \mathrm{sin}\theta A_{rr}B_{rr}=(8\pi /15)A_{ij}B_{ij}`$ and $`_{\mathrm{\Omega }_t}𝑑\theta 𝑑\varphi \mathrm{sin}\theta (2A_{\theta \varphi }B_{\theta \varphi }A_{\theta \theta }B_{\varphi \varphi }A_{\varphi \varphi }B_{\theta \theta })=(4\pi /3)A_{ij}B_{ij}`$ where the integrations are over the unit sphere. This result requires some interpretation. As the external field changes with time and thereby forces the self-gravitating body to change configuration, the work done by the external field can be split into time reversible and irreversible parts \[as seen in Eq. (18)\]. The reversible work represents work being done to increase the potential energy of the system and is recoverable. On the other hand the irreversible part represents work done to deform and/or heat up the system. This is the tidal heating that we are interested in. Further, from the quasilocal perspective, we expect to see an energy flow arising from fluctuations of the quasilocal surface within otherwise static fields. Of course this work would also be reversible. Thus, it is only the irreversible part that we are interested in and we have calculated that to be $`\frac{1}{2}_{ij}d_{ij}/dt`$ above. This is the same leading term obtained when one does the corresponding calculation in Newtonian gravity or with pseudotensors and it is independent of diffeomorphisms generated by $`\xi _j`$ which correspond to fluctuations of the quasilocal surface. Note however, that the time reversible and gauge dependent terms of equation (18) are dependent on those fluctuations and furthermore that dependence is different from that found in ref. using pseudo-tensor methods. Similarly other pseudo-tensor or quasilocal methods would obtain a different gauge dependence for these terms. What is important is that the physically relevant time irreversible term does not depend on the $`\xi _j`$-generated diffeomorphisms. Finally for completeness let us consider how this energy flow splits up into its components parts as considered in Eq. (5). Then to the order that we are interested the angular momentum term is zero and we are left with two terms $`dW_N/dt=𝑑\theta 𝑑\varphi \sqrt{\sigma }\epsilon \text{£}_tN`$ and $`dW_\sigma /dt=\frac{1}{2}𝑑\theta 𝑑\varphi \sqrt{\sigma }Ns^{ab}\text{£}_t\sigma _{ab}`$. We find $`{\displaystyle \frac{dW_N}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}_{ij}{\displaystyle \frac{d_{ij}}{dt}}+{\displaystyle \frac{\alpha }{15}}{\displaystyle \frac{d_{ij}}{dt}}_{ij}{\displaystyle \frac{\beta }{5}}_{ij}{\displaystyle \frac{d_{ij}}{dt}}{\displaystyle \frac{4\gamma }{5}}_{ij}{\displaystyle \frac{d_{ij}}{dt}}`$ (25) $`+{\displaystyle \frac{1}{60}}{\displaystyle \frac{d}{dt}}[2(4\gamma +\beta 2)r^5_{ij}_{ij}6_{ij}_{ij}`$ $`3(2\alpha 3)r^5_{ij}_{ij}].`$ The second term is a bit more complicated. It is $`{\displaystyle \frac{dW_\sigma }{dt}}`$ $`=`$ $`{\displaystyle \frac{\alpha }{15}}{\displaystyle \frac{d_{ij}}{dt}}_{ij}+{\displaystyle \frac{\beta }{5}}_{ij}{\displaystyle \frac{d_{ij}}{dt}}+{\displaystyle \frac{4\gamma }{5}}_{ij}{\displaystyle \frac{d_{ij}}{dt}}`$ (29) $`+{\displaystyle \frac{1}{30}}{\displaystyle \frac{d}{dt}}[(13\beta 2\beta ^2+4\gamma ^2+8\beta \gamma )r^5_{ij}_{ij}`$ $`+2(3\alpha +3\beta 6\gamma +4\alpha \gamma )_{ij}_{ij}`$ $`(2\alpha ^29\alpha +9)r^5_{ij}_{ij}].`$ Thus part of the work done is measured by deformations of the surface and part is measured by changes in how observers choose to measure the rate of passage of time. Note that individually the time irreversible sections of the two parts are gauge dependent, but when we combine them we reobtain Eq. (18) and the gauge dependence vanishes back into the reversible part where we would expect it. ## V Conclusions We have modified the quasilocal energy formalism of Brown and York so that it may be used to study non-stationary spacetimes where energy flows in and out through the quasilocal surface. As applications of this extension we have examined implications for the physics of relativistic thin shells of matter, the energy carried from a source to infinity by gravitational waves, and the transfer of energy to a body during gravitational tidal heating. The success of the formalism in all three applications provides further evidence that the Brown-York energy has physical content. Furthermore, in the tidal heating application we have seen how the quasilocal formalism provides a geometrical explanation of the gauge ambiguities that are also found in the Newtonian and pseudotensor approaches. ###### Acknowledgements. The authors would like to thank Robert Mann, Patricia Purdue, Alan Wiseman, and Kip Thorne for their useful comments and suggestions. This work was supported by the Natural Sciences and Engineering Research Council of Canada and NSF grants PHY-9728704 and AST-9731698.
warning/0003/hep-th0003178.html	ar5iv	text	# 1 Introduction ## 1 Introduction In this article we consider some applications of the general theory derived in article I . We show that nonabelian duality and Poisson-Lie duality are special cases of the general theory. In fact we show that nonabelian duality is a special case of a more general situation. The spirit of this paper is that there are natural geometric scenarios that need to be explored. We explore a few of the easier ones and see that they lead to open mathematical questions in the theory of Lie algebras. For example, Lie bialgebras and generalizations, and $`R`$-matrices naturally appear in this framework. References to equations and sections in article I are preceded by I, e.g., (I-8.3). ## 2 Examples with a flat $`\psi `$ connection ### 2.1 General remarks To best understand how to use the equations given in the Section I-8.1 it is best to do a few examples. The examples we will consider in Section 2 assume that the connection $`\psi _{ij}`$ is flat. We are mostly interested in local properties so we might as well assume $`P`$ is parallelizable. We can use parallel transport with respect to this connection to get a global framing. In this special framing the connection coefficients vanish and thus we can make the substitution $`\psi _{ij}=0`$ in all the equations in Section I-8.1. A previous remark tells us that in this case $`f_{ijk}`$ and $`\stackrel{~}{f}_{ijk}`$ are pullbacks respectively of tensors on $`M`$ and $`\stackrel{~}{M}`$. Consequently we have that $`df_{ijk}=f_{ijkl}^{}\theta ^l`$ and $`d\stackrel{~}{f}_{ijk}=\stackrel{~}{f}_{ijkl}^{\prime \prime }\stackrel{~}{\theta }^l`$, i.e., $`f_{ijkl}^{\prime \prime }=\stackrel{~}{f}_{ijkl}^{}=0`$. Because we are interested in mostly local considerations we might as well assume both $`M`$ and $`\stackrel{~}{M}`$ are parallelizable. ### 2.2 The tensor $`n_{ij}`$ is the pullback of a tensor on $`\stackrel{~}{M}`$ Here we assume that the connection $`\psi `$ is flat. In Section I-8.2 we considered what happened if $`n_{ij}`$ was covariantly constant in this section we relax this condition to one where $`n_{ij}`$ only depends on a natural subset of the variables. We assume that $`n_{ij}`$ is the pullback of a tensor on $`\stackrel{~}{M}`$. This means that $$dn_{ij}=n_{ijk}^{\prime \prime }\stackrel{~}{\theta }^k\text{or equivalently}n_{ijk}^{}=0.$$ (2.1) An equivalent formulation is that $`\frac{1}{2}n_{ij}\stackrel{~}{\theta }^i\stackrel{~}{\theta }^j`$ is the pullback of a $`2`$-form on $`\stackrel{~}{M}`$. We immediately see from (I-8.17) and (I-8.18) that $`n_{ijk}^{\prime \prime }`$ $`=`$ $`m_{kl}f_{lij},`$ (2.2) $`n_{kij}^{\prime \prime }n_{kji}^{\prime \prime }`$ $`=`$ $`m_{jl}f_{lki}m_{il}f_{lkj}=\stackrel{~}{f}_{lij}m_{lk}.`$ (2.3) A brief computation shows that $`d(m_{ji}\stackrel{~}{\theta }^j)=0`$ and therefore we can locally find $`n`$ functions $`p_1,\mathrm{},p_n`$ such that $$dp_i=m_{ji}\stackrel{~}{\theta }^j.$$ (2.4) Note that since $`m_{ij}`$ is invertible the differentials $`\{dp_1,\mathrm{},dp_n\}`$ are linearly independent. The functions $`p_1,\mathrm{},p_n`$ are the pullbacks of functions locally defined on $`\stackrel{~}{M}`$ and they define a local coordinate system on $`\stackrel{~}{M}`$. We immediately see that $$dn_{ij}=f_{kij}dp_k.$$ (2.5) The integrability condition $`0=d^2n_{ij}=f_{kijl}^{}\theta ^ldp_k`$ immediately tells us that $`f_{kijl}^{}=0`$ since $`\{\theta _1,\mathrm{}\theta _n,dp_1,\mathrm{},dp_n\}`$ are linearly independent. Since $`f_{ijkl}^{}=f_{ijkl}^{\prime \prime }=0`$ we see that $`f_{ijk}`$ are constants and thus $`M`$ is diffeomorphic to a Lie group $`G`$. The $`\theta ^i`$ are the pullbacks by $`\mathrm{\Pi }`$ of the left invariant Maurer-Cartan forms on $`G`$. The next question is whether the pullback metric $`\theta ^i\theta ^i`$ is invariant under the action of the group. Choose a $`G`$ left-invariant vector field $`X`$ “along” $`M`$, i.e., $`\iota _X\stackrel{~}{\theta }^i=0`$. A brief computation shows that $$_X(\theta ^i\theta ^i)=X^kf_{ijk}(\theta ^i\theta ^j+\theta ^j\theta ^i)$$ (2.6) The metric on $`M`$ is $`G`$-invariant if and only if $`f_{ijk}=f_{jik}`$, i.e., the structure constants are totally antisymmetric, a standard result. We can integrate (2.5) to obtain $$n_{ij}=n_{ij}^0+f_{kij}p_k$$ (2.7) where $`n_{ij}^0`$ are constants. Define the constant tensor $`m^0`$ by $`m_{ij}^0=\delta _{ij}+n_{ij}^0`$. Equation (2.3) immediately gives us an expression for $`\stackrel{~}{f}_{ijk}`$ in terms of $`f_{ijk}`$ and $`p_l`$. It is a straightforward computation to verify that $`\stackrel{~}{f}_{ijk}`$ satisfies the integrability conditions for $`d\stackrel{~}{\theta }^l=\frac{1}{2}\stackrel{~}{f}_{lij}\stackrel{~}{\theta }^i\stackrel{~}{\theta }^j`$. In general $`\stackrel{~}{f}_{ijk}`$ are not constants and are thus not generically Maurer-Cartan forms for a Lie group. To determine $`H`$ we use (I-8.15). Note that our hypotheses imply that the left hand side is automatically zero. After using the Jacobi identity satisfied by the $`f_{ijk}`$ we find that $$H_{ijk}=f_{lij}n_{lk}^0+f_{ljk}n_{li}^0+f_{lki}n_{lj}^0.$$ (2.8) If we define a left invariant $`2`$-form by $`n^0=\frac{1}{2}n_{ij}^0\theta ^i\theta ^j`$ then the above equation is $`H=dn^0`$ which tells us that $`H`$ is cohomologically trivial. We use (I-8.16) to determine $`\stackrel{~}{H}`$: $`\stackrel{~}{H}_{ijk}`$ $`=`$ $`(\stackrel{~}{f}_{ijk}+\stackrel{~}{f}_{jki}+\stackrel{~}{f}_{kij})`$ (2.9) $`+`$ $`H_{ijk}(f_{ijk}+f_{jki}+f_{kij}).`$ #### 2.2.1 Cotangent bundle duality We will see in this section that cotangent bundle duality is a special case of what we have been discussing in Section 2.2. In particular this corresponds to what is often called “nonabelian duality”. The cotangent bundle has a natural symplectic structure and thus we automatically have a candidate symplectic manifold $`P`$ for free. Assume the manifold $`M`$ is parallelizable. This means that the cotangent bundle $`T^{}M`$ is trivial, i.e., it is a product space, $`P=T^{}M=M\times ^n`$. The projections $`\mathrm{\Pi }`$ and $`\stackrel{~}{\mathrm{\Pi }}`$ will be taken to be the cartesian projections and therefore $`\stackrel{~}{M}=^n`$. If $`\theta ^i`$ is an orthonormal coframe on $`M`$ then the canonical $`1`$-form on $`T^{}M`$ may be written as $`\alpha =p_i\theta ^i`$. The $`p_i`$ are coordinates along the fibers of $`\mathrm{\Pi }`$. The fibers of $`\stackrel{~}{\mathrm{\Pi }}`$ (diffeomorphic to $`M`$) given by $`dp_i=0`$, $`i=1,\mathrm{},n`$ are the same as the fibers given by $`\stackrel{~}{\theta }^i=0`$, $`i=1,\mathrm{},n`$ therefore there must exist an invertible matrix defined by functions $`\widehat{m}_{ij}`$ such that $$dp_j=\widehat{m}_{ij}\stackrel{~}{\theta }^i.$$ (2.10) The canonical symplectic form is given by $`\beta =d\alpha =dp_j\theta ^j\frac{1}{2}p_kf_{kij}\theta ^i\theta ^j`$. In going from $`\beta `$ to $`\gamma `$ the $`\stackrel{~}{\theta }\theta `$ term is not changed and so we immediately learn that $`\widehat{m}_{ij}=m_{ij}=\delta _{ij}+n_{ij}`$. Taking the exterior derivative of (2.10) leads to $`n_{ijk}^{}=0`$ and (2.3). Thus $`n_{ij}`$ is the pullback of a tensor on $`^n`$ and we are back to our general discussion given in Section 2.2. This is an example of what is often called “nonabelian duality” which has generating function given by (I-2.2) with $`\alpha =p_k\theta ^k`$. We know that $`M`$ is a Lie group $`G`$, $`T^{}G=G\times 𝔤^{}`$ and thus $`\stackrel{~}{M}=𝔤^{}`$. The metric on $`𝔤^{}`$ is immediately computable from (2.10) since $`m_{ij}=\delta _{ij}+n_{ij}^0+p_kf_{kij}`$. It is worth remarking that because $`\stackrel{~}{M}=𝔤^{}=^n`$ there is an abelian Lie group action on $`\stackrel{~}{M}`$ which does not leave the metric invariant. Said differently the $`dp`$ are the Maurer-Cartan forms for the abelian Lie group $`𝔤^{}`$. This statement is made in anticipation of our discussion about Poisson-Lie duality in Section 3. “Nonabelian duality” usually refers to the case with $`n_{ij}^0=0`$ where (2.8) tells us that $`H_{ijk}=0`$. We can compute $`\stackrel{~}{H}`$ from (2.9) directly or it was already remarked in Section I-7.2 that $`\stackrel{~}{B}_{ij}`$ is easily determined. #### 2.2.2 Cotangent bundle duality and gauge invariance In this section we revisit cotangent bundle duality and try to understand geometrically the role of the $`B`$ field gauge transformation. We assume the manifold $`M`$ has a trivial cotangent bundle $`T^{}M=M\times ^n`$ with canonical $`1`$-form $`\alpha =p_i\theta ^i`$ and symplectic form $`\beta =d\alpha `$. We modify the discussion of Section 2.2.1 by demanding that the bifibration not be given by the cartesian projections. We want the vertical fibers to be original ones so we have $`\mathrm{\Pi }:(x,p)T^{}MxM`$. On the other hand the projection $`\stackrel{~}{\mathrm{\Pi }}:T^{}M\stackrel{~}{M}`$ will not be the canonical projection. The fibers of this projection are “slanted” relative to the fibers of the cartesian projection $`\stackrel{~}{\mathrm{\Pi }}_c`$. Note that the fibers of $`\stackrel{~}{\mathrm{\Pi }}`$ are the integral manifolds of the Pfaffian equations $`\stackrel{~}{\theta }^i=0`$. From general principles we know that $$dp_i=\widehat{m}_{ji}\stackrel{~}{\theta }^j+u_{ji}\theta ^j$$ (2.11) for some functions $`\widehat{m},u`$. Geometrically this is the statement that the fibers of $`\stackrel{~}{\mathrm{\Pi }}`$ are slanted relative to the fibers of $`\stackrel{~}{\mathrm{\Pi }}_c`$. As before the structure of the symplectic form leads to the result that $`\widehat{m}_{ij}=m_{ij}`$. The integrability condition $`d^2p_i=0`$ leads to the following equations $`n_{ijk}^{\prime \prime }n_{ikj}^{\prime \prime }`$ $`=`$ $`\stackrel{~}{f}_{ljk}m_{li},`$ (2.12) $`u_{jik}^{}u_{kij}^{}`$ $`=`$ $`f_{ljk}u_{li},`$ (2.13) $`u_{jik}^{\prime \prime }`$ $`=`$ $`n_{kij}^{}.`$ (2.14) Comparing (2.12) with (I-8.18) we see that $`n_{ijk}^{}=0`$ and thus $`n_{ij}`$ is the pullback of a tensor on $`\stackrel{~}{M}`$. The general discussion of Section 2.2 tells us that there exists function $`\widehat{p}_i`$ in $`T^{}M`$ such that $`d\widehat{p}_i=m_{ji}\stackrel{~}{\theta }^j`$. Geometrically this is just the statement that the fibers is given by $`\widehat{p}_i=\text{constant}`$. Note that (2.14) tells us that $`u_{ijk}^{\prime \prime }=0`$ and thus $`u_{ij}`$ is the pullback of a tensor on $`M`$. Equation (2.11) tells us that $`d(p_i\widehat{p}_i)=u_{ji}\theta ^j`$ and thus we see that $`p_i=\widehat{p}_i+k_i`$ where $`k_i`$ are pullbacks of functions on $`M`$, i.e., $`k_i=k_i(x)`$. Thus we we see that the canonical transformation generated by $`p_i\theta ^i`$ is gauge equivalent to the one generated by $`\widehat{p}_i\theta ^i`$ and corresponds to a different choice of fibration. #### 2.2.3 Can $`\stackrel{~}{M}`$ naturally be a Lie group? Is it possible for the dual manifold to be naturally a Lie group under the assumptions underlying the discussions in Section 2.2? By “naturally” we mean that the $`\stackrel{~}{\theta }^i`$ are the Maurer-Cartan forms for a Lie group $`\stackrel{~}{G}`$. Solving (2.3) for $`\stackrel{~}{f}_{ijk}`$ generally leads to a non-constant solution. There is a possibility that the solution may be constant, i.e., $`(\stackrel{~}{M},\stackrel{~}{g},\stackrel{~}{B})`$ is naturally a Lie group. Unfortunately we will see that both $`𝔤`$ and $`\stackrel{~}{𝔤}`$ must be abelian and so there are no new interesting examples. Inserting (2.7) into (2.3) and using the linear independence of the $`dp_i`$ leads to two equations: $`m_{jl}^0f^l{}_{ki}{}^{}m_{il}^0f^l_{kj}`$ $`=`$ $`\stackrel{~}{f}^l{}_{ij}{}^{}m_{lk}^{0},`$ (2.15) $`f^m{}_{il}{}^{}f_{}^{l}_{jk}`$ $`=`$ $`f^m{}_{il}{}^{}\stackrel{~}{f}_{}^{l}{}_{jk}{}^{}.`$ (2.16) To obtain the latter equation we used the Jacobi identity satisfied by $`f_{ijk}`$. The indices have also been set at their natural (co)variances for future convenience. The $`\stackrel{~}{f}^i_{jk}`$ satisfy the Jacobi identity if they satisfy the two equations above because of our remark about integrability after (2.7). You should think of $`M`$ as a Lie group $`G`$ with Lie algebra $`𝔤`$ and similarly for $`\stackrel{~}{M}`$. Let $`\mathrm{ad}_X(Y)=[X,Y]`$ be the adjoint action of $`𝔤`$. Let $`[,]^{}`$ be the Lie bracket on $`\stackrel{~}{𝔤}`$ with adjoint action denoted by $`\stackrel{~}{\mathrm{ad}}_{\stackrel{~}{X}}\stackrel{~}{Y}=[\stackrel{~}{X},\stackrel{~}{Y}]^{}`$. The reduction of the structure group of $`P`$ to $`\mathrm{O}(n)`$ meant that we could identify the horizontal tangent space with the vertical tangent space and thus we can identify $`𝔤`$ with $`\stackrel{~}{𝔤}`$. We should think of a single vector space $`V`$ with two Lie brackets giving us two Lie algebras $`𝔤=(V,[,])`$ and $`\stackrel{~}{𝔤}=(V,[,]^{})`$. In this notation (2.16) becomes $$\mathrm{ad}_X(\mathrm{ad}_Y\stackrel{~}{\mathrm{ad}}_Y)=0.$$ (2.17) There are some immediate important consequences of this equation. Let $`𝔡`$ be the vector subspace of $`𝔤`$ spanned by $`\mathrm{ad}_YZ\stackrel{~}{\mathrm{ad}}_YZ`$ for all $`Y,Z`$. The subspace $`𝔡`$ is contained in the center $`𝔷`$ of $`𝔤`$ by (2.17). If $`𝔡0`$, i.e., $`\mathrm{ad}\stackrel{~}{\mathrm{ad}}`$, then the center $`𝔷`$ is a nontrivial abelian ideal in the Lie algebra $`𝔤`$ and thus $`𝔤`$ is not semisimple. If $`𝔡=0`$ then we have that $`\mathrm{ad}=\stackrel{~}{\mathrm{ad}}`$ and we will show that $`𝔤`$ is abelian when we incorporate (2.15) into our reasoning. Note that if $`𝔷=0`$ then $`𝔡=0`$. Let us study the implications of $`\mathrm{ad}=\stackrel{~}{\mathrm{ad}}`$, i.e., $`𝔡=0`$. Equation (2.15) may be rewritten as $$2f_{kij}=(f_{ijk}+f_{jki}+f_{kij})+(n_{il}^0f_{ljk}+n_{jl}^0f_{lki}+n_{kl}^0f_{lij}).$$ (2.18) Notice that the right hand side is totally antisymmetric under permutations of $`i,j,k`$. This means that $`f_{kij}`$ is totally antisymmetric and thus the metric $`\theta ^i\theta ^i`$ is a bi-invariant positive definite metric on $`G`$. The proposition proven in Appendix A tells us that there is a decomposition into ideals $`𝔤=𝔨𝔞`$ where $`𝔨`$ is a compact semisimple Lie algebra and $`𝔞`$ is an abelian Lie algebra. The next part of the argument only involves the compact ideal $`𝔨`$ and without any loss we assume $`𝔞=0`$ for the moment. Using the antisymmetry of the structure constants the above may be rewritten as $$f_{ijk}=n_{il}^0f_{ljk}+n_{jl}^0f_{lki}+n_{kl}^0f_{lij}.$$ Using the remark made immediately after (2.8) we see that $`\frac{1}{3!}f_{ijk}\theta ^i\theta ^j\theta ^k=dn^0`$. The closed $`3`$-form $`\frac{1}{3!}f_{ijk}\theta ^i\theta ^j\theta ^k`$ is exact and this contradicts $`H^3(𝔨)0`$ as discussed in Appendix A. If $`\mathrm{ad}=\stackrel{~}{\mathrm{ad}}`$ then we conclude that $`𝔤=𝔞`$ is abelian and we are back in familiar territory. The next we consider the general case by exploiting the observation that $`𝔡𝔷`$. Since we have a positive definite metric on $`𝔤`$ there is an orthogonal direct sum decomposition $`𝔤=𝔷𝔷^{}`$. The orthogonal complement $`𝔷^{}`$ is not generally an ideal because the metric is not necessarily $`(\mathrm{ad}𝔤)`$-invariant. Nevertheless we can choose an orthonormal basis $`\{e_\alpha \}`$ for $`𝔷`$ and an orthonormal basis $`\{e_a\}`$ for $`𝔷^{}`$. Greek indices are associated with $`𝔷`$, indices in $`\{a,b,c,d\}`$ are associated with $`𝔷^{}`$; and indices in $`\{i,j,k,l\}`$ run from $`1,\mathrm{},dim𝔤`$ and are associated with all of $`𝔤`$. The Lie algebra $`𝔤`$ is given by $`[e_\alpha ,e_j]`$ $`=`$ $`0,`$ (2.19) $`[e_a,e_b]`$ $`=`$ $`f^c{}_{ab}{}^{}e_{c}^{}+f^\gamma {}_{ab}{}^{}e_{\gamma }^{}.`$ (2.20) Note that $`𝔥=𝔤/𝔷`$ is a Lie algebra since $`𝔷`$ is an ideal. It follows that the structure constants of $`𝔥`$ are $`f^c_{ab}`$. Also, $`𝔤`$ is a central extension of $`𝔥`$ with extension cocycle $`f^\gamma _{ab}`$. It is well known that the cocycle is trivial if $`f^\gamma {}_{ab}{}^{}=t^\gamma {}_{c}{}^{}f_{}^{c}_{ab}`$ corresponding to a redefinition of the basis given by $`e_ae_a+t^\gamma {}_{a}{}^{}e_{\gamma }^{}`$. The Lie algebra $`\stackrel{~}{𝔤}`$ must have the form below because the two Lie brackets are the same modulo the center $`𝔷`$: $`[e_\alpha ,e_j]^{}`$ $`=`$ $`\stackrel{~}{f}^\gamma {}_{\alpha j}{}^{}e_{\gamma }^{},`$ (2.21) $`[e_a,e_b]^{}`$ $`=`$ $`f^c{}_{ab}{}^{}e_{c}^{}+\stackrel{~}{f}^\gamma {}_{ab}{}^{}e_{\gamma }^{}.`$ (2.22) In equation (2.15) choose $`k=\gamma `$ then the left hand side vanishes and the equation becomes $`0=\stackrel{~}{f}^l{}_{ij}{}^{}m_{l\gamma }^{0}=\stackrel{~}{f}^d{}_{ij}{}^{}m_{d\gamma }^{0}+\stackrel{~}{f}^\delta {}_{ij}{}^{}m_{\delta \gamma }^{0}`$. Using the different choices for $`(i,j)`$ we find $`\stackrel{~}{f}^\delta _{\alpha j}`$ $`=`$ $`\stackrel{~}{f}^d{}_{\alpha j}{}^{}m_{d\gamma }^{0}((m^0)^1)^{\gamma \delta }=0,`$ $`\stackrel{~}{f}^\delta _{ab}`$ $`=`$ $`f^d{}_{ab}{}^{}m_{d\gamma }^{0}((m^0)^1)^{\gamma \delta }.`$ (2.23) These equations tell us that $`\stackrel{~}{𝔤}`$ is a central extension of $`𝔥`$ with a trivial cocycle. Next choose $`i=\alpha `$ in (2.15) with result $`0=m_{\alpha l}^0f^l_{kj}`$. Choosing $`k=a`$ and $`j=b`$ leads to $$f^\gamma {}_{ab}{}^{}=((m^0)^1)^{\gamma \alpha }m_{\alpha c}^0f^c{}_{ab}{}^{}=0$$ (2.24) which tells us that the cocycle $`f^\gamma _{ab}`$ is also trivial. Finally choose $`(i,j,k)=(a,b,c)`$ in (2.15) and substitute (2.23) and (2.24) for $`\stackrel{~}{f}^\delta _{ab}`$ and $`f^\gamma _{ab}`$ to obtain $`(m_{bd}^0m_{b\gamma }^0((m^0)^1)^{\gamma \alpha }m_{\alpha d}^0)f^d_{ca}`$ $``$ $`(m_{ad}^0m_{a\gamma }^0((m^0)^1)^{\gamma \alpha }m_{\alpha d}^0)f^d_{cb}`$ (2.25) $`=`$ $`f^d{}_{ab}{}^{}(m_{dc}^0m_{d\alpha }^0((m^0)^1)^{\alpha \gamma }m_{\gamma c}^0).`$ Next we observe that since $`m_{\alpha \beta }^0=\delta _{\alpha \beta }+n_{\alpha \beta }`$ we conclude that $`((m^0)^1)^{\alpha \beta }=s^{\alpha \beta }+a^{\alpha \beta }`$ where $`s^{\alpha \beta }`$ is symmetric and positive definite and $`a^{\alpha \beta }`$ is skew. In particular note that $`m_{b\gamma }^0((m^0)^1)^{\gamma \alpha }m_{\alpha d}^0=s^{\gamma \alpha }m_{\gamma b}^0m_{\alpha d}^0+a^{\gamma \alpha }m_{\gamma b}^0m_{\alpha d}^0`$ where the first term is symmetric and positive definite and the second is skew. Using this we immediately see that $$m_{bd}^0m_{b\gamma }^0((m^0)^1)^{\gamma \alpha }m_{\alpha d}^0=S_{bd}+N_{bd}$$ where $`S_{bd}`$ is symmetric and positive definite and $`N_{bd}`$ is skew. We can use $`S_{ab}`$ as a second metric on $`𝔤`$ and use it to “raise and lower” the indices in (2.25) leading to $$2f_{cab}=(f_{abc}+f_{bca}+f_{cab})+(N_{ad}f^d{}_{bc}{}^{}+N_{bd}f^d{}_{ca}{}^{}+N_{cd}f^d{}_{ab}{}^{}).$$ (2.26) This tells us that $`f_{abc}`$ is totally antisymmetric and therefore $`S_{ab}`$ is an invariant metric on $`𝔥`$. The same chain of arguments used after (2.18) tell us that $`𝔥`$ is abelian which implies that $`\stackrel{~}{f}^\delta {}_{ab}{}^{}=0`$ and $`f^\gamma {}_{ab}{}^{}=0`$ concluding that $`𝔤`$ and $`\stackrel{~}{𝔤}`$ are abelian Lie algebras. #### 2.2.4 Connection with $`R`$-matrices Equation (2.15) is closely related to the theory of $`R`$-matrices developed by Semenov-Tian-Shansky which is different than the one developed by Drinfeld . The discussion here suggests that Poisson-Lie groups may play a role in duality, see Section 3. We observe that (2.15) may be rewritten as $$(m^0)_{jl}^1f_{kil}(m^0)_{kl}^1f_{jil}=(m^0)_{jl}^1(m^0)_{km}^1\stackrel{~}{f}_{nlm}m_{ni}^0.$$ (2.27) Next we show that this equation describes a potential double Lie algebra structure on $`𝔤`$ à la Semenov-Tian-Shansky (B.11). Assume the Lie algebra $`𝔤`$ admits an invariant metric $`K`$. For example if the Lie algebra is semisimple then $`K`$ may be taken to be the Killing metric. Let us use the indices $`a,b,c,d,e`$ to denote generic components in a generic basis. In terms of a basis $`(e_1,\mathrm{},e_n)`$ the structure constants are given by $`[e_a,e_b]=f^c{}_{ab}{}^{}e_{c}^{}`$. For the moment it is best to forget about the orthonormal basis we were previously using before because the metric $`K`$ may not be related to the previous metric. The components of the invariant metric are given by $`K_{ab}=K(e_a,e_b)`$. We will use $`K`$ to identify $`𝔤`$ with $`𝔤^{}`$, i.e., raise and lower indices. The tensor with components $`K^{ab}`$ is the inverse of the invariant metric, i.e., the induced metric on $`𝔤^{}`$. The statement that $`K`$ is $`𝔤`$-invariant is equivalent to $`f_{abc}`$ being totally antisymmetric. With these assumptions equation (B.13) may be rewritten as $$(R^d{}_{e}{}^{}K_{}^{ea})f^b{}_{cd}{}^{}(R^d{}_{e}{}^{}K_{}^{eb})f^a{}_{cd}{}^{}=K^{aa^{}}K^{bb^{}}K_{cc^{}}(f^c^{}{}_{a^{}b^{}}{}^{})_R.$$ (2.28) The indices of $`R`$ and $`f`$ have their natural (co)variances. We are now going to compare (2.28) with (2.27). Remember that (2.27) is valid in an orthonormal basis with respect to a specific metric on $`𝔤`$ which may not be related to the invariant metric $`K`$. What we will have to do is express (2.28) in the orthonormal basis but this will be simple because we adjusted all the (co)variances correctly in the equation above. The indices $`i,j,k,l,m,n`$ refer to our orthonormal basis. We raise and lower indices using the Kronecker delta tensor. The only potential confusion is that we have to be careful and remember that $`K_{ab}`$ becomes $`K_{ij}`$ and the inverse invariant metric $`K^{ab}`$ becomes $`K_{ij}^1`$. The correspondence is made by choosing $`R`$ (in our orthonormal basis) to be defined by $$(m^0)_{ij}^1=R_{jl}K_{li}^1.$$ (2.29) The $`R`$-matrices we are considering are invertible because both $`m^0`$ and $`K`$ are always invertible. This is different that in the Drinfeld case where the $`R`$-matrix is skew adjoint and may not be invertible. If you think of $`K`$ as a map $`K:𝔤𝔤^{}`$ then the equation above is $`(m^0)^t=KR^1:𝔤𝔤^{}`$ which suggests that $`m^0`$ should be interpreted as a map $`m^0:𝔤^{}𝔤`$. Comparing the right hand sides of (2.27) and (2.28) we see that $$\stackrel{~}{f}_{ilm}R_{lj}R_{mk}=R_{il}(f_{ljk})_R.$$ (2.30) If $`R_{ij}`$ satisfies the modified Yang Baxter equation (B.10) then $`(f^i{}_{jk}{}^{})_R`$ are the structure constants of a Lie algebra $`𝔤_R`$ associated with a Lie group $`G_R`$. Let $`(\mu ^1,\mathrm{},\mu ^n)`$ be the left invariant Maurer-Cartan forms for $`G_R`$: $$d\mu _l=\frac{1}{2}(f_{ljk})_R\mu _j\mu _k.$$ Going to a different basis given by $`\lambda _i=R_{il}\mu _l`$ we see that $$d\lambda _i=\frac{1}{2}\stackrel{~}{f}_{ilm}\lambda _l\lambda _m.$$ We conclude that if $`R`$ is a solution of the modified classical Yang Baxter equation then we can construct the Lie algebra $`\stackrel{~}{𝔤}`$ with structure constant $`\stackrel{~}{f}_{ijk}`$ and an associated Lie group $`\stackrel{~}{G}`$. In our situation we also have to impose a second equation (2.16) and this is very restrictive. Our in-depth analysis that eliminated all except the abelian case. Logically, there is the possibility of non-trivial solutions by the use of $`R`$-matrices. The reason is that using (2.30) and (B.11) we can rewrite (2.16) as $$0=f^m{}_{il}{}^{}(f^l{}_{jk}{}^{}R_{}^{j}{}_{p}{}^{}R_{}^{k}{}_{q}{}^{}f^r{}_{nq}{}^{}R_{}^{n}{}_{p}{}^{}R_{}^{l}{}_{r}{}^{}f^r{}_{pn}{}^{}R_{}^{n}{}_{q}{}^{}R_{}^{l}{}_{r}{}^{}).$$ The stuff between the parentheses is the left hand side of modified classical Yang Baxter equation (B.12). If $`R`$ satisfies the modified classical Yang Baxter equation the above becomes $`0=cf^m{}_{il}{}^{}f_{}^{l}_{pq}`$. The solution to this equation is $`c=0`$ or $`\mathrm{ad}_X\mathrm{ad}_Y=0`$. We know from our general analysis that it is unnecessary to proceed along these lines. ### 2.3 Symmetric duality In the previous sections we had $`\alpha =p_i\theta ^i`$ for some functions $`p_i`$ on $`P`$. A consequence was that the $`p_i`$ were good coordinates on the fibers of $`\mathrm{\Pi }`$ and these fibers were also lagrangian submanifolds of $`\beta `$. There was a certain asymmetry in the way fibers of $`\mathrm{\Pi }`$ and of $`\stackrel{~}{\mathrm{\Pi }}`$ were treated. In this section we consider a more symmetric situation. To motivate the ensuing presentation let us go to the case of $`M=^n`$ and $`P=T^{}M`$ where we have the symplectic form $`\beta =dp_idq^i`$. Up to constant $`1`$-forms the most general antiderivative is $`\alpha =(1u)p_idq^iuq^idp_i`$ where $`u`$ is a constant. Let $`UP`$ be a neighborhood, we can always write $$\alpha =\stackrel{~}{q}_i\theta ^iq^i\stackrel{~}{\theta }_i,$$ (2.31) where $`q^i`$ and $`\stackrel{~}{q}_i`$ are local functions on $`U`$. We now make some special assumptions about the functions $`q`$ and $`\stackrel{~}{q}`$. Assume there exist matrix valued functions $`E`$ and $`\stackrel{~}{E}`$ such that $`dq^i`$ $`=`$ $`E_{ij}\theta ^j,`$ (2.32) $`d\stackrel{~}{q}_j`$ $`=`$ $`\stackrel{~}{\theta }_i\stackrel{~}{E}_{ij}.`$ (2.33) These functions are not arbitrary and must satisfy a variety of constraints. For example, the ranks are constrained by equation (2.41) below. For example if $`E`$ and $`\stackrel{~}{E}`$ are invertible then the functions $`(q,\stackrel{~}{q})`$ are independent in the sense that the map $`(q^1,\mathrm{},q^n,\stackrel{~}{q}_1,\mathrm{},\stackrel{~}{q}_n):U^{2n}`$ is of rank $`2n`$. Consequently the fibers of $`\mathrm{\Pi }`$ are locally described by $`(q=\text{constant})`$ and the fibers of $`\stackrel{~}{\mathrm{\Pi }}`$ by $`(\stackrel{~}{q}=\text{constant})`$. In general $`E`$ or $`\stackrel{~}{E}`$ will not be invertible. In this case not all the functions will be independent and $`(q=\text{constant})`$ will define a family of manifolds each diffeomorphic to $`\stackrel{~}{M}`$. As before we use the prime and double prime notation to denote derivatives in the appropriate directions. The equation $`d^2q^i=0`$ tells us that $`E_{ijk}^{\prime \prime }`$ $`=`$ $`0,`$ (2.34) $`E_{ijk}^{}E_{ikj}^{}`$ $`=`$ $`E_{il}f_{ljk}.`$ (2.35) Likewise the equation $`d^2\stackrel{~}{q}=0`$ tells us that $`\stackrel{~}{E}_{ijk}^{}`$ $`=`$ $`0,`$ (2.36) $`\stackrel{~}{E}_{kjl}^{\prime \prime }\stackrel{~}{E}_{ljk}^{\prime \prime }`$ $`=`$ $`\stackrel{~}{f}_{ikl}\stackrel{~}{E}_{ij}.`$ (2.37) Note that you can solve (2.35) for $`E_{ijk}^{}`$ as a function of $`E_{il}f_{ljk}`$ and you can solve (2.37) for $`\stackrel{~}{E}_{ijk}^{\prime \prime }`$ and as function of $`\stackrel{~}{f}_{ikl}\stackrel{~}{E}_{ij}`$. Next we compute $`\beta =d\alpha `$: $$\beta =(\stackrel{~}{E}_{ij}+E_{ij})\stackrel{~}{\theta }^i\theta ^j\frac{1}{2}\stackrel{~}{q}_if_{ijk}\theta ^j\theta ^k+\frac{1}{2}q^i\stackrel{~}{f}_{ijk}\stackrel{~}{\theta }^j\stackrel{~}{\theta }^k.$$ (2.38) In general neither the fibers of $`\mathrm{\Pi }`$ or of $`\stackrel{~}{\mathrm{\Pi }}`$ are lagrangian submanifolds of $`\beta `$. This is very different from the cotangent bundle cases previously discussed. Next we make we write $`E`$ and $`\stackrel{~}{E}`$ as $`E_{ij}`$ $`=`$ $`\sigma _{ij}+\nu _{ij},`$ (2.39) $`\stackrel{~}{E}_{ij}`$ $`=`$ $`\stackrel{~}{\sigma }_{ij}+\stackrel{~}{\nu }_{ij}`$ (2.40) where $`\sigma `$ and $`\stackrel{~}{\sigma }`$ are symmetric, and $`\nu `$ and $`\stackrel{~}{\nu }`$ are antisymmetric. Comparing to (I-6.5) we see that $$\delta _{ij}=\sigma _{ij}+\stackrel{~}{\sigma }_{ij}\text{and}n_{ij}=\nu _{ij}+\stackrel{~}{\nu }_{ij}.$$ (2.41) The matrix valued functions $`E`$ and $`\stackrel{~}{E}`$ cannot be arbitrary and must satisfy conditions given by the above. We remark that $`n_{ijk}^{}=m_{ijk}^{}=E_{ijk}^{}+\stackrel{~}{E}_{ijk}^{}=E_{ijk}^{}=\sigma _{ijk}^{}+\nu _{ijk}^{}`$. Since $`n_{ij}=n_{ji}`$ we immediately see that $$\sigma _{ijk}^{}=0,\text{and}E_{ijk}^{}=n_{ijk}^{}=\nu _{ijk}^{}.$$ (2.42) Also $`E_{ijk}^{\prime \prime }=0`$ and thus we see that $`\sigma _{ijk}^{\prime \prime }=0`$ and $`\nu _{ijk}^{\prime \prime }=0`$. This means that $`\sigma _{ij}`$ is constant. Similarly we conclude that $`\stackrel{~}{\sigma }_{ij}`$ is constant, $`\stackrel{~}{E}_{ijk}^{\prime \prime }=n_{ijk}^{\prime \prime }=\stackrel{~}{\nu }_{ijk}^{\prime \prime }`$ and $`\stackrel{~}{\nu }_{ijk}^{}=0`$. We can take the results given above and insert into (I-8.18) and (I-8.17) to obtain $`E_{ijk}^{}=n_{ijk}^{}=\nu _{ijk}^{}`$ $`=`$ $`\stackrel{~}{f}_{lij}E_{lk}.`$ (2.43) $`\stackrel{~}{E}_{ijk}^{\prime \prime }=n_{ijk}^{\prime \prime }=\stackrel{~}{\nu }_{ijk}^{\prime \prime }`$ $`=`$ $`+\stackrel{~}{E}_{kl}f_{lij}.`$ (2.44) We now insert the above into (2.35) and (2.37) to obtain the basic equations $`\stackrel{~}{E}_{kl}f_{lij}\stackrel{~}{E}_{jl}f_{lik}`$ $`=`$ $`\stackrel{~}{f}_{ljk}\stackrel{~}{E}_{li},`$ (2.45) $`\stackrel{~}{f}_{lij}E_{lk}\stackrel{~}{f}_{lik}E_{lj}`$ $`=`$ $`E_{il}f_{ljk}.`$ (2.46) Note that $`dE_{ij}`$ $`=`$ $`E_{ijk}^{}\theta ^k=\stackrel{~}{f}_{lij}E_{lk}\theta ^k=\stackrel{~}{f}_{lij}dq^l`$ (2.47) $`d\stackrel{~}{E}_{ij}`$ $`=`$ $`\stackrel{~}{E}_{ijk}^{\prime \prime }\stackrel{~}{\theta }^k=\stackrel{~}{E}_{kl}f_{lij}\stackrel{~}{\theta }^k=f_{lij}d\stackrel{~}{q}^l`$ (2.48) where we used (2.32) and (2.33). We know that $`df_{lij}=f_{lijk}^{}\theta ^k`$ and $`d\stackrel{~}{f}_{lij}=\stackrel{~}{f}_{lijk}^{\prime \prime }\stackrel{~}{\theta }^k`$. Therefore by taking the exterior derivatives of (2.47) and (2.48) we learn that $`\stackrel{~}{E}_{ml}(df_{lij})`$ $`=`$ $`0,`$ (2.49) $`(d\stackrel{~}{f}_{lij})E_{lm}`$ $`=`$ $`0.`$ (2.50) To make progress we have to make some assumptions. The simplest assumption is that $`E_{ij}=0`$. In this case we are back to the discussion given in Section 2.2.1. In this paragraph we use the notation that a Lie group $`G`$ via nonabelian duality is dual to $`\stackrel{~}{G}𝔤^{}`$. You can generalize cotangent bundle duality along the following lines. Consider matrices with with the same $`2\times 2`$ block form $$E=\left(\begin{array}{cc}E^{(1)}& 0\\ 0& 0\end{array}\right),\stackrel{~}{E}=\left(\begin{array}{cc}0& 0\\ 0& \stackrel{~}{E}^{(2)}\end{array}\right).$$ This will lead to a manifold $`M=\stackrel{~}{G}^{(1)}\times G^{(2)}`$ and $`\stackrel{~}{M}=G^{(1)}\times \stackrel{~}{G}^{(2)}`$. Next we look at the case where $`\stackrel{~}{E}_{ij}`$ is invertible. Equation (2.49) tells us that $`f_{ijk}`$ are constants and thus $`M`$ is naturally a Lie group $`G`$ with structure constants $`f_{ijk}`$. We can immediately integrate (2.48) obtaining $$\stackrel{~}{E}_{ij}=\stackrel{~}{\sigma }_{ij}+\stackrel{~}{\nu }_{ij}^0+f_{kij}\stackrel{~}{q}^k,$$ (2.51) where $`\stackrel{~}{\nu }^0`$ is a constant tensor. With this information we can use (2.45) to determine $`\stackrel{~}{f}_{ijk}`$. It is straightforward to verify that $`\stackrel{~}{f}_{ijk}`$ satisfy the integrability conditions for (I-8.2) with $`\psi _{ij}=0`$. All we have to do is to find a tensor $`E_{ij}`$ that satisfies (2.46) and (2.47). Note that (2.41) tells us that $`\sigma _{ij}=\delta _{ij}\stackrel{~}{\sigma }_{ij}`$. This case merits further analysis. If besides $`\stackrel{~}{E}_{ij}`$ being invertible we also impose that $`E_{ij}`$ is invertible then $`\stackrel{~}{f}_{ijk}`$ are constant (see (2.49)) and $`\stackrel{~}{M}`$ is naturally a Lie group $`\stackrel{~}{G}`$. We can integrate (2.47) to obtain $$E_{ij}=\sigma _{ij}+\nu _{ij}^0+\stackrel{~}{f}_{lij}q^l,$$ (2.52) where $`\nu ^0`$ is a constant tensor. It is convenient to define $`E^0`$ $`=`$ $`\sigma +\nu ^0,`$ (2.53) $`\stackrel{~}{E}^0`$ $`=`$ $`\stackrel{~}{\sigma }+\stackrel{~}{\nu }^0.`$ (2.54) We can insert (2.51) and (2.52) into (2.45) and (2.46) and expand both sides in powers of $`q`$ and $`\stackrel{~}{q}`$ to obtain: $`\stackrel{~}{E}_{jl}^0f^l{}_{ki}{}^{}\stackrel{~}{E}_{il}^0f^l_{kj}`$ $`=`$ $`\stackrel{~}{f}^l{}_{ij}{}^{}\stackrel{~}{E}_{lk}^{0},`$ (2.55) $`f^m{}_{il}{}^{}f_{}^{l}_{jk}`$ $`=`$ $`f^m{}_{il}{}^{}\stackrel{~}{f}_{}^{l}{}_{jk}{}^{},`$ (2.56) $`\stackrel{~}{f}^l{}_{ki}{}^{}E_{lj}^{0}\stackrel{~}{f}^l{}_{kj}{}^{}E_{li}^{0}`$ $`=`$ $`E_{kl}^0f^l{}_{ij}{}^{},`$ (2.57) $`\stackrel{~}{f}^m{}_{il}{}^{}\stackrel{~}{f}_{}^{l}_{jk}`$ $`=`$ $`\stackrel{~}{f}^m{}_{il}{}^{}f_{}^{l}{}_{jk}{}^{},`$ (2.58) where the Jacobi identity was used to simplify the above. We are now in a situation very similar to that in Section 2.2.3. The difference is that the relevant metrics are now $`\sigma _{ij}`$ and $`\stackrel{~}{\sigma }_{ij}`$. The difficulty arises in that we have lost positive definiteness of the metrics. The only constraint is that $`\sigma +\stackrel{~}{\sigma }=I`$. If either $`\sigma `$ or $`\stackrel{~}{\sigma }`$ is definite then an analysis along the lines of Section 2.2.3 leads to the conclusion that $`𝔤`$ and $`\stackrel{~}{𝔤}`$ are abelian. If both are indefinite or both are singular then the analysis previously provided breaks down. This situation also merits further study. ## 3 Poisson-Lie duality Here we discuss a beautiful example of scenario I-3 described at the end of Section I-8.1 where we are given a special symplectic bifibration and we have to construct the metrics and antisymmetric tensors on $`M`$ and $`\stackrel{~}{M}`$. The Drinfeld double Lie group is an example of a special symplectic bifibration. The metrics and antisymmetric tensors constructed in this manner correspond to the Poisson-Lie duality of Klimcik and Severa . The explicit duality transformation was obtained by Sfetsos . We can use our general framework to determine both by making educated guesses. The Drinfeld double $`G_D`$ is a Lie group whose Lie algebra $`𝔤_D`$ is a Lie bialgebra, see Appendix B.1. The bifibration is by Lie groups $`G`$ and $`\stackrel{~}{G}`$ with respective Lie algebras $`𝔤`$ and $`\stackrel{~}{𝔤}𝔤^{}`$. The Lie algebras are related by $`𝔤_D=𝔤\stackrel{~}{𝔤}=𝔤𝔤^{}`$. If $`\{T_a\}`$ is a basis for $`𝔤`$ and $`\{\stackrel{~}{T}^a\}`$ is the associated dual basis for $`𝔤^{}`$ then $`[T_a,T_b]`$ $`=`$ $`C^c{}_{ab}{}^{}T_{c}^{},`$ $`[T_a,\stackrel{~}{T}^b]`$ $`=`$ $`\stackrel{~}{C}_a{}_{}{}^{bc}T_{c}^{}C^b{}_{ac}{}^{}\stackrel{~}{T}_{}^{c},`$ $`[\stackrel{~}{T}^a,\stackrel{~}{T}^b]`$ $`=`$ $`\stackrel{~}{C}_c{}_{}{}^{ab}\stackrel{~}{T}_{}^{c},`$ where $`C^c_{ab}`$ and $`\stackrel{~}{C}_a^{bc}`$ are respectively the structure constants for $`G`$ and $`\stackrel{~}{G}`$. The two sets of structure constants must satisfy compatibility condition (B.4). To write down the symplectic structure in a convenient way we introduce some notation slightly different than the one given in . Let $`gG`$ then the adjoint representation on $`𝔤`$ is given by $`gT_bg^1=T_aa^a{}_{b}{}^{}(g)`$. One also has $`g\stackrel{~}{T}^ag^1=a^a{}_{b}{}^{}(g^1)(\stackrel{~}{T}^b+\mathrm{\Pi }^{bc}(g)T_c)`$ where $`\mathrm{\Pi }^{ab}=\mathrm{\Pi }^{ba}`$. Similarly one has that if $`\stackrel{~}{g}\stackrel{~}{G}`$ then $`\stackrel{~}{g}\stackrel{~}{T}^b\stackrel{~}{g}^1=\stackrel{~}{T}^a\stackrel{~}{a}_a{}_{}{}^{b}(\stackrel{~}{g})`$ and $`\stackrel{~}{g}T_a\stackrel{~}{g}^1=\stackrel{~}{a}_a{}_{}{}^{b}(\stackrel{~}{g}^1)(T_b+\stackrel{~}{\mathrm{\Pi }}_{bc}(\stackrel{~}{g})\stackrel{~}{T}^c)`$ where $`\stackrel{~}{\mathrm{\Pi }}`$ is antisymmetric. It is worthwhile to note that if $`g=e^{x^aT_a}`$ then $`\mathrm{\Pi }^{ab}(g)=x^c\stackrel{~}{C}_c{}_{}{}^{ab}+O(x^2)`$ and similarly for $`\stackrel{~}{\mathrm{\Pi }}_{ab}`$. Drinfeld shows that the bivector $`\mathrm{\Pi }^{ab}T_aT_b`$ on $`G`$ defines a Poisson bracket that is compatible with the group multiplication law . A Poisson-Lie group is a Lie group with a Poisson structure which is compatible with the group operation. Thus we have that $`G`$ is a Poisson-Lie group. Note that the Poisson bivector is degenerate. Similarly the bivector $`\stackrel{~}{\mathrm{\Pi }}_{ab}\stackrel{~}{T}^a\stackrel{~}{T}^b`$ makes $`\stackrel{~}{G}`$ a Poisson Lie group. Klimcik and Severa discovered that the sigma model defined on the Poisson-Lie group $`G`$ is dual to the sigma model defined on the Poisson-Lie group $`\stackrel{~}{G}`$ hence the name Poisson-Lie duality. To exhibit the duality transformation we write down the symplectic structure on $`G_D`$. We note that the Drinfeld double is not a Poisson-Lie group in the Poisson structure associated with the symplectic structure since the Poisson bivector would be nondegenerate. In the (perfect) Drinfeld double every element $`kG_D`$ can uniquely be written as $`k=gu`$ or $`k=vh`$ where $`g,hG`$ and $`u,v\stackrel{~}{G}`$. The inverse function theorem shows that you can choose $`h`$ and $`u`$ as local coordinates on $`G_D`$ near the identity. Let $`\lambda =h^1dh`$ and $`\stackrel{~}{\lambda }=u^1du`$ be respectively the left invariant Maurer-Cartan forms on $`G`$ and $`\stackrel{~}{G}`$. The symplectic structure may be written as $`2\beta `$ $`=`$ $`(duu^1)_a(g^1dg)^a+(v^1dv)_a(dhh^1)^a,`$ (3.1) $`=`$ $`\stackrel{~}{\lambda }_a[\lambda ^b+\stackrel{~}{\lambda }_c\mathrm{\Pi }^{cb}(h^1)](I\stackrel{~}{\mathrm{\Pi }}(u^1)\mathrm{\Pi }(h^1))^1{}_{b}{}^{}^a`$ $`+`$ $`[\stackrel{~}{\lambda }_b+\lambda ^c\stackrel{~}{\mathrm{\Pi }}_{cb}(u^1)]\lambda ^a(I\mathrm{\Pi }(h^1)\stackrel{~}{\mathrm{\Pi }}(u^1))^1{}_{}{}^{b}{}_{a}{}^{}.`$ Using $`\beta `$ we can construct the duality transformations. All we have to verify is that we get metrics and antisymmetric tensor fields on $`G`$ and $`\stackrel{~}{G}`$. Klimcik and Severa show that the symmetric and antisymmetric parts of the rank two tensors $`E=\lambda ^t(F^1+\mathrm{\Pi })^1\lambda `$ and $`\stackrel{~}{E}=\stackrel{~}{\lambda }^t(F+\stackrel{~}{\mathrm{\Pi }})^1\stackrel{~}{\lambda }`$ are respectively the metrics and antisymmetric tensors for the dual sigma models on $`G`$ and $`\stackrel{~}{G}`$. The coefficients $`F_{ab}`$ are constants. By an appropriate choice of basis for $`𝔤`$ one can always choose $`F=I+b`$ where $`b`$ is antisymmetric. For pedagogical reasons we first look at the special case where $`b=0`$ which is analogous to choosing $`n^0=0`$ in (2.7). By looking at (3.1) we make an educated guess and conjecture that our orthonormal bases should be given by $`\theta `$ in $`G`$ and $`\stackrel{~}{\theta }`$ in $`\stackrel{~}{G}`$ where $$\lambda =(I+\mathrm{\Pi })\theta \text{and}\stackrel{~}{\lambda }=(I+\stackrel{~}{\mathrm{\Pi }})\stackrel{~}{\theta }.$$ (3.2) We can verify that this is in agreement with the Klimcik and Severa data by noting that $`E=\theta ^t(I\mathrm{\Pi })\theta `$ and $`\stackrel{~}{E}=\stackrel{~}{\theta }^t(I\stackrel{~}{\mathrm{\Pi }})\stackrel{~}{\theta }`$. The symmetric parts on $`E`$ and $`\stackrel{~}{E}`$ in this basis are respectively $`\theta ^t\theta `$ and $`\stackrel{~}{\theta }^t\stackrel{~}{\theta }`$ and thus we see that we have orthonormal bases on $`G`$ and $`\stackrel{~}{G}`$. Likewise we see that the components of the antisymmetric tensors in this basis are given by $`\mathrm{\Pi }`$ and $`\stackrel{~}{\mathrm{\Pi }}`$ respectively. Note that $`\theta `$ ($`\stackrel{~}{\theta }`$) is well defined on $`G`$ ($`\stackrel{~}{G}`$) because $`\mathrm{\Pi }`$ ($`\stackrel{~}{\mathrm{\Pi }}`$) is defined on $`G`$ ($`\stackrel{~}{G}`$). We are now ready to verify that there is a duality transformation. Postulate that frames given by $`\theta `$ and $`\stackrel{~}{\theta }`$ are the orthonormal ones we need. Rewrite (3.1) in the orthonormal frame where you find $`m`$ $`=`$ $`(I\stackrel{~}{\mathrm{\Pi }})(I\mathrm{\Pi }\stackrel{~}{\mathrm{\Pi }})^1(I+\mathrm{\Pi }),`$ (3.3) $`\stackrel{~}{l}`$ $`=`$ $`(I\stackrel{~}{\mathrm{\Pi }})\mathrm{\Pi }(I\stackrel{~}{\mathrm{\Pi }}\mathrm{\Pi })^1(I+\stackrel{~}{\mathrm{\Pi }}),`$ (3.4) $`l`$ $`=`$ $`(I\mathrm{\Pi })\stackrel{~}{\mathrm{\Pi }}(I\mathrm{\Pi }\stackrel{~}{\mathrm{\Pi }})^1(I+\mathrm{\Pi }),`$ (3.5) using the notation in (I-3.5). A brief computation shows that $`m=I+n`$ where $`n`$ is antisymmetric and given by $`n`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left[(\mathrm{\Pi }\stackrel{~}{\mathrm{\Pi }})^k(\stackrel{~}{\mathrm{\Pi }}\mathrm{\Pi })^k\right]`$ $`+`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(\mathrm{\Pi }\stackrel{~}{\mathrm{\Pi }})^k\mathrm{\Pi }{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\stackrel{~}{\mathrm{\Pi }}(\mathrm{\Pi }\stackrel{~}{\mathrm{\Pi }})^k.`$ In this frame $`m`$ is already in normal form and we can proceed. Note that $`\stackrel{~}{n}=n`$ as follows from (I-3.15). Using (I-3.8) and (I-3.9) we see that $`B=lm+I=\mathrm{\Pi }`$ and $`\stackrel{~}{B}=\stackrel{~}{l}+mI=\stackrel{~}{\mathrm{\Pi }}`$. The important result here is that $`B`$ and $`\stackrel{~}{B}`$ are quantities which respectively live on $`G`$ and $`\stackrel{~}{G}`$ and thus we have constructed the Poisson-Lie duality of Klimcik and Severa for the special case $`b=0`$. The general solution for arbitrary $`b`$ is given by choosing the orthonormal frames to be given by $$\lambda =(I+\mathrm{\Pi }F)\theta \text{and}\stackrel{~}{\lambda }=(F+\stackrel{~}{\mathrm{\Pi }})\stackrel{~}{\theta }.$$ (3.6) We have $`E=\theta ^t(FF^t\mathrm{\Pi }F)\theta `$ and $`\stackrel{~}{E}=\stackrel{~}{\theta }(F^t\stackrel{~}{\mathrm{\Pi }})\stackrel{~}{\theta }`$ with the components of the antisymmetric tensor fields given by $`B=bF^t\mathrm{\Pi }F`$ and $`\stackrel{~}{B}=b\stackrel{~}{\mathrm{\Pi }}`$ in this basis. The components of the symplectic form in this basis are $`m`$ $`=`$ $`(F^t\stackrel{~}{\mathrm{\Pi }})(I\mathrm{\Pi }\stackrel{~}{\mathrm{\Pi }})^1(I+\mathrm{\Pi }F),`$ (3.7) $`\stackrel{~}{l}`$ $`=`$ $`(F^t\stackrel{~}{\mathrm{\Pi }})\mathrm{\Pi }(I\stackrel{~}{\mathrm{\Pi }}\mathrm{\Pi })^1(F+\stackrel{~}{\mathrm{\Pi }}),`$ (3.8) $`l`$ $`=`$ $`(IF^t\mathrm{\Pi })\stackrel{~}{\mathrm{\Pi }}(I\mathrm{\Pi }\stackrel{~}{\mathrm{\Pi }})^1(I+\mathrm{\Pi }F),`$ (3.9) A brief computation shows that $`m=I+n`$ where $`n`$ is antisymmetric and given by $`n`$ $`=`$ $`b+{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left[F^t(\mathrm{\Pi }\stackrel{~}{\mathrm{\Pi }})^k(\stackrel{~}{\mathrm{\Pi }}\mathrm{\Pi })^kF\right]`$ $`+`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}F^t(\mathrm{\Pi }\stackrel{~}{\mathrm{\Pi }})^k\mathrm{\Pi }F{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\stackrel{~}{\mathrm{\Pi }}(\mathrm{\Pi }\stackrel{~}{\mathrm{\Pi }})^k.`$ Using (I-3.8) and (I-3.9) we see that $`B=lm+I=bF^t\mathrm{\Pi }F`$ and $`\stackrel{~}{B}=\stackrel{~}{l}+mI=b\stackrel{~}{\mathrm{\Pi }}`$. We have reproduced the ansatz of Klimcik and Severa and the duality transformation of Sfetsos. ## 4 Infinitesimally homogeneous $`n_{ij}`$ ### 4.1 General theory In this section we address the question, “What if $`n_{ij}`$ is the same everywhere?” We will see that this is a much weaker condition than saying $`n^{}=n^{\prime \prime }=0`$ which we already studied in Section I-8.2 and lead to abelian duality. We will show that $`P`$ is a homogeneous space under certain assumptions. First we have to address the question of what does “same everywhere” mean. The best way to do this is to exploit some ideas developed by Singer for the study of homogeneous spaces. It is convenient to work in the bundle $`(P)`$ of the adapted orthogonal frames we have been using. This bundle has structure group $`\mathrm{O}(n)`$ and it admits a global coframing given by $`(\theta ^i,\stackrel{~}{\theta }^j,\psi _{kl})`$ where $`(\theta ^i,\stackrel{~}{\theta }^j)`$ are the canonical $`1`$-forms on the frame bundle and $`\psi _{ij}`$ is an $`\mathrm{O}(n)`$ connection. Remember that the Maurer-Cartan form on $`\mathrm{O}(n)`$ is the restriction of $`\psi `$ to a fiber of $`(P)`$. The relationships among the geometries of $`M,\stackrel{~}{M},P`$ are encapsulated in the Cartan structural equations for $`(P)`$: $`d\theta ^i`$ $`=`$ $`\psi _{ij}\theta ^j{\displaystyle \frac{1}{2}}f_{ijk}\theta ^j\theta ^k,`$ (4.1) $`d\stackrel{~}{\theta }^i`$ $`=`$ $`\psi _{ij}\stackrel{~}{\theta }^j{\displaystyle \frac{1}{2}}\stackrel{~}{f}_{ijk}\stackrel{~}{\theta }^j\stackrel{~}{\theta }^k,`$ (4.2) $`d\psi _{ij}`$ $`=`$ $`\psi _{ik}\psi _{kj}T_{ijlm}^{\prime \prime }\theta ^l\stackrel{~}{\theta }^m.`$ (4.3) where $`f_{ijk}=f_{ikj}`$, $`\stackrel{~}{f}_{ijk}=\stackrel{~}{f}_{ikj}`$ and $`T_{ijkl}^{\prime \prime }=T_{jikl}^{\prime \prime }`$. Here $`f,\stackrel{~}{f},T^{\prime \prime }`$ are all functions on $`(P)`$. Note that there is torsion arising from the reduction of the structure group. The ideal generated by $`\{\theta ^1,\mathrm{},\theta ^n\}`$ is a differential ideal with integral submanifolds being the restriction of $`(P)`$ to the fibers of $`\mathrm{\Pi }`$. The ideal generated by $`\{\stackrel{~}{\theta }^1,\mathrm{},\stackrel{~}{\theta }^n\}`$ is a differential ideal with integral submanifolds being the restriction of $`(P)`$ to the fibers of $`\stackrel{~}{\mathrm{\Pi }}`$. The degenerate quadratic forms $`\theta ^i\theta ^i`$ and $`\stackrel{~}{\theta }^i\stackrel{~}{\theta }^i`$ on $`(P)`$ are respectively pullbacks of the metrics on $`M`$ and $`\stackrel{~}{M}`$. The pullback of the Riemannian connection on the frame bundle of $`M`$ is schematically $`\psi +f\theta `$ and likewise for $`\stackrel{~}{M}`$. Said differently, when restricting $`\psi _{ij}`$ to a “horizontal fiber” you get an orthogonal connection on the fiber with torsion, etc. We remind the reader that a tensor in $`P`$ is a collection of functions on $`(P)`$ which transform linearly under the action of $`\mathrm{O}(n)`$ on $`(P)`$, i.e., as you change frames the “tensor” transforms appropriately. For future use the frame dual to the coframe $`(\theta ^i,\stackrel{~}{\theta }^j,\psi _{kl})`$ will be denoted by $`(e_i,\stackrel{~}{e}_j,E_{kl})`$. The horizontal vector fields with respect to $`\psi `$ are spanned by $`\{e_A\}=\{e_i,\stackrel{~}{e}_j\}`$. We are now ready to define the statement “$`n_{ij}`$ is the same everywhere”. Pick a point $`b(P)`$. If we go to a rotated frame $`Rb`$, $`R\mathrm{O}(n)`$, then $`n_{ij}(b)`$ becomes $`n_{ij}(Rb)=R_{ik}R_{jl}n_{kl}(b)`$. Notice that as we move along the fiber going through $`b`$ we will get the full orbit of $`n_{ij}(b)`$ under $`\mathrm{O}(n)`$. Thus to make sense of “$`n_{ij}`$ is the same everywhere” we should not really talk about $`n_{ij}`$ but about the invariants of antisymmetric tensors under the orthogonal group. We should be thinking in terms of the space of orbits of antisymmetric tensors under $`\mathrm{O}(n)`$. In the frame bundle, the functions $`n_{ij}`$ define a map $`n:(P)^2(^n)`$. We say that $`P`$ is $`n`$-homogeneous if the image of the map $`n:(P)^2(^n)`$ is a single $`\mathrm{O}(n)`$-orbit. Said differently, you get the same $`2\times 2`$ block diagonalization of $`n_{ij}`$ at each point of $`P`$. This is a weaker condition than covariantly constant $`n`$. If $`P`$ is simply connected then a covariantly constant $`n_{ij}`$ is determined by parallel transporting $`n_{ij}`$ from a reference point. The value of $`n`$ at the reference point determines $`n`$ everywhere. The condition that $`P`$ be $`n`$-homogeneous is not strong enough for us. This leads to the notion of “infinitesimally $`n`$-homogeneous” where not only is $`n`$ the same everywhere but also the first $`(N+1)`$ covariant derivatives of $`n`$. Let $`n`$ denote the covariant derivative of $`n`$, $`^2n`$ the second covariant derivative of $`n`$, etc. Consider the map $$\rho ^s=(n,n,^2n,\mathrm{},^sn):(P)^{J_s},$$ where $`J_s`$ is an integer we do not compute. We say that $`P`$ is infinitesimally $`n`$-homogeneous if image of the map $`\rho ^{N+1}`$ is a single $`\mathrm{O}(n)`$-orbit. The integer $`N`$ is determined inductively as follows. First we do a rough argument and afterwards we state Singer’s result. Look at $`\rho ^0=n:(P)^2(^n)`$ and pick a point $`n^0`$ in the orbit. Consider $`=\{b(P):n(b)=n^0\}`$. Note that $``$ is a sub-bundle of $`(P)`$ because $`n((P))`$ is a single orbit. If $`K_0^{}\mathrm{O}(n)`$ is the isotropy group of $`n^0`$ then the action of $`K_0^{}`$ on a point $`b`$ leaves you in $``$. Thus $``$ is a principal sub-bundle of $`(P)`$ with structure group $`K_0^{}`$. The choice of $`n^0`$ has broken the symmetry group to $`K_0^{}`$. Now let us be precise about Singer’s result. Pick a $`b_0(P)`$. There exists a principal sub-bundle $`_0(P)`$ containing $`b_0`$ such that $`n`$ is constant on $`_0`$ and the structure group $`K_0\mathrm{O}(n)`$ of $`_0`$ is the connected component of the isotropy group of $`n(b_0)`$. Note that for a generic orbit, $`K_0`$ will be a maximal torus of $`\mathrm{O}(n)`$. Next we invoke $`n`$ to reduce the symmetry group some more. We use $`\rho ^1`$ and apply Singer’s theorem to it. There exists a principal sub-bundle $`_1_0`$ containing $`b_0`$ such that $`(n,n)`$ is constant on $`_1`$ and the structure group $`K_1K_0`$ of $`_1`$ is the connected component of the isotropy group of $`(n(b_0),n(b_0))`$. We continue the procedure by looking at $`\rho ^2,\rho ^3,\mathrm{}`$ and finding a sequence of principal sub-bundles $`(P)_0_1\mathrm{}_N_{N+1}`$ with respective structure groups $`\mathrm{O}(n)K_0K_1\mathrm{}K_NK_{N+1}`$. Since $`\mathrm{O}(n)`$ is finite dimensional there exists an integer $`N`$ such that the chain of groups satisfies the property that $`K_0K_1\mathrm{}K_{N1}K_N=K_{N+1}`$. In fact Singer establishes that $`_N=_{N+1}`$, henceforth denoted by $`H`$, is a principal bundle with structure group $`K=K_N`$. Our arguments show that $`\rho ^{N+1}=(n,n,\mathrm{},^{N+1}n)`$ is constant on $`H`$. Note that structure group $`K`$ is the connected component of the isotropy group of $`\rho ^N(b_0)`$. The chain of groups tells us that $`N\frac{1}{2}n(n1)`$ and later on we will see that we also require $`N1`$. Next we show that $`H`$ is a Lie group and conclude that $`P=H/K`$ is a homogeneous space. The strategy is to write down the Cartan structural equations for the principal bundle $`H`$ and show that they are actually the Maurer-Cartan equations for a group. Pick a point $`b(P)`$ and observe that $`d\rho ^N`$ is tangent to the orbit because $`\rho ^N((P))`$ is a single orbit. The orbit is generated by the action of $`\mathrm{O}(n)`$ therefore for $`XT_b(P)`$ there exists a linear map $`\mathrm{\Xi }:T_b(P)𝔰𝔬(n)`$ such that $`d\rho ^N(X)=\mathrm{\Xi }(X)\rho ^N(b)`$. Use the standard metric on $`\mathrm{SO}(n)`$ to write an orthogonal decomposition $`𝔰𝔬(n)=𝔨𝔨^{}`$ where $`𝔨`$ is the Lie algebra of $`K`$ and $`𝔨^{}`$ is its orthogonal complement. Let $`b_0H`$ then we observe that if $`\mathrm{\Xi },\mathrm{\Xi }^{}`$ are such that for $`XT_{b_0}(P)`$ you have $`d\rho ^N(X)=\mathrm{\Xi }(X)\rho ^N(b_0)=\mathrm{\Xi }^{}(X)\rho ^N(b_0)`$ then $`\mathrm{\Xi }^{}(X)\mathrm{\Xi }(X)𝔨`$. At $`b_0H`$ you can uniquely specify $`\mathrm{\Xi }`$ by requiring that $`\mathrm{\Xi }(X)𝔨^{}`$. We will make this choice. Note that we allow $`X`$ to be in the full tangent space $`T_{b_0}(P)`$. In summary, for $`b_0H(P)`$ there exists a unique linear transformation $`\mathrm{\Xi }:T_{b_0}(P)𝔨^{}`$ such that $$d\rho ^N(X)=\mathrm{\Xi }(X)\rho ^N(b_0).$$ (4.4) The definition of the covariant derivative is $`d\rho ^N(X)`$ $`=`$ $`\psi (X)\rho ^N(b_0)`$ (4.5) $`+`$ $`((_Xn)(b_0),_X(n)(b_0),_X(^2n)(b_0),\mathrm{},_X(^Nn)(b_0)).`$ Under the decomposition $`𝔰𝔬(n)=𝔨𝔨^{}`$ we have $`\psi =\psi ^𝔨+\psi ^𝔨^{}`$. Upon restriction to $`H`$, $`\psi ^𝔨`$ is a $`K`$-connection on the principal bundle $`H`$ and $`\psi ^𝔨^{}`$ will become torsion. Since $`K`$ is the isotropy group of $`\rho ^N(b_0)`$ we conclude that $`\psi ^𝔨\rho ^N(b_0)=0`$. If we restrict (4.5) to $`H(P)`$ and choose $`XT_{b_0}H`$ then $`d\rho ^N(X)=0`$ because $`\rho ^N`$ is constant on $`H`$. Thus for $`XT_{b_0}H`$ we have that $$\psi ^𝔨^{}(X)\rho ^N(b_0)=((_Xn)(b_0),_X(n)(b_0),_X(^2n)(b_0),\mathrm{},_X(^Nn)(b_0)).$$ (4.6) If we think of the above as a series of linear equations for $`\psi ^𝔨^{}(X)`$ then it is easy to see that if a solution exists then in must be unique. Next we show that the solution exists. To do this we observe that the covariant derivative (with connection $`\psi `$) of $`\rho ^N`$ in direction $`e_A`$ is given by $`d\rho ^N(e_A)=\mathrm{\Xi }(e_A)\rho ^N(b_0)`$, see (4.4), (4.5). Thus we have $$\mathrm{\Xi }(e_A)\rho ^N(b_0)=((_Xn)(b_0),_X(n)(b_0),_X(^2n)(b_0),\mathrm{},_X(^Nn)(b_0)).$$ Comparing this with (4.6) and using the uniqueness of the solution we conclude that the torsion $`\psi ^𝔨^{}(e_A)=\mathrm{\Xi }(e_A)`$. Also note that the right hand side of (4.6) is constant on $`H`$ and thus the $`\psi ^𝔨^{}(e_A)`$ must be constant on $`H`$ by uniqueness. On restriction to $`H`$ we have $$\psi _{ij}^𝔨^{}=\tau _{kij}^𝔨^{}\theta ^k+\stackrel{~}{\tau }_{kij}^𝔨^{}\stackrel{~}{\theta }^k,$$ (4.7) where $`\tau _{kij}^𝔨^{}`$ and $`\stackrel{~}{\tau }_{kij}^𝔨^{}`$ are constant on $`H`$. We are now almost ready to write down the Cartan structural equations for $`H`$. First we observe that certain functions are constant. If we let $`F`$ denote $`f_{ijk}`$ or $`\stackrel{~}{f}_{ijk}`$ then equations (I-8.17) and (I-8.18) may schematically be written $`(I+n)F=n`$. By differentiating we learn that $`^sF`$ is a function of $`(n,n,\mathrm{},^{s+1}n)`$ only. If $`N1`$ then $`f_{ijk},\stackrel{~}{f}_{ijk},T_{ijkl}^{\prime \prime }`$ in (4.1), (4.2) and (4.3) are constant on $`H`$ since $`\rho ^{N+1}`$ is constant on $`H`$. We require $`N1`$. The first Cartan structural equations (4.1) and (4.2) become $`d\theta ^i`$ $`=`$ $`\psi _{ij}^𝔨\theta ^j{\displaystyle \frac{1}{2}}c_{ijk}\theta ^j\theta ^k+\stackrel{~}{\tau }_{kij}^𝔨^{}\theta ^j\stackrel{~}{\theta }^k,`$ (4.8) $`d\stackrel{~}{\theta }^i`$ $`=`$ $`\psi _{ij}^𝔨\stackrel{~}{\theta }^j{\displaystyle \frac{1}{2}}\stackrel{~}{c}_{ijk}\stackrel{~}{\theta }^j\stackrel{~}{\theta }^k+\tau _{kij}^𝔨^{}\stackrel{~}{\theta }^j\theta ^k,`$ (4.9) where $`c,\stackrel{~}{c},\tau ^𝔨^{},\stackrel{~}{\tau }^𝔨^{}`$ are constant and $`c_{ijk}=c_{ikj}`$, $`\stackrel{~}{c}_{ijk}=\stackrel{~}{c}_{ikj}`$. Think of the $`ij`$ indices of $`T_{ijlm}^{\prime \prime }`$ as taking values in the $`𝔰𝔬(n)`$ Lie algebra and denote the projection of $`T_{ijlm}^{\prime \prime }`$ onto $`𝔨`$ by $`K_{ijlm}^𝔨`$. Note that $`K_{ijlm}^𝔨`$ is constant on $`H`$. The second Cartan structural equation (4.3) may be written as $`d\psi _{ij}^𝔨`$ $`=`$ $`\psi _{ik}^𝔨\psi _{kj}^𝔨(\psi _{ik}^𝔨\psi _{kj}^𝔨^{})^𝔨(\psi _{ik}^𝔨^{}\psi _{kj}^𝔨)^𝔨`$ (4.10) $``$ $`(\psi _{ik}^𝔨^{}\psi _{kj}^𝔨^{})^𝔨+K_{ijlm}^𝔨\theta ^l\stackrel{~}{\theta }^m,`$ where you substitute (4.7) for $`\psi ^𝔨^{}`$ in the above. The important lesson is that $`H`$ has a coframing given by $`(\theta ,\stackrel{~}{\theta },\psi ^𝔨)`$ and that the structural equations (4.8), (4.9) and (4.10) only involve constants and thus are the Maurer-Cartan equations for a Lie group. We have shown that if $`P`$ is infinitesimally $`n`$-homogeneous with $`N1`$ then $`P`$ is a homogeneous space $`H/K`$ where the Lie group $`H`$ is a sub-bundle of the frame bundle $`(P)`$. ### 4.2 The case of $`K=\{e\}`$ This is the situation where the residual symmetry group $`K`$ is broken all the way down to the identity group $`\{e\}`$. In this case $`P=H`$, the symplectic manifold $`P`$ is a Lie group, and $`𝔨^{}=𝔰𝔬(n)`$. The Maurer-Cartan equations for $`P`$ are $`d\theta ^i`$ $`=`$ $`{\displaystyle \frac{1}{2}}c_{ijk}\theta ^j\theta ^k+\stackrel{~}{\tau }_{kij}\theta ^j\stackrel{~}{\theta }^k,`$ (4.11) $`d\stackrel{~}{\theta }^i`$ $`=`$ $`{\displaystyle \frac{1}{2}}\stackrel{~}{c}_{ijk}\stackrel{~}{\theta }^j\stackrel{~}{\theta }^k+\tau _{kij}\stackrel{~}{\theta }^j\theta ^k,`$ (4.12) where $`c_{ijk}=c_{ikj}`$, $`\stackrel{~}{c}_{ijk}=\stackrel{~}{c}_{ikj}`$. Also $`\tau _{kij}`$ and $`\stackrel{~}{\tau }_{kij}`$ are antisymmetric under $`ij`$ for arbitrary $`i,j`$ reflecting that $`𝔨^{}=𝔰𝔬(n)`$. Note that $`\stackrel{~}{\theta }`$ generates a differential ideal and thus $`\stackrel{~}{\theta }=0`$ defines a fibration $`\stackrel{~}{\mathrm{\Pi }}:P\stackrel{~}{M}`$ with fibers isomorphic to a Lie group $`G`$ with structure constants $`c_{ijk}`$. Likewise, $`\theta `$ generates a differential ideal and thus $`\theta =0`$ defines a fibration $`\mathrm{\Pi }:PM`$ with fibers isomorphic to a Lie group $`\stackrel{~}{G}`$ with structure constants $`\stackrel{~}{c}_{ijk}`$. We also remark that $`G`$ and $`\stackrel{~}{G}`$ are Lie subgroups of $`P`$. If $`\tau =0`$ then $`G`$ is a normal subgroup of $`P`$. If $`\stackrel{~}{\tau }=0`$ then $`\stackrel{~}{G}`$ is a normal subgroup of $`P`$. Let $`(e_i,\stackrel{~}{e}_j)`$ be the basis dual to $`(\theta ^i,\stackrel{~}{\theta }^j)`$. The Maurer-Cartan equations may be reformulated as the Lie algebra relations $`[e_i,e_j]`$ $`=`$ $`c_{kij}e_k`$ (4.13) $`[\stackrel{~}{e}_i,\stackrel{~}{e}_j]`$ $`=`$ $`\stackrel{~}{c}_{kij}\stackrel{~}{e}_k`$ (4.14) $`[e_i,\stackrel{~}{e}_j]`$ $`=`$ $`\tau _{ikj}\stackrel{~}{e}_k\stackrel{~}{\tau }_{jki}e_k.`$ (4.15) All the statements made in the previous paragraph also follow from the above. Consider the left invariant vector fields $`X=X^ie_i`$, $`Y=Y^ie_i`$. Note that $`_X\stackrel{~}{\theta }^i=X^k\tau _{kij}\stackrel{~}{\theta }^j`$. Let $`\tau _{ij}(X)=X^k\tau _{kij}`$ then using the identity $`[_X,_Y]\stackrel{~}{\theta }^i=_{[X,Y]}\stackrel{~}{\theta }^i`$ you obtain $`[\tau (X),\tau (Y)]=\tau ([X,Y])`$. Thus we have a Lie algebra representation $`\tau :𝔤𝔰𝔬(n)`$. This means we have a representation of $`G`$ by real orthogonal $`n\times n`$ matrices. Likewise, $`\stackrel{~}{\tau }:\stackrel{~}{𝔤}𝔰𝔬(n)`$ is a Lie algebra representation and we have a representation of $`\stackrel{~}{G}`$ by real orthogonal matrices. This does not mean that $`G`$ is a compact group if $`\tau 0`$. A comment made in Appendix A tells us that $`𝔤/(\mathrm{ker}\tau )`$ is a Lie algebra of the form “compact semisimple \+ abelian”. We know nothing at all about the ideal $`\mathrm{ker}\tau `$ so we cannot make a stronger statement about the structure of $`𝔤`$. Similar remarks apply to $`\stackrel{~}{G}`$. A Drinfeld double $`G_D`$ admits the following geometric characterization: 1. It is a Lie group of dimension $`2n`$ with a bi-invariant quadratic form of type $`(n,n)`$. 2. It is a bifibration with the property that the fibers are isotropic submanifolds of $`G_D`$. The fibers are also isomorphic to Lie groups $`G`$ and $`\stackrel{~}{G}`$. If we apply the above to our situation by requiring that the quadratic form $`\stackrel{~}{\theta }^i\theta ^i+\theta ^i\stackrel{~}{\theta }^i`$ be bi-invariant then we learn that $`P=H`$ is a Drinfeld double, $`\tau _{ijk}=c_{ijk}`$, $`\stackrel{~}{\tau }_{ijk}=\stackrel{~}{c}_{ijk}`$, $`\tau _{kij}=\tau _{ikj}`$ and $`\stackrel{~}{\tau }_{kij}=\stackrel{~}{\tau }_{ikj}`$. It follows that both $`c_{ijk}`$ and $`\stackrel{~}{c}_{ijk}`$ are totally antisymmetric and thus the quadratic forms $`\theta ^i\theta ^i`$ and $`\stackrel{~}{\theta }^i\stackrel{~}{\theta }^i`$ give bi-invariant positive definite metrics on $`G`$ and $`\stackrel{~}{G}`$ respectively. Thus $`G`$ and $`\stackrel{~}{G}`$ are of the type “compact semisimple + abelian”. The symplectic structure on $`P`$ (if it exists) appears to be different than the standard symplectic structure on the Drinfeld double, see (3.1). In the examples I am familiar, if $`G`$ is simple and compact the its dual $`\stackrel{~}{G}`$ constructed via $`R`$-matrices is neither simple nor compact. I do not know what is known in this more general case “compact semisimple + abelian”. Returning to the general case we remark that the equations $`d^2\theta =0`$ and $`d^2\stackrel{~}{\theta }=0`$ leads to the conclusion that $`c_{ijk}`$ and $`\stackrel{~}{c}_{ijk}`$ are respectively the structure constants for Lie groups $`G`$ and $`\stackrel{~}{G}`$, $`\tau :𝔤𝔰𝔬(n)`$ and $`\stackrel{~}{\tau }:\stackrel{~}{𝔤}𝔰𝔬(n)`$ are Lie algebra representations and $`\tau _{kjm}\stackrel{~}{\tau }_{jil}\tau _{ljm}\stackrel{~}{\tau }_{jik}+c_{ijk}\stackrel{~}{\tau }_{mjl}c_{ijl}\stackrel{~}{\tau }_{mjk}+c_{jkl}\stackrel{~}{\tau }_{mij}`$ $`=`$ $`0,`$ (4.16) $`\stackrel{~}{\tau }_{kjm}\tau _{jil}\stackrel{~}{\tau }_{ljm}\tau _{jik}+\stackrel{~}{c}_{ijk}\tau _{mjl}\stackrel{~}{c}_{ijl}\tau _{mjk}+\stackrel{~}{c}_{jkl}\tau _{mij}`$ $`=`$ $`0.`$ (4.17) These equations are generalizations of the corresponding equations (B.4) in the bialgebra case. These equations follow just from the structure equations for the group $`P=H`$. There are additional constraints which follow from duality considerations such as $`d\gamma =H\stackrel{~}{H}`$ which lead to $`c_{ijk}n_{il}+c_{ikl}n_{ij}+c_{ilj}n_{ik}`$ $`=`$ $`+H_{jkl},`$ (4.18) $`\stackrel{~}{c}_{ijk}n_{il}+\stackrel{~}{c}_{ikl}n_{ij}+\stackrel{~}{c}_{ilj}n_{ik}`$ $`=`$ $`\stackrel{~}{H}_{jkl},`$ (4.19) $`m_{jl}\tau _{kli}m_{kl}\tau _{jli}+n_{jl}\stackrel{~}{\tau }_{ilk}n_{kl}\stackrel{~}{\tau }_{ilj}`$ $`=`$ $`m_{il}c_{ljk},`$ (4.20) $`m_{jl}\stackrel{~}{\tau }_{kli}m_{kl}\stackrel{~}{\tau }_{jli}+n_{jl}\tau _{ilk}n_{kl}\tau _{ilj}`$ $`=`$ $`+\stackrel{~}{c}_{ljk}m_{li}.`$ (4.21) We do not know if there are non-trivial solutions to these equations. ## Acknowledgments I would like to thank O. Babelon, L. Baulieu, T. Curtright, L.A. Ferreira, D. Freed, S. Kaliman, C-H Liu, R. Nepomechie, N. Reshetikhin, J. Sánchez Guillén, N. Wallach and P. Windey for discussions on a variety of topics. I would also like to thank Jack Lee for his Mathematica package Ricci that was used to perform some of the computations. I am particularly thankful to R. Bryant and I.M. Singer for patiently answering my many questions about differential geometry. This work was supported in part by National Science Foundation grant PHY–9870101. Appendices ## Appendix A Some Lie groups facts For clarification purposes we make some remarks about the left and right actions on a Lie group. For notational simplicity we restrict to matrix Lie groups. The identity element will be denoted by $`I`$. Let $`G`$ be a Lie group. Let $`aG`$ then the left and right actions on $`G`$ are respectively defined by $`L_ag=ag`$ and $`R_ag=ga`$ for $`gG`$. We take $`𝔤`$, the Lie algebra of $`G`$, to be the left invariant vector fields and we identify it with the tangent space at the identity $`T_IG`$. The left invariant Maurer-Cartan forms are $`\theta =g^1dg`$. They satisfy the Maurer-Cartan equations $`d\theta =\theta \theta `$. Pick a basis $`(e_1,\mathrm{},e_n)`$ of left invariant vector fields for $`𝔤`$ with bracket relations $`[e_i,e_j]=f^k{}_{ij}{}^{}e_{k}^{}`$. If the dual basis of left invariant forms is $`(\theta ^1,\mathrm{},\theta ^n)`$ then $`d\theta ^i=\frac{1}{2}f^i{}_{jk}{}^{}\theta _{}^{j}\theta ^k`$. Naively you would expect that if $`X`$ is a left invariant vector field then $`_X\theta =0`$ since $`\theta `$ is left invariant. In fact a brief computation shows that $`_X\theta ^i=(X^kf^i{}_{kj}{}^{})\theta ^j`$ which is the adjoint action. The answer to this conundrum is that the left invariant vector fields generate the right group action. The easiest way to see this is to use old fashioned differentials. Let $`vT_IG`$ then the left invariant vector field $`X`$ at $`g`$ which is $`v`$ at the identity is given by $`X=gv`$. The infinitesimal action of this vector field at $`g`$ is given by $`gg+ϵX=g+ϵgv=g(I+ϵv)ge^{ϵv}`$ which is the right action of the group. Thus we see that $`g^1dge^{ϵv}(g^1dg)e^{ϵv}`$ which is the adjoint action in accordance with the Lie derivative computation. Take any metric $`h_{ij}`$ at the identity then $`h_{ij}\theta ^i\theta ^j`$ is a left invariant metric on $`G`$. In general this metric is not invariant under the right action of the group. The right invariance condition is $`h_{il}f^l{}_{jk}{}^{}+h_{jl}f^l{}_{ik}{}^{}=0`$ which means that the structure constants are totally antisymmetric if the indices are lowered using $`h_{ij}`$. In such a situation the metric is bi-invariant. Assume you have a Lie algebra $`𝔤`$ with an invariant positive definite metric then you have an orthogonal decomposition into ideals $`𝔤=𝔨𝔞`$ where $`𝔨`$ is a compact semisimple Lie algebra and $`𝔞`$ is abelian. The proof is straightforward and involves putting together a variety of observations. The invariance of the inner product $`(,)`$ is the statement that the adjoint representation $`\mathrm{ad}_XY=[X,Y]`$ is skew adjoint with respect to the metric: $`(\mathrm{ad}_XY,Z)+(Y,\mathrm{ad}_XZ)=0`$. The skew adjointness immediately leads to the decomposition of $`𝔤`$ into irreducible pieces. It is an elementary exercise in linear algebra to show that if $`𝔥`$ is a non-trivial ideal in $`𝔤`$ then its orthogonal complement $`𝔥^{}`$ is also an ideal. Since $`𝔤=𝔥𝔥^{}`$ we conclude that by continuing this process the Lie algebra decomposes into irreducible ideals $`𝔤=𝔤_1𝔤_2\mathrm{}𝔤_N`$. Let us collate all the abelian subalgebras together into $`𝔞`$ and rewrite the decomposition as $`𝔤=𝔨_1𝔨_2\mathrm{}𝔨_M𝔞`$ where each $`𝔨_j`$ is a simple lie algebra. Since this is a decomposition into ideals we have that $`𝔞`$ is the center of the Lie algebra. Let $`(,)_j`$ be the restriction of the invariant inner product to $`𝔨_j`$. An application of Schur’s Lemma tells us that an invariant bilinear form on a simple Lie algebra is a multiple of the Killing form. Thus we conclude that $`(,)_j=\lambda K_j(,)`$ where $`\lambda `$ is a non-zero scalar and $`K_j`$ is the Killing form<sup>3</sup><sup>3</sup>3The sign of the Killing form is chosen such that it is positive on a compact simple Lie algebra. on $`𝔨_j`$. Since $`(,)`$ is positive definite, the Killing form $`K_j`$ must be definite and this is only possible if the Lie algebra $`𝔨_j`$ is of compact type. This concludes the proof of the proposition in the opening sentence. Closely related to the above is the following. If a Lie algebra $`𝔤`$ has a faithful representation $`\tau :𝔤𝔰𝔬(n)`$ by skew adjoint matrices then $`𝔤=𝔨𝔞`$ where $`𝔨`$ is a compact semisimple Lie algebra and $`𝔞`$ is abelian. The proof follows from the observation that because the representation is faithful we can think of $`𝔤`$ as a matrix Lie subalgebra of $`𝔰𝔬(n)`$. We know that $`𝔰𝔬(n)`$ has a positive definite invariant metric so restriction to $`𝔤`$ induces a positive definite invariant metric on $`𝔤`$. We now use the proposition from the previous paragraph. We remark that if the representation $`\tau `$ in the previous paragraph is not faithful then $`(\mathrm{ker}\tau )𝔤`$ is a nontrivial ideal in $`𝔤`$. The Lie algebra $`𝔤/(\mathrm{ker}\tau )`$ is of the form $`𝔨𝔞`$ but we can say nothing about the Lie algebra $`(\mathrm{ker}\tau )`$. It is easy to see that the space of all invariant positive definite metrics on a Lie algebra is a convex set. In a simple Lie algebra $`𝔤`$, the third cohomology group $`H^3(𝔤)`$ is one dimensional and generated by the three form $`\omega (X,Y,Z)=K(X,[Y,Z])`$ where $`K`$ is the killing form which may be written in terms of the structure constants as $`\omega _{ijk}=K_{il}f^l_{jk}`$. If $`h`$ is an invariant metric then $`h(X,[Y,Z])`$ is also a closed three form and by the convexity of the space of positive definite invariant metrics it must be in the same cohomology class as $`\omega `$. ## Appendix B A primer on classical $`R`$-matrices There are two main nonequivalent approaches to classical $`R`$-matrices. The more familiar one is due to Drinfeld and based on the study of Lie bialgebras . The other due to Semenov-Tian-Shansky is based on double Lie algebras is the one directly related to our work. Here we discuss the interconnections between these two approaches. For an introduction to $`R`$-matrices and Poisson-Lie groups see the book by Chari and Pressley or the article . ### B.1 The Drinfeld Approach Drinfeld bases his approach on the notion of a Lie bialgebra. We begin with a down to earth approach. Assume you have a Lie algebra $`𝔤`$ with basis $`(e_1,\mathrm{},e_n)`$ and Lie bracket relations $`[e_a,e_b]=f^c{}_{ab}{}^{}e_{c}^{}`$. If $`X,Y𝔤`$ then the adjoint action by $`X`$ is given by $`\mathrm{ad}_XY=[X,Y]`$. The adjoint action extends naturally to tensor products of $`𝔤`$. Let $`𝔤^{}`$ be the vector space dual with corresponding basis $`(\omega ^1,\mathrm{},\omega ^n)`$. The Lie algebra $`𝔤`$ acts on $`𝔤^{}`$ via the coadjoint representation $`\mathrm{ad}_{e_a}^{}\omega ^b=f^b{}_{ac}{}^{}\omega _{}^{c}`$. In general $`𝔤^{}`$ is not a Lie algebra but there is a natural Lie algebra structure on $`𝔤𝔤^{}`$ given by $`[e_a,e_b]`$ $`=`$ $`f^c{}_{ab}{}^{}e_{c}^{},`$ $`[e_a,\omega ^b]`$ $`=`$ $`f^b{}_{ac}{}^{}\omega _{}^{c}`$ $`[\omega ^a,\omega ^b]`$ $`=`$ $`0.`$ The most famous example is combining $`𝔤=𝔰𝔬(3)`$ and its Lie algebra dual $`𝔤^{}^3`$ into the Lie algebra of the euclidean group $`E(3)`$. The situation becomes much more interesting when $`𝔤^{}`$ is a Lie algebra in its own rights $`[\omega ^a,\omega ^b]=\widehat{f}_c{}_{}{}^{ab}\omega _{}^{c}`$. You observe that $`𝔤^{}`$ acts via its coadjoint action on its dual $`(𝔤^{})^{}𝔤`$. Thus one can consider the following more symmetric structure which takes into account the respective coadjoint actions $`[e_a,e_b]`$ $`=`$ $`f^c{}_{ab}{}^{}e_{c}^{},`$ (B.1) $`[e_a,\omega ^b]`$ $`=`$ $`\widehat{f}_a{}_{}{}^{bc}e_{c}^{}f^b{}_{ac}{}^{}\omega _{}^{c},`$ (B.2) $`[\omega ^a,\omega ^b]`$ $`=`$ $`\widehat{f}_c{}_{}{}^{ab}\omega _{}^{c}.`$ (B.3) This will be a Lie algebra if the following conditions are satisfied: $$f^e{}_{ab}{}^{}\widehat{f}_{e}^{}{}_{}{}^{cd}=\widehat{f}_b{}_{}{}^{ed}f_{}^{c}{}_{ae}{}^{}+\widehat{f}_b{}_{}{}^{ce}f_{}^{d}{}_{ae}{}^{}\widehat{f}_a{}_{}{}^{ed}f_{}^{c}{}_{be}{}^{}\widehat{f}_a{}_{}{}^{ce}f_{}^{d}{}_{be}{}^{}.$$ (B.4) According to Drinfeld a Lie bialgebra is a Lie algebra with Lie brackets (B.1), (B.2) and (B.3). Drinfeld gives a more abstract formulation which is more suitable for studying the abstract properties of a bialgebra and seeing the origins of classical $`R`$-matrices. Assume you have a Lie algebra $`𝔤`$ and a “cobracket” $`\mathrm{\Delta }:𝔤^2𝔤`$. Drinfeld requires that the cobracket defines a Lie algebra on $`𝔤^{}`$. The structure constants on $`𝔤^{}`$ are related to the cobracket by $`\mathrm{\Delta }(e_a)=\frac{1}{2}\widehat{f}_a{}_{}{}^{bc}e_{b}^{}e_c`$. Compatibility condition (B.4) is incorporated via a cohomological argument. The complex in question is $$^2𝔤𝔤^{}^2𝔤\stackrel{^{}}{}^2𝔤^{}^2𝔤\mathrm{}.$$ The coboundary operator (differential) is given by $$(^{}\mathrm{\Delta })(X,Y)=\mathrm{ad}_X(\mathrm{\Delta }(Y))\mathrm{ad}_Y(\mathrm{\Delta }(X))\mathrm{\Delta }([X,Y]).$$ (B.5) The condition (B.4) that glues the Lie algebras into a Lie bialgebra is seen to be equivalent to the cocycle condition $`^{}\mathrm{\Delta }=0`$. In this language, a Lie bialgebra is a Lie algebra $`𝔤`$ along with a cobracket $`\mathrm{\Delta }`$ such that $`𝔤^{}`$ a Lie algebra and the cobracket is a $`1`$-cocycle. A Lie bialgebra is exact if the cocycle is exact. This means that there exists a $`r^2𝔤`$ such that $`\mathrm{\Delta }(X)=^{}r=\mathrm{ad}_Xr`$. A computation shows that $`r`$ defines a bialgebra structure if and only if the Schouten bracket $`[[r,r]]`$ is $`\mathrm{ad}(𝔤)`$-invariant: $`[[r,r]](^3𝔤)_{\mathrm{inv}}`$. The Schouten bracket is defined by $$[[WX,YZ]]=[W,Y]XZ[W,Z]XY[X,Y]WZ+[X,Z]WY.$$ The condition $`[[r,r]](^3𝔤)_{\mathrm{inv}}`$ is called the modified classical Yang-Baxter equation (MCYBE) and $`[[r,r]]=0`$ is called the classical Yang-Baxter equation (CYBE). If $`𝔤`$ is semisimple then the Whitehead lemma states that $`H^1(𝔤,^2𝔤)=0`$ and thus the cocycle $`\mathrm{\Delta }`$ is always a coboundary $`\mathrm{\Delta }=^{}r`$. Thus in this case we need to understand the set of all $`r^2𝔤`$ which satisfy MCYBE. If $`𝔤`$ is simple then $`(^3𝔤)_{\mathrm{inv}}`$ is one dimensional and is generated by the three index tensor obtained by raising two indices on the structure constants using the Killing metric. If we call this object $`B_\mathrm{K}`$ then MCYBE becomes $`[[r,r]]=aB_\mathrm{K}`$ for some constant $`a`$. Let us work in a basis. If $`r=\frac{1}{2}r^{ab}e_ae_b`$. Then $`\mathrm{\Delta }(e_a)=\mathrm{ad}_{e_a}r=\frac{1}{2}r^{bc}\mathrm{ad}_{e_a}(e_be_c)=\frac{1}{2}(r^{dc}f^b{}_{ad}{}^{}+r^{bd}f^c{}_{ad}{}^{})e_b`$. This tells us that $$\widehat{f}_a{}_{}{}^{bc}=r^{dc}f^b{}_{ad}{}^{}+r^{bd}f^c{}_{ad}{}^{}.$$ (B.6) Note that $`\widehat{f}_a{}_{}{}^{bc}=\widehat{f}_a^{cb}`$ because $`r`$ is skew. It is possible to generalize the above by allowing the cocycle (now called $`C`$) to be in $`𝔤𝔤`$. If you write $`C=C^{ab}e_ae_b`$ and you let $`\mathrm{\Delta }(X)=\mathrm{ad}_XC`$. In this case you find you get a Lie bialgebra if the following two conditions are satisfied: 1. $`(C^{ab}+C^{ba})e_ae_b`$, the symmetric part of $`C`$, is $`\mathrm{ad}(𝔤)`$-invariant, 2. $`[[C,C]][C^{12},C^{13}]=[C^{12},C^{23}]+[C^{13},C^{23}]𝔤𝔤𝔤`$ is $`\mathrm{ad}(𝔤)`$-invariant. We use standard quantum group notation where for example $`C^{13}=C^{ab}e_aIe_b`$, etc. The bracket $`[[,]]`$ above reduces to the Schouten bracket if $`C`$ is skew symmetric. We remark that the equation $`[[C,C]]=0`$ is also called the classical Yang-Baxter equation. $`C`$ or $`r`$ are called classical $`R`$-matrices (by Drinfeld). The modified classical Yang-Baxter equation is $`[[C,C]](𝔤𝔤𝔤)_{\mathrm{inv}}`$. A brief computation shows that $$\widehat{f}_a{}_{}{}^{bc}=C^{dc}f^b{}_{ad}{}^{}+C^{bd}f^c{}_{ad}{}^{}.$$ (B.7) Let us write $`C=\frac{1}{2}s^{ab}(e_ae_b+e_be_a)+\frac{1}{2}r^{ab}e_ae_b`$ where $`s`$ is symmetric and $`r`$ is antisymmetric. Since $`s^{ab}`$ is an $`\mathrm{ad}(𝔤)`$-invariant tensor we have that $$\widehat{f}_a{}_{}{}^{bc}=r^{dc}f^b{}_{ad}{}^{}+r^{bd}f^c{}_{ad}{}^{}.$$ (B.8) Thus the Lie algebra structure on $`𝔤^{}`$ only depends on the antisymmetric part of $`C`$. Note that $`\widehat{f}_a^{bc}`$ will be skew under $`bc`$ as required. The effect of the symmetric part $`s_{ab}`$ is to change the $`\mathrm{ad}𝔤`$-invariant term in the right hand side of the MCYBE. Equation (B.7) may be interpreted as giving a Lie algebra homomorphism $`C:𝔤^{}𝔤`$. ### B.2 Semenov-Tian-Shansky approach Semenov-Tian-Shansky’s approach is directly influenced by classical integrable models where he needs that a single Lie algebra admits two different Lie brackets. Let $`𝔤`$ be a Lie algebra and let $`R:𝔤𝔤`$ be a linear transformation (not necessarily invertible). Define a skew operation $`[,]_R`$ by $$[X,Y]_R=[RX,Y]+[X,RY].$$ (B.9) If $`[,]_R`$ is a Lie bracket then $`R`$ is called a classical R-matrix by Semenov-Tian-Shansky. The Jacobi identity for $`[,]_R`$ may be written as $$[X,[RY,RZ]R([Y,Z]_R)]+\text{cyclic permutations}=0$$ A Lie algebra $`𝔤`$ with two Lie brackets $`[,]`$ and $`[,]_R`$ is called a double Lie algebra by Semenov-Tian-Shansky. The equation is $`[RY,RZ]R([Y,Z]_R)=0`$ is also called the CYBE. The equation $$[RY,RZ]R([Y,Z]_R)=c[Y,Z]$$ (B.10) where $`c`$ is a constant is also called the MCYBE. Solutions to either of these satisfy the Jacobi identity displayed above. In a basis we have that $`Re_a=e_bR^b_a`$, $`[e_a,e_b]_R=(f^c{}_{ab}{}^{})_Re_c`$ and consequently (B.9) becomes $`[e_a,e_b]_R=(R^d{}_{a}{}^{}f_{}^{c}{}_{db}{}^{}+R^d{}_{b}{}^{}f_{}^{c}{}_{ad}{}^{})e_c`$. Thus the new structure constants are $$(f^c{}_{ab}{}^{})_R=R^d{}_{a}{}^{}f_{}^{c}{}_{db}{}^{}+R^d{}_{b}{}^{}f_{}^{c}{}_{ad}{}^{}.$$ (B.11) The Semenov-Tian-Shansky version of the MCYBE (B.10) is $$f^c{}_{de}{}^{}R_{}^{d}{}_{a}{}^{}R_{}^{e}{}_{b}{}^{}f^e{}_{db}{}^{}R_{}^{d}{}_{a}{}^{}R_{}^{c}{}_{e}{}^{}f^e{}_{ad}{}^{}R_{}^{d}{}_{b}{}^{}R_{}^{c}{}_{e}{}^{}=cf^c{}_{ab}{}^{}.$$ (B.12) ### B.3 Relating Drinfeld and Semenov-Tian-Shansky To relate the Semenov-Tian-Shansky approach and the Drinfeld approach one needs an $`\mathrm{ad}(𝔤)`$-invariant metric on $`𝔤`$. If $`𝔤`$ is semisimple then one can take the $`\mathrm{ad}(𝔤)`$-invariant metric to be the Killing metric. The $`\mathrm{ad}(𝔤)`$-invariant metric is used to identify $`𝔤`$ with $`𝔤^{}`$. By lowering indices $`f_{abc}`$ is completely antisymmetric (due to the $`\mathrm{ad}(𝔤)`$-invariance). We wish to identify $`f_R`$ with $`\stackrel{~}{f}`$. Note that by rearranging indices we have $`(f_c{}_{}{}^{ab})_R`$ $`=`$ $`R^{da}f_{cd}{}_{}{}^{b}+R^{db}f_c{}_{}{}^{a}_d`$ (B.13) $`=`$ $`R^{da}f^b{}_{cd}{}^{}R^{db}f^a_{cd}`$ If we are in a situation where $`R^{ab}=r^{ab}`$ then we have related $`f_R`$ to $`\widehat{f}`$. Said differently $`R:𝔤𝔤`$ is a skew-adjoint operator with respect to the invariant metric. In fact there is a theorem which states that if $`𝔤`$ has an $`\mathrm{ad}(𝔤)`$-invariant metric and if $`R:𝔤𝔤`$ is skew-adjoint then the double Lie algebra is isomorphic to a Lie bialgebra and all the structures in the Semenov-Tian-Shansky approach (CYBE, MYBE) go into the structures in the Drinfeld approach (CYBE, MYBE). The isomorphism is given by thinking of the metric as giving a map $`𝔤𝔤^{}`$, i.e., lowering/raising indices. We are in a different situation because not all our $`R`$-matrices are skew adjoint.
warning/0003/hep-th0003259.html	ar5iv	text	# References POINCARE ANOMALY IN PLANAR FIELD THEORY Subir = Ghosh<sup>1</sup><sup>1</sup>1email:sghosh@isical.ac.in;subir@boson.bose.res.in P. A. M. U., Indian Statistical Institute, 203 B. T. Road, Calcutta 700035, India. Abstract: We show the presence of Poincare anomaly in Maxwell-Chern-Simons theory = with an explicit mass term, in 2+1-dimensions. The ubiquitous use of 2+1-dimensional field theories in condensed matter = systems, where the dynamics normal to a plane is severely restricted, = has enlarged the scope of lower dimensional physics from being just a = toy model of the 3+1-dimensional world. The topologically massive Maxwell-Chern-Simons (MCS) gauge theory was first thoroughly analysed by Deser, Jackiw and Templeton (DJT) in = their seminal work , where the subtle interplay between = Poincare invariance and an unambiguous determination of the spin of the = excitations in a vector theory was revealed. It was shown that correct = space-time transformation of the gauge invariant observables, such as = electric and magnetic fields, were induced by Poincare generators which = obey an anomalous algebra among themselves. However, a phase = redefinition of the creation and annihilation operators removed the = commutator anomaly and yielded the spin contribution in a single stroke. = In the present Letter, we consider the MCS model with an explicit mass = term, i.e. the MCS-Proca (MCSP) model. We show that in the = presence of two mass scales, the topological one $`(\mu )`$, generated by = the Chern-Simons term, and the explicit (gauge symmetry breaking one) = one $`(m)`$, the anomaly in the Poincare transformationsin the = eletromagnetic fields can not be removed, even though the equations of = motion are manifestly Lorentz covariant. This is our main result. The MCSP model was studied previously in . It also appears = naturally in the large fermion mass limit of the bosonization of gauged = massive Thirring model . In , it was argued that the = self-dual factorisation of the equation of motion leads to two self and = anti-self dual excitations of different masses, thereby accounting for = the parity violation induced by the topological term. Recent results = indicate, in the path integral formalism, that the above naive = conclusions are invalid. As shown in , the fact that MCSP model = is a result of a fusion between self and anti-self dual models explains = the self dual factorisation of the equation of motion. But in the = process, the self and anti-self dual property of the MCSP model is no = longer manifest. This controversy demands an indepth analysis of the model. The MCSP model, with the metric being $`g_{\mu \nu }=3Ddiag(+),ϵ=_{12}=3D1`$, is $$_{MCSP}=3D\frac{1}{4}A_{\mu \nu }A^{\mu \nu }+\frac{\mu }{4}ϵ_{\mu \nu \lambda }A^{\mu \nu }A^\lambda +\frac{m^2}{2}A_\mu A^\mu ,A_{\mu \nu }=3D_\mu A_\nu =_\nu A_\mu .$$ (1) Taking $`m^2=3D0`$ reproduces the MCS theory, which being a gauge theory = is amenable to gauge fixing conditions. This simplifies the model = considerably and makes the field content transparant. We also try to = implement similar parametrizations as in and hence convert = the above gauge non-invariant theory to a gauge invariant one by the = Stuckelberg prescription, $$_{St}=3D\frac{1}{4}A_{\mu \nu }A^{\mu \nu }+\frac{\mu }{=4}ϵ_{\mu \nu \lambda }A^{\mu \nu }A^\lambda +\frac{m^2}{2}(A_\mu =_\mu \theta )(A^\mu ^\mu \theta ),$$ (2) where $`\theta `$ is the Stuckelberg field. We define the conjugate = momenta and the Poisson bracket algebra as, $$\frac{_{St}}{\dot{A}^i}\mathrm{\Pi }^i=3D\dot{=}A_i+_iA_0\frac{\mu }{2}ϵ_{ij}A_j;\frac{_{St}}{\dot{A}^0}\mathrm{\Pi }=^0=3Dm^2\theta ;\frac{_{St}}{\dot{\theta }=}\mathrm{\Pi }_\theta =3Dm^2\dot{\theta },$$ $$\{A_\mu (x),\mathrm{\Pi }_\nu (y)\}=3Dg_{\mu \nu }\delta (xy),\{\theta =(x),\mathrm{\Pi }_\theta (y)\}=3D\delta (xy).$$ (3) The Hamiltonian is $$_{St}=3D\mathrm{\Pi }^\mu \dot{A}^\mu +\mathrm{\Pi }_\theta \dot{\theta }=_{St}$$ $$=3D\frac{1}{2}\mathrm{\Pi }_i^2+\frac{1}{4}A_{ij}A_{ij}+(\frac{m^2}{2}+\frac{\mu =^2}{8})A_iA_i=20\frac{\mu }{2}ϵ_{ij}\mathrm{\Pi }_iA_j$$ $$+\frac{1}{2m^2}\mathrm{\Pi }_\theta ^2+\frac{m^2}{2}_i\theta _i\theta +m^2(=_iA_i)\theta A_0(_i\mathrm{\Pi }_i+\frac{\mu }{2}ϵ_{ij}=_iA_j+\frac{m^2}{2}A_0),$$ (4) where a total derivative term has been dropped. The two involuting first = class constraints, (in the Dirac sense of classification), are $$\chi _1\mathrm{\Pi }_0m^2\theta ,\chi _2_i\mathrm{\Pi }_i+\frac{\mu =}{2}ϵ_{ij}_iA_j+m^2A_0+\mathrm{\Pi }_\theta .$$ (5) The unitary gauge, $`\varphi _1\mathrm{\Pi }_\theta ;\varphi _2\theta `$, establishes gauge equivalence between the embedded model and the = original MCSP model. This ensures that in the gauge invariant sector, = results obtained in any convenient gauge will be true for the MCSP = theory. We invoke the rotationally symmetric Coulomb gauge $$\psi _1A_0;\psi _2_iA_i.$$ (6) The $`(\chi _i,\psi _j)`$ system of four constraints are now second = class, meaning that the constraint algebra metrix is invertible. The = Dirac brackets, defined in the conventional way are given below, $$\{A_i(x),\mathrm{\Pi }_j(y)\}^{}=3D(\delta _{ij}\frac{_i_j}{=^2})\delta (xy);\{\mathrm{\Pi }_i(x),\mathrm{\Pi }_j(y)\}^{}=3D\frac{\mu }{=2}ϵ_{ij}\delta (xy)$$ $$\{\mathrm{\Pi }_i(x),\theta (y)\}^{}=3D\frac{_i}{=^2}\delta (xy);\{\mathrm{\Pi }_i(x),\mathrm{\Pi }_0(y)\}^{}=3Dm^2\frac{_i}{=^2}\delta (xy).$$ (7) The remaining brackets are same as the Poisson brackets. The reduced = Hamiltonian in Coulomb gauge is $$_S=3D\frac{1}{2}\mathrm{\Pi }_i^2+\frac{1}{2}_iA_j=_iA_j+(\frac{m^2}{2}+\frac{\mu ^2}{8})A_iA_i=20\frac{\mu }{2}ϵ_{ij}\mathrm{\Pi }_iA_j+\frac{1}{2m^2}\mathrm{\Pi }_\theta ^2+\frac{m^2}{2}_i\theta _i\theta .$$ (8) Although somewhat tedious, it is straightforward to verify that the = following combinations, $`\varphi =3D((ϵ_{ij}_iA_j),(ϵ_{ij}_i\mathrm{\Pi }=_j),\mathrm{\Pi }_\theta ,\theta )`$ obey the higher derivative equation $$(\mathrm{}+M_1^2)(\mathrm{}+M_2^2)\varphi =3D0;M_1^2(M_2^2)=3D\frac{1}{=2}[2m^2+\mu ^2\pm \mu \sqrt{\mu ^2+4m^2}].$$ (9) The spectra agrees with . Note that for $`\mu ^2=3D0`$, the roots collapse to $`M_1^2=3DM_2^2=3Dm^2`$, which is = just the Maxwell-Proca model, whereas for $`m^2=3D0`$ the roots are $`M_1^2=3D\mu ^2,M_2^2=3D0`$, indicating the = presence of only the topologically massive mode, since the Stuckelberg = field $`\theta `$ is no longer present. Prior to fixing the $`\psi _2`$ gauge, the gauge invariant sector is = identified as, $$E_i=3D\mathrm{\Pi }_i+\frac{\mu }{2}ϵ_{ij}A_j;B=3Dϵ=_{ij}_iA_j;\mathrm{\Pi }_\theta ;A_i+_i\theta ,$$ (10) where $`E_i`$ and $`B`$ are the conventional electric and magnetic field. In the reduced space, the Hamiltonian and spatial translation generators = are gauge invariant, $$_{St}=3D\frac{1}{2}(E_i^2+B^2+\frac{\mathrm{\Pi }_\theta ^2}{=m^2}+m^2(A_i+_i\theta )^2),$$ $$𝒫_{St}^i=3Dϵ_{ij}E_jB\mathrm{\Pi }_\theta (A_i+_i\theta =).$$ (11) Defining the boost transformation as $`M^{i0}=3Dtd^2x=𝒫_{St}^i(x)+d^2xx^i_{St}(x)`$, the Dirac brackets with the = gauge invariant variables are easily computed. They will contain non = canonical pieces in order to be consistent with the constraints. = However, changing to a new set of variables by the following canonical = transformations,=20 $$Q_1(Q_2)=3D\frac{1}{\sqrt{2^2}}[ϵ_{ij}=_iA_j\pm \frac{1}{m}\mathrm{\Pi }_\theta ];P_1(P_2)=3D[\frac{1}{\sqrt{2^2}}ϵ_{ij}_i\mathrm{\Pi }=_j\frac{m}{2}\sqrt{2^2}\theta ],$$ (12) we can convert our system to a nearly decoupled one. Passing on to the quantum theory, the redefined variables satisfy the = canonical algebra, $$i\{P_i,Q_j\}=3D\delta _{ij}\delta (xy);\{Q_i,Q_j\}=3D\{P_i,P_j\}=3D0.$$ (13) The electric and magnetic fields and the translation generators are = rewritten as, $$B=3D\frac{\sqrt{2^2}}{2}(Q_1+Q_2);E_i=3D\frac{1}{\sqrt{2^2}}[ϵ_{ij}=_j(P_1+P_2)+(\mu +m)_iQ_1+(\mu m)_iQ_2],$$ (14) $$H_{St}=3Dd^2x[\frac{1}{2}(P_1^2+_iQ_1=_iQ_1+M_1^2Q_1^2)+\frac{1}{2}(P_2^2+_iQ_2=_iQ_2+M_2^2Q_2^2)+\frac{\mu ^2}{2}Q_1Q_2]$$ $$P_{St}^i=3Dd^2x[P_1^iQ_1+P_2^iQ_2]$$ (15) In order to drive home the peculiarities of MCSP theory, let us briefly consider the special cases, $`m^2=3D0`$ or $`\mu ^2=3D0`$. In the former limit, giving the MCS theory, as we noted before, $`\theta =`$ field is absent, which makes the $`(Q_1,P_1)`$ pair identical to the = $`(Q_2,P_2)`$ pair, leading to the following relations, with = $`i[p(x),q(y)]=3D\delta (xy)`$, $$B=3D\sqrt{^2}q,E_1=3D\frac{1}{\sqrt{=^2}}(ϵ_{ij}_jp+\mu _iq),$$ $$H=3Dd^2x\frac{1}{2}(p^2+_iq_iq+\mu =^2q^2),P^i=3Dd^2x(p^iq).$$ (16) This set of relations is identical to those in and hence the results obtained by DJT will follow trivially. The latter case, $`\mu ^2=3D0`$, refers to the Proca model, where = $`M_1^2=3DM_2^2=3Dm^2`$, and we get, $$B=3D\frac{\sqrt{2^2}}{2}(Q_1+Q_2);E_i=3D\frac{1}{\sqrt{2^2}}[ϵ_{ij}=_j(P_1+P_2)+m(_iQ_1_iQ_2)],$$ $$H=3Dd^2x[\frac{1}{2}(P_1^2+_iQ_1=_iQ_1+m_1^2Q_1^2)+\frac{1}{2}(P_2^2+_iQ_2=_iQ_2+m_2^2Q_2^2)],$$ $$P^i=3Dd^2x[P_1^iQ_1+P_2^iQ_2].$$ (17) Following the prescription of DJT given in , the boost = generator $`M^{i0}`$ should be reinforced by the additional terms, $$mϵ_{ij}d^2x(\frac{P_1_jQ_1}{^2}\frac{P_2_jQ_2}{^2}),$$ such that the electromagnetic fields transform correctly. This addition, = however, generates a zero momentum anomaly in the boost algebra, $$i[M^{i0},M^{j0}]=3Dϵ^{ij}(M\mathrm{\Delta }),\mathrm{\Delta }=3D\frac{m^3}{4\pi =}\{(Q_1)^2(Q_2)^2\}+\frac{m}{4\pi }\{(P_1)^2(=P_2)^2\},$$ (18) where $`M`$ is the rotation generator=20 $$M=3Dd^2x(P_1ϵ^{ij}x^i=_jQ_1+P_2ϵ^{ij}x^i_jQ_2)$$ . Making the mode expansions, $$Q_1(x)(Q_2(x))=3D\frac{d^2k}{2\pi \sqrt{2\omega =(k)}}[e^{ikx}a(k)(b(k))+e^{ikx}a^+(k)(b^+(k))],$$ (19) and effecting the phase redefinitions, $$ae^{i\frac{m}{m}\theta }a,b=e^{i\frac{m}{m}\theta }b,$$ (20) where $`\theta =3Dtan^1k_2/k_1`$, one recovers the full angular = momentum as $$M=3Dd^2k(a^+(k)\frac{1}{i}\frac{}{\theta }a(k)+b^+(k)\frac{1}{i}\frac{}{\theta }b(k))+\frac{m}{=m}d^2k(a^+(k)a(k)b^+(k)b(k)),$$ (21) where the second term is the spin. Now comes the intriguing part, i.e. what happens when both $`\mu `$ and $`m`$ are nonzero. First of all, for simplicity, let us neglect $`O(\mu ^2)`$ terms, which makes $`_{St}`$ a decoupled sum of ”1” and ”2” variables. But even then, similar extensions in $`M^{i0}`$, as done in the previous cases, will not have the desired effect since the parameters present in $`_{St}`$, $`M_1^2(M_2^2)_{\mu =^2=3D0}=3Dm^2\pm m\mu `$ are different from the parameters appearing in = the electric field, $`(\mu \pm m)^2_{\mu ^2=3D0}`$. Obviously, if we = keep the $`O(\mu ^2)`$ terms as well, the situation will worsen since = $`_{St}`$ is no longer decoupled. This constitutes the main result of this Letter. In an earlier work , it was argued that the anomaly in = appeared only because of the mapping of the system in terms = of a scalar variable. However, it has been demonstrated in = how to overcome this problem, leading to the correct spin value of the = excitation in the process. As we have shown, this scheme is untenable in = the MCSP model. To conclude, We have shown that in the Maxwell-Chern-Simons-Proca model, = where two mass scales, topological and non-topological or explicit, are = present simultaneously, the electromagnetic field transforms anomalously = under Poincare transformations. The conventional way of = redefining the phases of the creation and annihilation operators of the = basic fields to remove the anomaly is inadequate in the present case. A deeper understanding of this pathological behaviour is necessary. However, in applications of condensed matter physics, where Poincare or Lorentz invariance is generally not a big issue, these = models can still play an important role.
warning/0003/astro-ph0003219.html	ar5iv	text	# The Northern ROSAT All-Sky (NORAS) Galaxy Cluster Survey I: X-ray Properties of Clusters Detected as Extended X-ray Sources 1footnote 11footnote 1Results reported here are based on observations made with the Multiple Mirror Telescope, a joint facility of the Smithonian Institution and the University of Arizona ## 1 Introduction Galaxy clusters are important tracers of the large-scale structure of the matter distribution in the Universe. As the evolution of clusters is closely linked to the overall evolution of the cosmic large-scale structure, important tests of cosmological models can be performed with statistical data on the cluster population. The mass distribution and the spatial clustering of clusters are particularly interesting measures in such studies (e.g. Henry et al. 1992, Bahcall & Cen 1992, Eke et al. 1998, Thomas et al. 1998, Borgani et al. 1999). The construction of well defined cluster catalogues and the compilation of their properties is therefore an important task for observational cosmology. Galaxy clusters were first detected and are continued to be cataloged from optical observations of galaxy density enhancements in the sky (e.g. Abell 1958, Abell et al. 1989 (ACO), Zwicky et al. 1961 - 68, Shectman 1985, Dalton et al. 1992, Lumbsden 1992, Collins et al. 1995, Couch et al. 1991, Bower et al. 1994, Postman et al. 1996, Olsen et al. 1999, Scodeggio et al. 1999). X-rays have also successfully been used to detect galaxy clusters and to conduct clusters surveys (e.g. Picinotti et al. 1982, Kowalski 1984, Lahav et al. 1989, Gioia et al. 1984, 1990, Edge et al. 1990, Henry et al. 1992, Romer et al. 1994, Pierre et al. 1994, Ebeling et al. 1996, 1998, Castander et al. 1995, Rosati et al. 1995, 1998, Burns et al. 1996 Collins et al. 1997, Burke et al. 1997, Vihklinin et al. 1998, Scharf et al. 1997, Jones et al. 1998, Böhringer et al. 1998, De Grandi et al. 1999, Ledlow et al. 1999, Romer et al. 1999). The use of samples of clusters detected and characterized by their X-ray emission for cosmological studies has two major advantages over samples based on optical observations. First, the optical observations (without very extensive redshift measurements) provide only the projected galaxy distribution and not all galaxy density enhancements in the sky are bound, three-dimensional entities. In fact, in the course of the ESO Nearby Abell Cluster Survey (Katgert et al. 1996, Mazure et al. 1996) it was found that of the order of 10% of the rich clusters from the catalogue of Abell, Corwin, and Olowin (1989) in the nearby redshift range $`z0.1`$ were spurious clusters without obvious clustering peaks in redshift space. For the optical surveys the reliability has improved, however, with the advent of multi-color surveys and machine based matched-filter selection techniques (e.g. Postman et al. 1996, Olsen et al. 1999). Extended X-ray emission from the hot intra-cluster plasma of galaxy clusters is a more clear indication of the presence of a large gravitationally bound mass aggregate since otherwise the hot plasma would have been dispersed immediately. And secondly, the X-ray luminosity is a parameter much more tightly correlated with the mass of clusters than the usual richness parameter measured in the optical (e.g. Reiprich & Böhringer 1999). Thus X-ray emission gives evidence for the presence of galaxy clusters within a certain mass interval. (The correlation of the X-ray luminosity and cluster mass actually shows a dispersion of about a factor of 1.6 if one wishes to determine the mass for a given luminosity - Reiprich & Böhringer, in preparation). The one exception is the case where the X-ray emission is not clearly extended and where the cluster emission could be confused with the emission of an AGN within the cluster or with a possible foreground or background source. This confusion is a problem for a very small fraction of the cluster sources, but in general X-rays are a very useful indicator of a true cluster. In addition projection effects are minimized in X-ray surveys since the X-ray surface brightness is more centrally concentrated than the galaxy distribution. The ROSAT All-Sky Survey (RASS), the only large scale X-ray survey conducted with an X-ray telescope (Trümper 1993, Voges et al. 1999), provides an ideal data base to detect large numbers of clusters and to compile an all-sky cluster catalogue with homogeneously applied selection criteria. To exploit this unique data base we are conducting an optical follow-up identification program and redshift survey of RASS X-ray clusters in the northern hemisphere, the Northern ROSAT All-Sky (NORAS) Cluster Survey project. A complementary survey, the REFLEX (ROSAT-ESO Flux Limited X-ray) Cluster Survey, is conducted for the southern part of the RASS (Böhringer et al. 1998, Guzzo et al. 1999). The NORAS identification program was started in 1992 in a first step with a list of extended X-ray sources extracted from the general source list of the first RASS processing (RASS I; Voges et al. 1992). Apart from the selection for extent the following extraction criteria were used: northern declination, a minimum distance of 20 degrees to the galactic plane, and a minimum count rate of 0.06 cts s<sup>-1</sup> in the ROSAT broad band (0.1 to 2.4 keV). The criterion of X-ray source extent was chosen for the selection of promising cluster candidates, because early tests have shown that such a sample would be highly enriched (by about 70 - 80%) in galaxy clusters. Contrary to the cluster selection scheme used for the REFLEX Survey which is based on the correlation of X-ray sources with galaxy overdensities, the present sample selection is purely based on X-ray information. With this different bias the NORAS Survey has also the potential to find more distant and possibly “opitcally dark” clusters. The identification of these sources is now complete. In this paper we present a catalogue of the X-ray properties of the 378 cluster sources and 117 non-cluster sources of the primary candidate list. An accompanying paper by Huchra et al. (1999) provides a detailed catalogue of the optical identifications and redshift measurements of this sample, and scientific aspects of this survey are discussed in a paper by Giacconi et al. (1999). While this survey was ongoing some of the brighter sources of this sample as well as some X-ray emitting Abell and Zwicky clusters were spectroscopically observed for the ”BCS program” (Ebeling et al. 1998) by Allen et al. (1992) and Crawford et al. (1995, 1999). The region with the deepest exposure in the northern RASS, the north excliptic pole with exposure times ranging from 2000 to over 40000 sec, is also the subject of a dedicated survey which has identified all X-ray sources (c.f. Henry et al. 1995; Gioia et al. 1995; Bower et al. 1996; and Henry et al. 1997). Further studies on the X-ray properties of a sample of these extended sources in $`2\mathrm{deg}\times 2\mathrm{deg}`$ sky fields extracted from the RASS revealed, that in the first standard processing of the RASS the count rate and the extent of the cluster sources are severely underestimated (see also Ebeling et al. 1996, DeGrandi et al. 1997). Therefore a detailed reanalysis of the sources in the present sample was necessary. Here we also report the results of the detailed reanalysis of the sources using a new X-ray source characterization technique. In 1996, a second revised processing of the RASS (RASS II; Voges et al. 1999) with greatly improved attitude quality and with a fully merged photon data base became available. Our reanalysis is based on these new data. Since the incomplete assessment of the X-ray count rate and extent in RASS I not only leads to an underestimate of the X-ray fluxes for extended sources but also to an incompleteness of the sample extracted from the data base with certain limiting parameters, we have also used the new RASS data base to explore the incompleteness of the present cluster sample in terms of a flux-limited X-ray selected sample of galaxy clusters. In this study we selected a subregion covering the right ascension range from 9<sup>h</sup> to 14<sup>h</sup>. In a first step we use the Abell cluster catalogue to study the completeness provided by the RASS I extent criterion. We further study the prospects of finding more clusters with a more comprehensive extent criterion based on the new X-ray source analysis technique. This study also points the way to a more complete selection of galaxy clusters from the RASS X-ray sources. We now apply this algorithm in the ongoing NORAS cluster redshift survey. Since we will probably not be able tp rapidly complete the identifications of the newly found sources, we decided to publish the first part of the survey for which identifications are now complete and redshifts are nearly (all but 9) complete. The main emphasis here is not to publish a catalog of a complete, flux-limited sample, but to compile a cluster catalogue with reliable identifications based on a wealth of X-ray and optical data which are included in the identification process in a comprehensive way. The present sample contains many newly found objects, some of which are interesting targets for further astrophysical studies. The paper is organized as follows. In Section 2 we summarize the properties of the primary RASS I source list of extended sources, and in Section 3 we describe the techniques used to reanalyze the X-ray properties of the sample sources. In Section 4 we present the X-ray source catalogue with detailed X-ray properties of the 495 sources. Major X-ray properties of the sources which help in the identification of the objects are discussed in Section 5. The completeness of the sample is addressed in Sections 6 and 7 where we report the results of a rigorous search for X-ray emission from all ACO clusters and search for more extended X-ray sources with an improved analysis algorithm in a test region ranging in right ascensions from 9<sup>h</sup> to 14<sup>h</sup>. In Section 8 we compare the present results to the previous ROSAT Bright Cluster Survey by Ebeling et al. (1998). Section 9 provides a summary and conclusions. Throughout the paper we are using a Hubble parameter of $`H_0=50`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and $`h_{50}=H_0/(50`$ km s<sup>-1</sup> Mpc$`{}_{}{}^{1})`$ and further $`\mathrm{\Omega }_0=1`$ for the density parameter and $`\mathrm{\Lambda }_0=0`$ for the cosmological parameter. ## 2 The RASS I list of extended sources During the ROSAT mission the first All-Sky Survey was conducted with an X-ray telescope (Trümper 1992, 1993). The RASS was performed over a period of six months from August 1990 to January 1991 with two follow-up auxiliary survey missions carried out to fill the gaps in the survey in February and August 1991. The first processing of the survey (RASS I) provided a source list of 49441 sources (Voges et al. 1992, 1996). For this first analysis the survey data received were sorted into one of 90 2-degree wide strips (oriented in the direction of constant ecliptic longitude) while the satellite was still scanning the sky. As a consequence of this, strips are overlapping in regions outside the equator. Photons are exclusively sorted into only one of the overlapping regions. Therefore the exposure time across the strips is quite homogenous, but no advantage can be taken of the high total exposure in the ecliptic pole regions. The survey product resulting from this first processing of the RASS will be referred to as RASS I data base. (The situation is different in RASS II, the second processing, where the $`6.4\mathrm{deg}\times 6.4\mathrm{deg}`$ sky regions contain the full exposure data from the RASS). The RASS was conducted with the X-ray telescope (Aschenbach et al. 1988) and Position Sensitive Proportional Counter (PSPC; Pfeffermann et al. 1986) providing a high sensitivity and a very low internal background. A histogram of the exposure time distribution of the NORAS survey area in RASS I is shown in Fig. 1 (where it is also compared to the exposure distribution in RASS II). For this statistic we use the maximum exposure in any of the strips for those cases where a sky pixel is covered by several survey strips. The mean and median exposure times are 397 and 402 sec, respectively. The source detection procedure was based on detections using two alternative sliding window techniques plus a subsequent evaluation of the source detection significance and quality based on a maximum likelihood method (Voges et al. 1992, for aspects of the maximum likelihood method see also Cruddace et al. 1991). Only sources with a likelihood of detection larger than $`L=10`$ were accepted into the RASS I source list. (The likelihood value here and throughout the paper is defined as $`L=\mathrm{ln}P`$, where P is the probability for a spurious source detection – or a spurious extent detection in the case of the extent likelihood.) This threshold was chosen such that an estimated fraction of less than 1% spurious sources enter the source catalogue. The detections are based on the source counts in the broad ROSAT PSPC energy band covering the detector channels 11 – 240 which roughly corresponds to an energy range of 0.1 to 2.4 keV. Further qualities of the sources evaluated during the maximum likelihood assessment in three energy bands comprise the source count rate, an estimated source extent in excess of the broadening of the sources due to the telescope-detector point spread function (PSF), and two hardness ratios based on the counts measured in the soft (channel 11 – 40) and hard energy band (channel 52 – 201) or in the hard band 1 (channel 52 – 91) and hard band 2 (channel 92 – 201), respectively. The source extent which is of special importance here was determined within the maximum likelihood analysis by assuming for the source image shape a convolution of two two-dimensional Gaussian functions for the PSF and the source shape (Voges et al. 1999), respectively. The result of this analysis is then a value for the excess extent in terms of a $`\sigma `$-radius of the second Guassian. In this approximation the Gaussian wings are less extended than both the wings of the PSF and the wings of a King-type surface brightness model (e.g. Cavaliere & Fusco-Femiano 1976, Jones & Forman 1984). This is one, probably minor, reason for the effect that some of the X-ray flux in the outer X-ray halos is underestimated in the RASS standard analysis. First tests of the source quality parameters in RASS I and ready identifications with existing source catalogues showed that many clusters of galaxies featured significantly extended X-ray emission in the RASS. It was known from previous X-ray surveys (e.g. the Einstein Medium Sensitivity Survey, e.g. Gioia et al. 1994, Stocke et al. 1994) that slightly more than about 10% of the X-ray sources at the depth of the RASS should be galaxy clusters emitting in X-rays (see also Böhringer et al. 1991). Therefore it was clear that an X-ray sample highly enriched in galaxy clusters could be obtained by selecting those RASS sources featuring a significant source extent. For example, this is demonstrated by a comparison of the extent properties of the X-ray sources for ACO clusters, stars, and AGN taken from the ROSAT Bright Source catalogue and the correlation with optical catalogues from Voges et al. (1999) shown in Fig. 2 for RASS II data (a comparable figure is also shown in Ebeling et al. 1996 for RASS I results). Even though the present study is concerned with RASS I results we are showing a statistic for RASS II in Fig. 2 because there is no principle difference and there is a larger data base of correlations with catalogued objects available for RASS II. About half of the galaxy clusters occupy an almost exclusive parameter space characterized by an excess source extent larger than 25 arcsec with a reasonably high extent likelihood (with a value of $`L=7`$). Only about 2-3 % of the non-cluster sources are found in this region. Since galaxy clusters account for about $`1015\%`$ of all X-ray sources we can expect a contamination of the order of $`2030\%`$ by non-cluster sources if the sample is selected from the upper right quadrant. This is approximately what is found below. That this small value of 25 arcsec for the excess extent radius threshold shows a significant effect is somewhat surprising, since the mean half power radius of the survey point spread function is 70 arcsec, much larger than this threshold. This can be explained by the fact that in the RASS analysis likelihoods are calculated separately for each photon before they are summed and therefore each photon can be weighted by its own PSF according to the place in the detector where it was registered. In this way photons registered in the central part of the detector with a half power radius of the PSF of 15 - 20 arcsec give a high weight to the maximum likelihood analysis. This makes the RASS source analysis very sensitive to the recognition of small deviations from the expected shape of point sources. In the following analysis we will not make use of the information on the detector positions of individual photons. This is a disadvantage when compared to the standard RASS maximum likelihood analysis, but other advantages more than compensate for this. A first search for X-ray selected clusters was made with RASS I sources flagged as extended. The selection criteria were as follows: For the extent parameters a minimum threshold of 25 arcsec for the extent radius and a minimum value of 7 for the extent likelihood was chosen, as indicated by the dividing lines in Fig. 2. Further a lower count rate limit of 0.06 cts s<sup>-1</sup> in the ROSAT broad band was set and the sky area was restricted to the region $`\delta 0\mathrm{deg}`$ and $`\|b_{II}\|20\mathrm{deg}`$. The sources at the count rate threshold are thus typically characterized by about 25 source photons. This leads to a source fraction of 76% galaxy clusters among the sources selected. The advantage of this approach is that it yields a low fraction of contaminating sources, provides an effective way to detect galaxy clusters, and involves relatively simple selection criteria. The disadvantage is that the selection by source extent is much more difficult to quantify and to model than for example a purely flux-limited selection technique. In total 537 sources matching the selection criteria were extracted from the RASS I data base. 40 of these sources have been found to be detections of secondary maxima or fragments of clusters which are already in the list due to a detection at the main maximum. The largest fraction of these fragment sources is located in the very extended, diffuse emission region of the Virgo cluster (Böhringer et al. 1994). The fragment sources were removed from the list after a careful check that they are not associated with another distinct X-ray source in the line-of-sight. We have also excluded from the present catalogue the two detections in the Virgo cluster at the position of M87 and M86, because a useful flux measurement in the Virgo region requires a more detailed approach. Thus we report results for these parts of Virgo separately. In the following we will therefore discuss the analysis and identification of the remaining 495 sources. Early 1996 a new ROSAT Survey product, RASS II, became available at MPE. This version which is based on a greatly revised attitude solution for the pointing of the satellite during the survey and also uses a much more stringent quality threshold for the times with acceptable attitudes was used to create a new RASS II source list from which the RASS Bright Source Catalogue (RASS BSC) was created (Voges et al. 1999). In this survey product the data are sorted in 1378 sky fields with sufficient overlap ($`0.23`$ degrees) to guarantee an undiscriminating source assessment in the boundary regions. Each field now contains all the photons registered for this part of the sky during the entire survey. The resulting exposure distribution in the NORAS survey area is also shown in Fig. 1. As expected this distribution features a tail of high exposures up to about 40000 s. The reanalysis of the X-ray sources of the present sample makes almost exclusively use of the RASS II data base. ## 3 Reanalysis of the X-ray source properties To reanalyse the sources we apply a novel technique that is essentially based on measuring background-corrected source counts as a function of a growing circular aperture and checking for saturation to determine the observed source counts. The growth curve of the counts as a function of aperture radius is also used subsequently to analyse further source properties. We therefore term this method the growth curve analysis, GCA. We preferred to apply this method over techniques applied in earlier studies. The Steepness Ratio Technique used in De Grandi et al. (1997) has some similarity to the present analysis, but makes only restricted use of the available photon data in only extracting the source counts in aperture radii of 3 and 5 arcmin. A comparison shows that the uncertainties in the determined count rates are usually higher for that technique than for GCA. Our preference of the GCA method over Voronoi-Tesselation and Percolation (Ebeling et al. 1996, 1998) is due to the fact that the GCA technique is simple to reproduce in models and simulations, the resulting count rates are quoted for a known aperture radius for each source allowing a better assessment of the results in subsequent modeling, and the GCA technique provides a set of very essential diagnostic plots which make the interactive evaluation of the reliability of the GCA results easy and transparent. In addition the VTP technique needs two counteracting steps to correct for the unobserved flux. The first step relies on an extrapolation based on the assumption of spherically symmetric sources and leads to a significant overcorrection which is then compensated in a second step with an a posteriori recalibration based on a comparison with pointed data. The present method achieves a good agreement with pointed data in one relatively minor ab initio correction in a first step as shown below. The presentation of more details and tests of the GCA method is planned for a future publication (Böhringer et al. in preparation), while the essential features of this method are described in the following. The reanalysis of the X-ray properties was conducted for all 495 X-ray sources in the sample using RASS II data in fields of $`1.5\mathrm{deg}\times 1.5\mathrm{deg}`$ centered on each X-ray source. For 17 nearby clusters featuring a large extent the analysis is performed in larger fields of $`4\mathrm{deg}\times 4\mathrm{deg}`$ or $`8\mathrm{deg}\times 8\mathrm{deg}`$. For 7 sources, where the exposure in RASS II is less than 70 sec, data fields were extracted from RASS I which features a higher exposure for these cases. The reason for the reduced exposure in RASS II is the tight quality constraint which leads to the rejection of some RASS photon data in RASS II as compared to RASS I. The 7 fields of RASS I used here were carefully checked and did not show any peculiarities as e.g. double or otherwise distorted images of bright sources which would indicate a problem with the attitude control during the observation. The primary data set used for each field consists of a photon event file containing all data of the photons registered in the field area and the corresponding exposure map, providing the exposure time as a function of sky position with a resolution of 45 arcsec pixels. The exposure maps include a broad band correction for vignetting and the effect of the shadowing of the support structure of the PSPC window. (The possible difference between the broad band vignetting correction and the proper correction for the specific source spectrum introduces an error no larger than 2% and no further correction is applied). The photon event files provide information on the sky position and energy channel as well as the time and the detector position for each registered photon. In the following analysis only the first two parameters are used. The three energy bands defined for our data reduction procedure are identical to those used in the RASS analysis (Voges et al. 1999): broad band, soft band, and hard band. For the source count rate determination and the shape characterization we use exclusively the photon counts in the hard energy band. In this band the soft X-ray background is reduced to about one fourth, while – depending on the value of the interstellar column density – 60 to 100% of the cluster emission is detected. Therefore the analysis in this energy band provides the highest signal-to-noise ratio and the most reliable count rates. Another quite important aspect of the choice of this energy band is that it minimizes the contribution of contaminating sources to the count rate. Since the majority of all sources in the RASS are softer than the cluster X-ray sources, their contribution to the hard band is usually less significant than the contribution to the broad band counts. ### 3.1 Source position and count rate Prior to the evaluation of the source count rate the source center position and the sky background brightness is determined. The input field centers chosen are the X-ray positions provided by the maximum likelihood technique of RASS II. These positions are not optimal for extended sources and in particular extended sources are sometimes multiply detected in the RASS. Therefore a redetermination of the source position is performed based on a moment method which determines the two-dimensional “center of mass” of the photon distribution within an aperture of 3 arcmin around the input value for the center. This procedure is iterated with the newly found center position until the process converges to a stable center position. The small aperture of 3 arcmin gives a high weight to local maxima. We check all the center positions interactively and correct those cases in which this method has settled on a secondary maximum or where the local maximum is obviously offset from the large scale symmetry of the cluster. Those 5 sources are marked in the catalogue. We have also applied the moment method for the determination of the center position using larger apertures (5 and 7.5 arcmin). In 7 cases we preferred to quote these centers in the catalogue. Also these cases are marked. The background of the field is then determined from a ring area centered on the source with an inner and outer ring radius of 20 arcmin and 41.3 arcmin, respectively, as shown in Fig. 3. The inner ring size has been chosen such that it is outside the outer radius of the X-ray emission for the majority of the clusters. The outer radius is chosen to make almost full use of the $`1.5\times 1.5`$ degree fields. This large background area ensures that the number of photons used for the determination of the background surface brightness is large and introduces an almost negligible Poissonian error into the source flux determination. There are 17 nearby clusters in the sample which exceed the inner background ring radius in size. For these clusters larger fields have been extracted from RASS II and larger background rings have been used in the analysis. The ring is subdivided in 12 sectors. The photons in each of the three energy bands in all sectors are counted and the exposure time for each photon position is obtained from the exposure map. The count rate in each sector and the surface brightness are then calculated by averaging in count rate: $$C=\underset{i}{}\frac{1}{t_i},$$ (1) where $`C`$ is the count rate, $`t_i`$ is the exposure time at each photon position, and the summation is over all photon events in the sector. The surface brightness is obtained by division with the sector area. An uncertainty for the surface brightness in each sector is calculated from Poisson statistics. To avoid that discrete sources located in the background ring are included in the background estimate, the median of the sector count rates is determined and sectors featuring a larger than $`2.3\sigma `$ deviation from the median are discarded from the further calculations. Even though sources are only expected to cause large enhancements, we also exclude sectors that have count rates which are too low by more than $`2.3\sigma `$, since some of the positive deviations are due to fluctuations, and to avoid a negative bias in the background the negative fluctuations have to be discarded for reasons of symmetry. The clipping threshold of 2.3 $`\sigma `$ guarantees a successful removal of sources with typical count rates above about 0.04 s<sup>-1</sup> which would otherwise introduce a typical error $`1\%`$ in the background determination. The chosen clipping threshold leads in general to the clipping of not more than 1 - 3 sectors which preserves most of the background area for averaging resulting in a small photon statistical error in the background of typically about 5%. This also shows that the variations in the background on that scale are generally small and hardly larger than what is expected from Poisson statistics. The procedure of the background determination is illustrated in Fig. 3. For the data in Fig. 3a there is no interference of background sources and the background is smooth enough that all sectors were included in the background measurement, while for the data in Fig. 3b one of the sectors had to be discarded due to the presence of a significantly disturbing source. The two figures show the photon distribution of the hard band counts for each source field. The two sources, which have net source counts of 75.2 and 127.9, respectively, will be used in the following to illustrate the further analysis. (These two sources selected as the first sources in the list illustrating the features we like to show: RXCJ0004.9+1142 is the first source in the list with less than 100 photons which features a small but significant extent and RXCJ0020.6+2840 is the first source showing the search for the plateau in the count rate in the presence of steps in the plateau region, see below). The cumulative source count rate as a function of radius starting from the earlier determined central position is then found by integrating the source counts in concentric rings outwards while subtracting the background contribution. The integration is performed using a ring-width of 0.5 arcmin (a reasonable resolution for the given PSF of the instrument). The source count rate is determined for each ring by weighting each photon with the local exposure time according to eq.(1). The integration is performed for the three selected energy bands. The results are count rate profiles as shown in Figs. 4a and 4b. The uncertainty corridors in these count rate profiles resulting from Poisson photon statistics are indicated as dashed lines in the figures. They also include the Poisson error of the background determination. In most of these cumulative profiles the count rate levels off to a plateau value which gives the total observed source count rate. The total observed count rate is determined in the automated source characterization program in two alternative ways. In the first approach we determine the radius outside of which the source signal increases less than the $`1\sigma `$ uncertainty in the count rate. This radius, which we call the outer radius of significant X-ray emission, $`R_x`$, is indicated by the vertical dashed line in Figs. 4a and 4b. The count rate with its statistical uncertainty at this radius provides the value of the significantly detected count rate of the source. The radius $`R_x`$ is shown for the two sources as the inner circle in Figs. 3a and 3b . Alternatively the total source count rate is measured by getting an estimate of the plateau level. In the simplest case it is the average of the flat plateau outside $`R_x`$. In practice we determine the mean value of the plateau as well as the slope by means of a linear regression method for the profile part outside $`R_x`$. If the slope of the plateau is less than $`0.8`$% of the total count rate per arcmin radius, the plateau value is accepted. If the plateau is decreasing the count rate is determined from the mean of three bins around $`R_x`$. If the plateau is increasing, another effort is made to find the best flat part of the plateau by iteratively excluding the outermost and in a second step also some of the innermost bins. This procedure helps in excluding an outer rise of the count rate profile due to a neighboring source or by skipping a few bins if the count rate curve has not completely saturated to a plateau at $`R_x`$. Both effects can be seen in Fig. 4b, where the outer radius of the considered plateau region before a secondary rise of the profile is indicated. The source analysis is checked in each case on the basis of diagnostic plots as those shown in Figs. 3 - 6. About 35% of the plateaus can be characterized by the first step and about 83% by the further iterative trials. There is a residual fraction of about 17% of the sources for which no satisfactory plateau can be established automatically often as a result of contaminating nearby sources. These cases are analysed individually. For all sources we determine the radius, $`R_{out}`$, out to which the plateau count rate was measured by simply following the profile until the plateau value is reached. For the count rates and their uncertainties quoted in the present catalogue we have adopted the following approach. For the count rates we take the results from the fitted plateau values, which are in general higher by a few percent than the count rates determined at $`R_x`$. This is mainly due to an insignificant further rise of the cumulative count rate profile beyond $`R_x`$. This rise is in most cases much smaller than the statistical error in the count rate. The error in the count rate is then calculated from the root mean square of the shot noise error and the deviation of the plateau value from the value at $`R_x`$. Taking the root mean square would be justified for errors which are statistically independent and Gaussian distributed. Since the second error is a systematic deviation, this does not apply in this case. Nevertheless this is for the present case a practical approach which integrates the two errors by putting a larger weight on the larger one of the two uncertainties. The correction of the measured count rate of the source to a total count rate is discussed below. ### 3.2 Spectral hardness ratio and source extent For the determination of the spectral hardness ratio the count rate in the soft band is also determined for the same radius as for the hard band value of $`R_x`$. The hardness ratio, HR, used here as well as in the RASS data base is defined as $$HR=\frac{HS}{H+S},$$ (2) where $`H`$ is the hard band and $`S`$ the soft band source count rate. The expected values for the hardness ratio for cluster sources is roughly in the range $`01`$ as shown in Fig. 8. The source extent is addressed in two ways: quantifying the source size and testing the probability that the extent is real, respectively. In the first analysis a King profile with the surface brightness distribution, $$S_x(R)=S_0\left(1+\frac{R^2}{r_c^2}\right)^{1.5\beta +0.5},$$ (3) (where $`R`$ is the projected radial distance from the soure center and $`r_c`$ is the core radius of the X-ray surface brightness distribution) convolved with the averaged survey PSF (as calculated by G. Hasinger from the ROSAT XRT/PSPC PSF averaged over the detector area with a correction for the vignetting effect) and azimuthally integrated is fitted to the differential count rate profile. For the profile parameters a fixed value $`\beta =2/3`$ is taken according to the most typical value found in X-ray cluster observations (e.g. Jones & Forman 1984). In the $`\chi ^2`$ fit the core radius is varied in steps of 0.5 arcmin, and the normalization is a free fitting parameter. Due to the very low count statistics we have thus limited the fitting parameters per step to one, the normalization, by choosing the most common value for $`\beta `$. The results for the two example sources are shown in Figs. 5a and 5b, and the best fitting profiles are indicated by dotted lines in Figs. 4a and 4b. Note that, while the figures show fits to the cumulative profiles for a better diagnosis of the results, the actual calculations are conducted for the differential profiles to assure statistical independence of the count rates in rings as required by the $`\chi ^2`$ fitting method. Together with the best fitting value we also keep the minimum radius which is still consistent within the $`2\sigma `$ uncertainty limit. Comparing the $`\chi ^2`$ values to the 1, 2, and 3$`\sigma `$ limits shown in Fig. 5 we find that for the two examples the first source is only marginally extended ($`2\sigma `$ result), while the second source features a clear and large extent. The second, more sensitive method is used as a test for the probability that the source has an extent at all. For this we use a Kolmogorov-Smirnov test comparing the expected cumulative count rate profile of a point source including background with the given instrument PSF with the radially sorted, cumulative, and unbinned photon counts out to a radius of 6 arcmin. This radius is larger than the 90% power radius of the survey PSF and thus provides enough leverage to display the deviations of extended sources, but is small enough to minimize the influence of possible background errors. This test does not depend on the assumption that the uncertainties are Gaussian distributed. Since the pure photon counts also contain the background counts the previously determined background surface brightness has also to be added to the expected point source profile. Examples of the expected point source profile and measured curves are shown in Figs. 6a and 6b. The background contribution to the expected profile is also indicated and we can see that it is generally a minor contribution at these small source radii. Tests with known point sources have shown that the misclassification of point sources as extended sources is generally less than about 5% if we take an upper limit for the KS probability of 0.01 to classify a source as extended (see Böhringer 1999, in preparation). For the two examples we find probability values of $`0.005`$ and $`1.3\times 10^{14}`$, respectively, and therefore both sources will be classified as extended in the catalogue. In the following we will use these results in the form of the extent parameter defined as $`P_{ext}=\mathrm{log}_{10}`$(KS probability). ### 3.3 Deblending and analysis of very extended sources The visual inspection of the diagnostic plots of the GCA results for all the sources showed that in 19 cases the source analysis suffered from the blending of the cluster source with another nearby source probably not associated with the intracluster X-ray emission. These sources were scheduled for another analysis including a deblending technique. The correction by deblending was performed for all the sources where the contamination was clearly recognized as due to point sources. Tests show that a single contaminating source is usually easily recognized if its contribution to the total hard band count rate is larger than 5 - 10% and if the source is outside the central 3 arcmin radius. The deblending is not applied to irregular clusters or clusters with substructure where the non-symmetric emission region is most probably part of the diffuse intracluster X-ray emission (and not likely to be due to a point source). In this second analysis the source region is divided in two sets of twelve sectors for the radial region 3 to 8 and 8 to 15 arcmin. The variation of the surface brightness in the different ring sectors is analyzed in a similar way as done for the background ring sectors. Contaminating sources are best flagged by selecting those sectors with a more than $`3.5\sigma `$ deviation from the median. For weaker sources this detection threshold corresponds roughly to sources with fluxes larger than $`3\times 10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup> and in general guarantees the deblending of sources with contributions larger than 10%. The clipping technique is illustrated in Figs. 7a and 7b. In the further analysis the marked sectors are interpolated, that is, they are assigned a value for the surface brightness equal to the mean of the remaining sectors. The automated clipping works well for most sources but in about 30% of the cases the clipping either did not remove the contaminating source completely or removed in addition other parts of the cluster. For these cases we preferred to determine the count rate of the contaminating source directly (with the analysis centered directly on the contaminating source) with the same deblending algorithm and subtracted it from the blended source count rate. The 19 sources that required a deblending are marked in the catalogue. Since the outer radius of the region in which the source profiles are analyzed is fixed in the automated source analysis routine, and since the cut out regions per source are restricted to a size of $`1.5\mathrm{deg}\times 1.5\mathrm{deg}`$ , some nearby clusters are too extended to be covered completely by the analysis. For these clusters, 17 in total, larger files of photon data from RASS II were requested with fields covering $`4\mathrm{deg}\times 4\mathrm{deg}`$ or even $`8\mathrm{deg}\times 8\mathrm{deg}`$ around the source. To these data the same source analysis was applied as described above with an extended radial range. The clusters that required a reanalysis in a larger sky field are also marked in the catalogue and can also be recognized by their large values of $`R_x`$ and $`R_{out}`$. ### 3.4 Flux and luminosity determination To determine the flux and the luminosity of a source we first obtain the value of the interstellar hydrogen column density as measured at 21cm (Dickey & Lockman 1990, Stark et al. 1992) for the direction of the source by means of the EXSAS software package routine (Zimmermann et al. 1994). The flux is then determined in a first step by calculating the conversion factor from count rate to flux for a source with a thermal spectrum and a temperature of 5 keV (based on a modern version of the radiation code by Raymond & Smith 1977), a metal abundance of 0.3 of the solar value (Anders & Grevesse 1989), a redshift of zero, and an interstellar absorption according to the measured 21 cm value (Dickey & Lockman 1990) and the absorption tables of Morrison & McCammon (1983). We calculate the flux in the nominal ROSAT energy band: 0.1 - 2.4 keV. All fluxes and luminosities quoted in this paper and in the catalogue refer to this energy band. Fig. 8a shows the conversion factors as a function of absorbing column density for the three plasma temperatures 2.5, 5 and 8 keV. We note that the main dependence is on the column density. The variation with temperature makes a difference of less than $`7\%`$ in the temperature range (2 - 10 keV). Only for temperatures below 1.5 keV larger corrections occur, which applies for the smallest groups of galaxies in the sample. The dependence on the metal abundances is even less, about 1%. Prior to any further knowledge about the redshift and the nature of the cluster source the estimated flux value represents a good first approximation. Note, however, that for non-cluster sources a different spectral shape and thus a different conversion factor is expected. Once the redshift of the cluster source is known we can determine its luminosity. The luminosity is also calculated for the ROSAT energy band. The calculation is performed iteratively. In a first step we calculate a trial luminosity from the estimated flux value and the redshift and use its value to estimate a cluster plasma temperature using the X-ray luminosity-temperature relation of Markevitch (1998) $$T_x=2.34L_{44}^{1/2}\times h_{50},$$ (4) where $`T_x`$ is in keV, $`L_{44}`$ is the X-ray luminosity in units of $`10^{44}`$ erg s<sup>-1</sup> (in the 0.1 - 2.4 keV band). We make use of the relation which was derived by Markevitch without any correction for cooling flows in $`L_{44}`$ and $`T_x`$. This applies to our case since we are only dealing with integral count rates and average spectral properties. (Note that we have approximated the exponent of 1/2.02 found by Markevitch by 1/2). The temperature estimate allows the calculation of a new count rate-flux conversion factor for which we now also take the redshift of the source spectrum into account and calculate the source rest frame value for the X-ray luminosity (which involves the equivalent to the cosmic “K-correction”). The K-correction term for example increases up to about 6% out to a redshift of $`z=0.3`$ for clusters with a temperature of about 2 keV and up to about 15% for clusters with 10 keV. The iteration is repeated twice but found to actually converge to the final solution in the first step. The final rest frame X-ray luminosity is the value quoted in the catalogue. We have compared our results for the count rate to flux conversion based on EXSAS software and special programs using a similar data base and radiation code as EXSAS with results obtained from XSPEC, the online PIMMS software, and the conversion factors used in Ebeling et al. (1998). The differences are always less than $`3\%`$. Thus the use of different flux evaluation software does not constitute a significant source of potential differences between different RASS cluster surveys. ### 3.5 Estimates of the total flux The X-ray fluxes determined from the observations may still be biased low compared to the total flux coming from the cluster, since part of the flux in the faint outer regions is lost in the background. We can use the fact that we obtained the flux within a well defined angular aperture to make an estimate of the flux that may lie outside this aperture. Note that we use this approach here as a tentative estimate of the flux lost and we will therefore make no further effort in this paper to use the results for a correction, since the underlying assumption that all clusters have the same self-similar shape is not realized precisely enough to make a case by case correction useful without further tests and justifications. The aperture radius that corresponds to the generally used plateau value of the count rate is $`R_{out}`$. To obtain a rough estimate of the flux possibly missed outside the aperture we adopt the following generic cluster model characterized by a $`\beta `$-model surface brightness distribution as given in eq.(3) with $`\beta =2/3`$ to extrapolate the surface brightness profile outside $`R_{out}`$. For the core radius we are not using the results of the $`\chi ^2`$ fit, since they have too large uncertainties. We rather prefer to make a rough estimate of the cluster size from its X-ray luminosity. From the studies by Reiprich & Böhringer (1999) we find that the cluster mass is well correlated with the X-ray luminosity according to the relation $$L_xM_{grav}^{1.2}.$$ (5) We further assume that the self-similar relation of the core radius and mass of the form $`r_cM_{grav}^{1/3}L_x^{1/3.6}`$ holds (see e.g. Kaiser 1986). Taking a Coma-type cluster with $`L_x7\times 10^{44}`$ erg s<sup>-1</sup> (0.1 - 2.4 keV) and a core radius of 300 kpc to normalize the relation we find $$r_c=0.3\mathrm{Mpc}\left(\frac{L_x}{7\times 10^{44}\mathrm{ergs}^1}\right)^{1/3.6}.$$ (6) Rather than integrating the X-ray flux of the $`\beta `$-model to infinite radius, we stop the integration at 12 core radii which is about as large as the virial radius of a Coma-type model cluster. The difference between an integration to infinity (as for example performed in Ebeling et al. 1998 and De Grandi et al. 1999) to the cut-off radius at 12 $`r_c`$ is about 8%. This overestimate of the flux for integration to infinity is in general smaller than the individual uncertainties but not negligible if one is concerned with the global bias of the sample. Applying this model to our cluster sample we can calculate the fraction of the X-ray flux missed for each of the sample sources. Fig. 9 shows these missing fractions as a function of the source luminosity and of the detected number of source photons. The mean missing flux is about 8.3%. (The most discrepant point with a missing flux of about 50% and $`L_x10^{44}`$ erg s<sup>-1</sup> in Fig. 9a is for example a distant cluster observed at low flux ($`0.510^{12}`$ erg s<sup>-1</sup> cm<sup>-2</sup>) at low exposure which would be excluded in a proper flux limited sample). The missing flux fraction, $`fr`$, features only a very weak dependence on X-ray luminosity. The linear regression fit to the function $`fr=f(\mathrm{log}(L_x)`$ shown in Fig. 9a decreases from 8.3% for $`L_x=10^{43}`$ erg s<sup>-1</sup> to 7.9% for $`L_x=10^{45}`$ erg s<sup>-1</sup>. A more significant dependence is found for the number of source photons, as could be expected since this is the main parameter determining the significance of the source detection and how far out the count rate integration can be performed. Here the linear regression fit of Fig. 9b, $`fr=f(\mathrm{log}(N_{ph})`$ shows a decrease of $`fr`$ from 9.5% for 30 source photons to 6.7% for 500 source photons. Also this dependence is weak. An exception to the relatively small values for the missing fraction constitute some of the low luminosity sources as displayed in Fig. 9a. These are elliptical galaxies or very small groups dominated by elliptical galaxies, which have much smaller core radii - as indicated by the GCA King-model fits - compared to the values assumed after eq.(6). Thus for these small objects the model assumption seems to break down and the actual values for the missing flux is much smaller than estimated here. The validity of this generic model for the extrapolation of the total flux will be pursued in more detail in a future publication. A comparison with pointed observations performed below gives already an encouraging confirmation of these estimates. Alternatively we explore a second approach to estimate the missing X-ray flux by fixing the core radius in the $`\beta `$ model to 250 kpc. This approach was used in earlier studies of X-ray cluster samples (e.g. Henry et al. 1992). Figs. 10a and 10b show the results corresponding to the results of Fig. 9. One clearly notes a very steep increase in the missing flux for decreasing X-ray luminosity in Fig. 10a. The dependence on the photon number for which we would expect the strongest dependence is much less pronounced and not essentially different from the results in Fig. 9. One notes, however, that the scatter in Fig. 10b has approximately doubled. The obvious interpretation of these results is that the strong dependence on X-ray luminosity seen in Fig. 10a is artificial and results from an overestimate of the core radius for the less luminous objects. This inappropriate choice of the core radius also increases the scatter in Fig. 10b. Thus we conclude that this approach is clearly inappropriate for our study and the above used scaling of the core radius is a reasonable choice. ### 3.6 Comparison with the results of RASS I A comparison of the count rates measured in RASS I for the 378 cluster sources listed in Table 1 with the results of the GCA reanalysis is shown in Figs. 11a and 11b. For the comparison the RASS I count rates measured in the broad band have been converted to hard band counts by means of the measured hardness ratio. Sources which feature a significant extent according to our new analysis are marked in the plot. There is a large fraction of sources for which the count rates measured in RASS I or RASS II are underestimated by up to an oder of magnitude, which are essentially the sources marked as extended by the GCA method. The pointlike sources scatter around the line of equal count rate with an increasing scatter with decreasing count rate. The increase of the RASS I to GCA count rate ratio for low count rates seen in Fig. 11b is most probably an artefact produced by the previously set count rate limit in the selection of the RASS I sources for this sample (as indicated in the figure by the dotted line). The source of disagreement for the extended sources results from the design of the source analysis technique used for the RASS which is tuned to work optimally for point sources. Two effects discriminate against the proper accounting of the count rate of extended sources: i) the Gaussian kernel of the source shape fitting of the maximum likelihood analysis is bound to miss the outer wings of a typical cluster surface brightness distribution and ii) part of the outer X-ray halos of extended clusters may be treated as background by the background spline fitting process as used in the RASS standard analysis. A source analysis technique tuned to process the extended sources properly is therefore required to avoid these problems and to obtain correct X-ray parameters for the objects in our sample. The reanalysis of all the X-ray sources in the sample was therefore a necessary prerequisite for the compilation of an X-ray cluster catalogue to be used for astronomical and cosmological studies. ### 3.7 Comparison to pointed observations To test the count rate determination of our new analysis technique against a more reliable standard we have analyzed 80 clusters of our sample in pointed observations in which the better photon statistics allows a more detailed and precise analysis. The results for the brighter sources were taken from the compilation of Reiprich & Böhringer (1999) of the ROSAT clusters with the highest flux. 13 of the clusters from this sample were also analyzed in very large RASS survey fields ($`4\mathrm{deg}\times 4\mathrm{deg}`$ or $`8\mathrm{deg}\times 8\mathrm{deg}`$) because of the large cluster sizes. In the analysis by Reiprich & Böhringer (1999) a more refined source analysis was performed (where contaminating sources are excised in a wide region in and around the cluster and a possibly badly measured background is iteratively corrected by parabolic fits to the azimuthally integrated background surface brightness profile outside the cluster). Therefore it is also interesting to keep these objects in the list for comparison. The other data were retrieved from the ROSAT archive. The analysis technique used is similar to the one described above with a main difference that contaminating sources are excised interactively. We are using the count rate at $`R_x`$. The values for $`R_x`$ found in the pointed observations are generally larger than $`R_{out}`$ found in the RASS data as the flux can usually be traced further out into the background in the deeper observations. Fig. 12 shows a comparison of the count rates found in the two data sets. The objects analyzed in large RASS fields are marked with open symbols. The mean deviation of the count rates determined for the pointed data and the present results is 8.6%. Thus we conclude that the missing flux is on average about $`710\%`$ without a significant bias as a function of X-ray flux. This result is in excellent agreement with the estimates for the missing flux in section 3.5. The validity of the GCA approach is thus confirmed in two ways: the missing flux fraction is relatively small compared for example to the measurement errors and the ab initio estimates for the missing fraction are approximately correct. ## 4 The catalogue of RASS I extended sources: clusters and non-cluster sources In the following we are presenting the catalogue of the sample of 495 RASS I extended X-ray sources identified as galaxy clusters or as non-cluster X-ray sources. We discuss further characteristics of the source properties and the sample in the subsequent sections. We split the catalogue in three parts, the list of 378 sources identified as clusters, the list of 99 non-cluster sources, and a list of 17 X-ray AGN and one star located in clusters or in the line-of-sight of clusters, where the cluster is not the main source of the X-ray emission. The first tables, Table 1 - 4 list the major properties of the sources. Further X-ray parameters for the galaxy clusters are given in Table 5. The parameters of the table columns of Table 1 are described as follows. Column (1) lists the source name given by the following scheme: we use the prefix RXCJ for the reanalyzed RASS sources that have been identified with a galaxy cluster, where the ”C” stands for cluster. This prefix will be exclusively used for all RASS clusters analyzed by the above GCA technique. In particular this prefix has also been assigned to the RASS clusters identified within the REFLEX Survey (Böhringer et al. 1998, Guzzo et al. 1999). The remaining part of the name refers to the source coordinates for the epoch J2000 in hours (RA) or degrees (DEC), minutes, and fractions of a minute. Since the cluster sources are usually extended by more than one arcmin we use an arcmin precision in converting the cluster coordinates into the remaining part of the source name, thus the total name has 15 digits. Note that the coordinates in the cluster source name can deviate from the official RASS source catalogue coordinates reflecting the difference in the source analysis technique used. Column (2) gives alternative source names for previously catalogued optical counterparts to the X-ray sources, mainly Abell and Zwicky cluster names, names of NGC and UGC galaxies forming the central dominant galaxies of groups, and previously identified RASS or other X-ray sources. The Zwicky cluster names given in the table conform with the convention of the NED data base (note that the coordinate reference in the name refers here to epoch B1950). Columns (3) and (4) give the source position in decimal degrees for the epoch J2000. Column (5) gives the redshift of the cluster. Column (6) lists the measured count rate in units of counts s<sup>-1</sup>. Columns (7) and (8) give the X-ray flux in units of $`10^{12}`$ erg s<sup>-1</sup> cm<sup>-2</sup> for the flux estimated for a temperature of 5 keV in the first step and the finally calculated corrected flux, respectively. This correction includes the recalculation of the count rate flux conversion for the best temperature estimate and the redshift of the spectrum, but not the addition of the estimated missing flux. The fractional uncertainty in percent for the count rate, the fluxes and the luminosity are given in Column (9). Column (10) gives the rest frame X-ray luminosity of the clusters in the 0.1 - 2.4 keV energy band. Column (11) lists the value assumed for the absorbing column density in the line-of-sight to the source in units of $`10^{20}`$ cm<sup>-2</sup>, and column (12) indicates with the sign “?” if the identification as a cluster leaves some residual doubts. In some cases a second flag provides information on the way the source position was determined as explained below. Column (12) also lists the flags for the sources which have been deblended, flag $`B`$, analyzed in extra large fields, flag $`L`$, and sources which may be contaminated, flag $`C`$. The last column, (13), gives the reference number for the redshift as listed in the table caption. The objects marked by “?” or $`C`$ are commented below. For the position we have used the coordinates determined by the moment method with a 3 arcmin aperture radius. As an exception in the case of 12 cluster sources we have chosen to either redetermine the center position by a moment method with a larger aperture of 5 or 7.5 arcmin or we have determined a center position by hand. In these sources the automatic detection has selected a maximum which is significantly offset from the global center of symmetry of the clusters. These clusters are marked by the flags b, c, o in column (13) for the 5 arcmin and 7.5 arcmin moment method and the determination by eye, respectively. Table 2 lists the equivalent parameters for the RASS I extended sources identified as non-cluster objects. Here column (2) gives the source type of the identification except for cases with popular object names. A more detailed identification will be given in the second paper by Huchra et al. (1999). Columns (3) to (8) have the same definition as these columns in Table 1. Columns (9) to (14) give the extent parameter $`P_{ext}`$ from the KS test, the best fitting core radius, its minimal value consistent with the $`2\sigma `$ uncertainty limit, the hardness ratio, its error, and the deviation from the expected value of $`HR`$ in units of $`\sigma `$. The hardness ratio does occationally exceed the value of one in the tables, which occurs when the soft photon counts in the source region fall short of the background expectation. Note that for the parameters $`P_{ext}`$ and $`\mathrm{\Delta }HR`$ an upper limit to the numerical value of 30 and 10, respectively, was set, which was also used in plotting the data. Column (15) gives the interstellar HI column density. In the last column (16) objects with no certain optical identification are marked with “?”. The identification strategy for the non-cluster sources is explained in section 5. The last column also contains the flag for the alternative center determination. Table 3 lists X-ray sources where the X-ray emission is obviously originating from an AGN and in one case from a star, but where a cluster or group of galaxies is also visible on the Palomar Sky Survey images or on the CCD frames taken for this project. In several cases which are commented in Table 4 the redshift information indicates in addition the existence of a cluster and that the AGN is at the same redshift. The meaning of the columns of Table 3 is the same as that for Table 2. Table 4 provides the redshifts of the AGN as far as known and comments on the source identification. Figure 13 shows the distribution of the sources from Table 1 - 3 in the sky. Since we have not yet imposed a strict flux limit to the survey and also due to the incompleteness of the sample we cannot necessarily expect a homogeneous coverage of the sources in the sky. Near the coordinates $`RA=270\mathrm{deg}`$ and $`DEC=70\mathrm{deg}`$ we note a significant concentration of sources, which is due to the concentration of overlapping survey strips at the north ecliptic pole (NEP). (Note that this is not due to a pile up of exposure time at the NEP, since each survey strip was analysed independently, but due to the fact that the chance for detection of low flux, extended sources is increased in the multiple strips. A homogeneous coverage is achieved here only after imposing a proper flux cut.) There are also several low density regions in the source distribution as for example at the northern tip of the south galactic cap region. The distribution in X-ray luminosity and redshift of the sources in the cluster sample is shown in Fig. 14. The parabolic curves indicate limiting fluxes of $`10^{12}`$ erg s<sup>-1</sup> cm<sup>-2</sup> and $`3\times 10^{12}`$ erg s<sup>-1</sup> cm<sup>-2</sup>, respectively. There are 5 clusters with redshifts larger than 0.4 and 18 clusters with redshifts between 0.3 and 0.4. About half (11) of these clusters have luminosities in excess of $`10^{45}`$ erg s<sup>-1</sup> and belong to the most X-ray luminous clusters in the Universe. ### 4.1 Comments on individual objects RXCJ0106.8+0103 has an active galaxy in the cluster noted in the survey by Romer et al. (1994). The RASS image is consistent with a point source and a cluster spectrum. The available ROSAT HRI image has a small but significant extent. If the emission is mainly due to the cluster it would be very compact for the given high luminosity. A CCD exposure shows a nice cluster image and therefore we expect that this object is an X-ray cluster with a severe X-ray contamination by the AGN. RXCJ0255.8+0918 is a pointlike X-ray source with spectral properties consistent with intracluster medium emission centered on an elliptical galaxy. The galaxy is also classified as Sy2/LINER (Pietsch et al. 1999). They also note that the galaxy is located in a group. Without further information it is difficult to distingish between AGN or hot gas halo emission. The X-ray luminosity of $`L_x=0.7\times 10^{43}`$ erg s<sup>-1</sup> could well be that of a small group. RXCJ0311.5+0714 is not certainly confirmed as an X-ray cluster and the redshift is uncertain since only one galaxy redshift is available. RXCJ0728.9+2935 is a cluster candidate in a crowded stellar field featuring a pointlike X-ray source. Only one galaxy redshift is available for the X-ray source region. RXCJ0736+3925 shows extended X-ray emission, mostly due to intracluster medium emission, but obviously also some contribution by a softer central point source. RXCJ0921.1+4538 coincides with the radio galaxy 3C219. The X-ray source is point like but the hardness ratio is consistent with hot gas emission. 3C219 is located in a galaxy group. Without further information we cannot definitely decide if the X-ray emission comes from the AGN or the group. RXCJ1009.3+7110 is a case similar to RXCJ0106.8+0103 where a Seyfert 2 galaxy is listed in the Veron catalogue (Veron-Cetti & Veron 1998) and the ROSAT HRI image shows most probably a point source with a very compact halo. Thus this object is also classified as cluster contaminated by an X-ray AGN. RXCJ1122.2+6712 shows in the available HRI image a small extended halo around the galaxy VII Zw 392. The RASS source is contaminated by a nearby point source. The X-ray halo is with a luminosity of less than $`10^{43}`$ erg s<sup>-1</sup> quite faint corresponding to a very small group. It is one of the faintest sources in the sample. RXCJ1157.3+3336: For this X-ray source two cluster identifications are possible. It is associated with Abell 1423 at a redshift of 0.0761 but there is also the possibility that a cluster is associated with the radio galaxy 7C 1154+3353 located at the center of the X-ray emission. The 7C galaxy is classified as cD galaxy. We adopt the identification of an X-ray cluster associated to the 7C galaxy because of the better positional coincidence. Crawford et al. (1999) note the same identification. RXCJ1229.7+0759 and RXCJ1243.6+1133 are the X-ray halos of the two Virgo cluster galaxies M49 and M60, respectively. They are included in this catalogue, even though these halos are embedded within the low surface brightness structure of the Virgo cluster emission. But since the two local halos stand out from the low surface brightness emission environment and can be reasonably characterized by the present source characterization technique we have not excluded them from our catalogue as we have excluded M86 and M87. RXCJ1510.1+3330, RXCJ1556.1+6621, RXCJ1700.7+6412 show a contamination by a point source in the available HRI image by no more than 10 - 15%. RXCJ1518.7+0613 is contaminated by X-ray emission from an AGN in the cluster, which is also indicated by a softer hardness ratio than expected for a galaxy cluster. RXCJ1554.2+3237 is not yet definitely confirmed as cluster and the redshift is uncertain since there is only one galaxy redshift. RXCJ1644.9+0140 and RXCJ1647.4+0441 are cluster candidates in a crowded stellar field where no conclusive redshift has been ontained so far. RXCJ1644.9+014 is possibly a distant cluster. RXCJ1738.1+6006 is also a cluster candidate in a crowded stellar field with extended X-ray emission at the significance threshold. No redshift has yet been obtained. A Seyfert galaxy is known with a distance of 1 arcmin from the X-ray maximum which is most probably not the X-ray source because the offset would be unusually large. RXCJ1800.5+6913 has a complex structure and is obviously contaminated by emission from point sources. But these sources contribute less than 20% to the overall emission from the cluster. RXCJ1832.5+6848 contains a BL Lac in the cluster center, which most probably severely contaminates the cluster X-ray emission. RXCJ1854.1+6858 is a cluster candidate featuring an extended X-ray source. The redshift is uncertain since only one galaxy redshift is available. RXCJ2035.7+0046 is an extended low surface brightness source with a detection significance of 3 to 4 $`\sigma `$. It is a cluster candidate in a very crowded stellar field for which no redshift is available yet. RXCJ2041.7+0721 is a cluster candidate in a crowded stellar field with only one available redshift. ## 5 Discussion of the major source properties Further properties of the X-ray cluster sources are provided by Table 5. The columns of the table are as follows. Column (1) and (2) repeat the name and rest frame ROSAT band X-ray luminosity of the sources from the previous tables. Columns (3) and (4) give the radius out to which the source count rate has been integrated, $`R_{out}`$, in units of arcmin and Mpc, respectively. Column (5) gives the probability result of the Kolmogorov-Smirnov test for the source to be a point source in terms of the parameter $`P_{ext}`$. For very low probabilities for the consitency with a point source the parameter $`P_{ext}`$ was limited to a value of 30 for the entry in the table. A source is considered very likely to be extended if this parameter has at least a value of 2 (point source excluded with 99% probability). Column (6) and (7) give the best fitting core radius for the King model fit and the minimal core radius still consistent within 2$`\sigma `$ error limits, respectively. Note that the core radii determined here are a only qualitative measure for the source extent, since in general the errors are very large and the fitting grid was coarsely spaced. Therefore we do not recommend to use the results for $`r_c`$ as a measure of the statistics of the cluster shapes at this point. Column (8) and (9) give the spectral hardness ratio defined by eq.(2) and its Poissonian error. Column (10) finally indicates the deviation of the measured hardness ratio from the expectation value calculated for the given $`N_H`$ and for a temperature of 5 keV. This deviation parameter is given in units of the $`1\sigma `$-error of the hardness ratio and the listed and plotted values were limited to a maximum numerical value of 10. One of the most interesting first questions about these parameters concerns the discrimination power of the spectral and extent parameters in distinguishing between cluster and non-cluster sources. This is analysed in detail for statistical samples of sources with known identifications in the forthcoming paper by Böhringer et al. (in preparation). The present source sample should be highly biased since it was preselected from the RASS I data with the criterion of showing a significant extent in the RASS I analysis. We have to expect for example that the subsample of point sources within the present sample is already enriched in pathological cases including e.g. double and very bright sources. Therefore the present sample is used here only for a qualitative discussion while statistical numbers can only be obtained from the analysis to be published in the following paper. Fig. 15 displays the distribution of the spectral and spatial extent parameters for the cluster and non-cluster sources for comparison. For the spectral discrimination we plot the difference of the actually observed hardness ratio to the theoretically expected hardness ratio as calculated for a 5 keV cluster as shown in Fig. 8. The difference is thereby quantified in terms of the $`\sigma `$-deviation accounting for the uncertainty of the hardness ratio measurement. To quantify the extent we use the probability result of the KS-test for a point source in terms of the extent parameter $`P_{ext}`$. As we can see most of the clusters occupy the upper part of the figure and cover a wide range of extent values. As expected, non-cluster sources are concentrated in the region of small values of $`P_{ext}`$. Clusters have harder spectra than the average of the other X-ray sources being mostly stars or AGN, and therefore clusters should be separable by the hardness ratio from the softer part of the non-cluster population. On the other hand the extent parameter should help in the discrimination against the non-cluster sources since they are to the vast majority point-like sources at the angular resolution of ROSAT. This is clearly seen in Figs. 16 and 17 which show a blow-up of Fig. 15 and display the cluster and non-cluster sources separately. Only 12 clusters out-off 378 show a more than $`3\sigma `$ deviation to the X-ray soft side. These are to about one half very bright clusters (including Coma) where the relative deviation is small in absolute sense and could reflect for example an incorrect value for the column density. For some of the sources we suspect some contamination by AGN to the overall emission. One of these cases is for example A1722 for which the HRI observation confirms the contribution of a point source to the extended cluster emission by an amount of about 30 - 40%. Anyway these pathological cases comprise only about 3% of the sample and the hardness ratio provides indeed a very powerful diagnostics for the source identification. A similarly interesting result is found for the extent parameter, $`P_{ext}`$ for the case of the non-cluster sources shown in Fig. 17. Only 22 out-off 117 sources show an extent if we adopt the extent threshold criterion defined above, that is, a value of $`P_{ext}2`$. Some of these sources are very bright and soft and are far from the expectation for a cluster type spectrum. If we exclude these sources with deviations larger than $`8\sigma `$, there are only 13 sources for which the nature of the extent should be investigated to rule out a cluster nature. For more than half of these sources we find a reason why they are featuring an extent: 2 are nearby galaxies, M82, M106, which show extended X-ray emission, 7 of the sources are double sources where the main, catalogued source is consistent with a point source. Some of the remaining sources are very bright in the RASS (more than 500 photons) and, therefore, the relative deviation corresponds to a very small absolute deviation which can be caused by small systematic effects. A pathological object in this subsample is the BL Lac RXJ1456.0+5048 which has a value for $`P_{ext}`$ of 5.4. It has been analyzed in a RASS-follow-up HRI observation by Nass (1998) and shows at most a very marginal extent in this observation with a ten times higher angular resolution. Therefore the extent seen in the present RASS analysis has to be spurious. We have also checked the mean off-axis position of the source photons as a signature of an imperfect scanning of the source which could produce a deviation of the mean PSF for the source observation compared to the average survey PSF, but found no obvious deviation. Therefore this case demonstrates that systematic effects in the broadening of the PSF for individual sources in the RASS exist, but they are obviously extremely rare. Subtracting the cases for which the extent is real, we find that less than 10% of the sources feature a spurious extent. Applying a similar test to an unbiased point source sample shows that the misclassification is usually less than about 5% (see Böhringer 1999 in preparation). In Fig. 17 we have also marked the non-cluster X-ray sources which seem to be associated with an optical cluster either in projection or located in the cluster. There is no indication that these sources show a different distribution in the source quality parameters. Therefore the identification that in most of these sources the AGN or star is clearly the dominant X-ray source is supported by this result. These results on the source quality can be compared with the prime selection criteria of this sample as extracted from RASS I. We recall that only sources which were characterized by a significantly large extent parameter obtained in the maximum likelihood analysis were included in the present sample. Using the new technique $`25\%`$ of the sources do not feature an extent. This could be partly due to the fact that the present method is sometimes less sensitive in recognizing the source extent for very compact sources since it does not use the detector position information and thus does not weight for photons imaged with different sharpness as does the maximum likelihood method. Looking at the nature of the sources we conclude that about 23% of the sources are genuine point sources and thus the failure rate of the RASS I extent classification is at least about 20%. The results displayed in Figs. 15 - 17 show that the present method provides a great improvement concerning the failure rate of the method. This justifies the use of these source quality parameters to assist the identification of the sources. We have made use of the above results in the identification process. The criterion for classifying a source as a non-cluster source was one of the following: i) The optical counterpart is a bright star ($`12.5`$ mag) and the source is not extended. This is justified by the finding of Voges et al. 1999 (in preparation) that there is a much less than 1% chance coincidence of a RASS X-ray source with such a bright star. ii) There is an optically identified AGN at the source position and the source is not extended. No cluster is seen on the CCD images taken for this project. iii) The source is a known, previously identified non-cluster X-ray source. iv) The source is clearly point-like and has a soft hardness ratio deviation larger than $`3\sigma `$. There is no signature of an optical cluster in the digitized Palomar Sky Survey image or the CCD images taken for the project. We noted 17 cases, where there is a signature of an optical cluster, but the other criteria are inconsistent with a cluster identification. These cases have been listed separately in Tables 3 and 4. This identification scheme is only made possible by the extensive spectroscopic follow-up and almost complete CCD imaging of the targets in the sample of extended sources (no spectroscopy was done for those objects for which this information is already available from the literature or the CfA archives). A classification only by identification of clusters on optical images or exclusion of sources with spectroscopically identified AGN without further inspection would have failed in several cases. ## 6 Further analysis of the 9 - 14<sup>h</sup> study region While the above results show that the source identification has reached a high level of reliability, there are serious concerns about the completeness of the present sample in terms of a purely flux-limited X-ray cluster sample, because the sample rests on the RASS I source extent selection criterion. There are two principle sources of incompleteness: i) medium distant and distant clusters may have too compact X-ray emission regions to be resolved as extended X-ray objects in the RASS and ii) the RASS maximum likelihood analysis may fail to recognize all the spatially resolved sources in the RASS as extended X-ray sources. Both sources of incompleteness affect the present sample. Figs. 2 and 16 show that there are galaxy clusters which appear as point-like X-ray sources for both analysis techniques, the standard RASS analysis and the method used here. In addition we will find below that there is a large fraction of significantly extended sources which are missed by RASS I. In an attempt to check for the incompleteness of the RASS I sample we conduct two studies: we investigate what fraction of ACO clusters with X-ray emission is missing in the present sample and we search for extended sources in the RASS II data base with the new analysis technique and inspect their nature to find additional clusters or cluster candidates. Since this is meant as a statistical test we restrict the analysis to a subregion of the study area: the region of the northern sky between $`9^h`$ and $`14^h`$. The sky area of this region is 1.309 ster as compared to the total NORAS survey area of 4.134. In this right ascension range the $`\|b_{II}\|20\mathrm{deg}`$ band of the Milky Way is completely located at negative declinations. We restrict the cluster search by imposing an X-ray flux limit of $`1.6\times 10^{12}`$ erg s<sup>-1</sup> cm<sup>-2</sup>. Thus in the first step we have run our analysis on the sky positions in RASS II of the 901 clusters listed by Abell, Corwin, & Olowin (1989) in the sub-survey region. Ten clusters of this sample were missed in the analysis since RASS II has a too low exposure in the corresponding sky fields for a significant source detection (these clusters are: A917, A974, A986, A996, A999, A1011, A1042, A1057, A1128, A1554). 93 additional ACO clusters are detected above the flux limit. A closer inspection of the detections with the same procedures as described above revealed that 8 of the detections have most probably an AGN as the dominant X-ray source. These ACO clusters are: A763, A924, A1030, A1225, A1575, A1593, A1739, A1774. The results for the other 85 detections, where we identify the X-ray emission to originate in the cluster, are listed in Tables 6 and 7. The columns of these tables have the same meaning as those of Tables 1 and 5, respectively. In the second step we have selected all the sources from the RASS II data base in the sub-survey region, reanalyzed them, and extracted all the sources with a flux $`1.6\times 10^{12}`$ erg s<sup>-1</sup> cm<sup>-2</sup> and an extent parameter $`P_{ext}2`$. In total there are 377 sources of which a large fraction is already contained in the NORAS I sample and the supplementary ACO sample. Again we find another fraction which is identified as double point sources or very bright point sources having a small absolute extent. An inspection of all the remaining sources on the digitized sky images and literature data from data bases leads to a sample of 52 very promising candidates of which 21 can readily be identified as previously listed galaxy clusters. The latter list of identified cluster sources is given in Table 8 and 9. For the positive identification of the remaining candidates further optical observations within the ongoing NORAS survey are in progress. To understand which of the two sources of incompleteness discussed above is more important for the loss of these additional cluster sources in the primary candidate list of the NORAS Survey, we plot in Fig. 18 the distribution of the extent parameters of the X-ray clusters of the two additional cluster lists. As we can see most of the additionally found ACO clusters (74.2%) feature an extent according to the new analysis technique. The second supplementary list consists by definition only of extended cluster sources. Therefore we have to conclude, that the extent flag in the RASS I data base is not only not very reliable but a large fraction of the well extended sources is also missed. ## 7 Discussion of the sample completeness To test the completeness of the catalogue, a comparison can be made to the southern RASS cluster survey project, the REFLEX Survey (Böhringer et al. 1998). This sample has been constructed in a different way making extensive use of the COSMOS data base to correlate X-ray sources with the galaxy distribution. Therefore this sample does not rely on existing cluster catalogues nor on the selection of X-ray sources featuring an extent. Internal statistical estimates for the REFLEX sample suggest a completeness larger than 90% for the flux limit of $`F_X=3\times 10^{12}`$ erg s<sup>-1</sup> cm<sup>-2</sup>. The number counts of the REFLEX survey and the NORAS survey are compared in Fig. 19. We note that the NORAS sample reaches a fraction of 50% of the REFLEX number counts at the REFLEX flux limit. A large fraction of the missing clusters has been recovered by the supplementary sample in the study area. In Fig. 20 various subsamples of the combined cluster catalog in the 9<sup>h</sup> to 14<sup>h</sup> region are compared to the REFLEX Survey. While the cluster sample from Table 1 recovers with 40% an even smaller fraction of the sky density compared to the REFLEX Survey, 70% are reached if the ACO clusters are added and 82% are obtained for the combined sample. Thus there is still a fraction of clusters missing. This is easily understood, since the REFLEX sample contains a fraction of 22% of all objects which are not resolved as extended by the GCA analysis. We can expect that about 10 - 20% of the clusters for the REFLEX flux limit are neither listed as ACO clusters nor recognized as extended sources. Note that the combined sample in the study region completely recovers the previous RASS cluster sample by Ebeling et al. (1998) as discussed below. Most clusters not contained in Table 1 are easily found as ACO clusters and 7 further clusters show a clear extent. Therefore to achieve a higher completeness in our continuing northern cluster survey we are, in addition to including known catalogued clusters with X-ray emission and newly classified extended X-ray sources, conducting further imaging of promising non-identified X-ray sources to recover the compact X-ray clusters missed in the previous searches. ## 8 Comparison with the northern Brightest Cluster Sample In total 166 X-ray sources analyzed here overlap with the previous RASS galaxy cluster compilation by Ebeling et al. (1998). 142 sources coincide with the sources of Table 1, 22 sources with sources in Tables 6 and 8 for the 9<sup>h</sup> to 14<sup>h</sup> region, and two sources are identified with AGN in our analysis. Thus 40 sources compiled by Ebeling et al. are not included here (including the Virgo cluster). Since the two compilations are made with the same intention we compare the results of the two samples in some detail. Fig. 21 shows a comparison of the X-ray fluxes determined by the two different methods used in the two surveys for the 166 cluster in common. Note that in our compilation the fluxes have not been corrected for the missing flux outside the measurement aperture. Therefore we show both results from the work of Ebeling et al. in the plots: the uncorrected measured fluxes (with open circles) and the corrected final fluxes (full circles). The agreement at high fluxes is very good and there is also quite good agreement for most sources. The scatter is increasing, however, towards lower fluxes. Mostly at lower fluxes we find a few sources with fluxes larger by factors up to 2 as determined by the VTP method compared to the present results. Note the bias introduced at low fluxes by the flux limit set in the Ebeling et al. sample that suppresses the cases where the VTP to GCA flux ratio is lower than one. Thus, what appears like a high bias of the VTP results at low flux in Fig. 21b is most probably explained by an increased scatter in the flux ratio with decreasing flux. Thus in general the agreement is good. Two sources in the cluster list by Ebeling et al. (1998), A2318 and Z2701, are classified as AGN in our work. The first source appears as RXJ1905.7+7804 in Table 3. We have optically identified a Seyfert 1 galaxy at the source position. The X-ray source is pointlike and has a spectral hardness ratio that is too soft for a cluster with a 3.8$`\sigma `$ deviation. The X-ray source is also significantly offset from the optical cluster center. Thus we identify the X-ray source with the AGN. The second source was found in our search for additional extended sources in the 9<sup>h</sup> \- 14<sup>h</sup> region and is not catalogued here. It was dismissed because it is identified by as AGN by Bade et al. 1998. The spectral hardness ratio shows again a 4.5$`\sigma `$ deviation to the soft side. The KS-test yields a small extent of the source and an inspection of the X-ray image indicates a smaller emission contribution from the region of the optical cluster center. But the main contribution comes obviously from the AGN. We can also make a comparison of the completeness of both surveys in the 9<sup>h</sup> \- 14<sup>h</sup> region. While we recover all the clusters in Ebeling et al. (except for Z2701 which is classified as AGN), two additional clusters are found in our search with a flux larger than $`L_X4.7\times 10^{12}`$ erg s<sup>-1</sup> cm<sup>-2</sup> (in addition the two Virgo galaxies M60 and M49 which are summed into the Virgo cluster in Ebeling et al.). Over the whole right ascension range we find 12 additional objects above the flux limit of Ebeling et al. These sources are RXCJ0005.3+1612 (A2703), RXCJ0209.5+1946 (A311), RXCJ0736.4+3925 (contains a BL Lac), RXCJ1121.7+0249 (SHK352), RXCJ1242.8+0241 (NGC4636), RXCJ1447.4+0827, RXCJ1501.2+0141, RXCJ1506.4+0136, RXCJ1617.5+3458 (NGC6107), RXCJ1718.1+7801 (A2271), RXCJ1742.8+3900, RXCJ1900.4+6958 (A2315). All the sources except for RXCJ1447.4+082 have a significant extent. The Abell clusters have fluxes within about 15% of the flux limit of Ebeling et al. (1998) and most of the remaining sources are nearby groups of galaxies being less rich than Abell clusters. This number of extra sources is still compatible with the claimed completeness of the Ebeling et al. (1998) sample. The above comparison with the results of the southern RASS analysis in the REFLEX sample indicates, however, that both samples are still less complete than the southern survey. ## 9 Discussion and Summary Using the early results of the RASS we compiled an X-ray sample of galaxy clusters by selecting the RASS sources which featured a significant extent in the first standard processing of the RASS. The complete spectroscopic identification of these sources (as far as no identification was already available) leads to a compilation of a catalogue of 378 X-ray cluster sources. A reanalysis of the X-ray properties of these sources shows that the X-ray source properties can successfully be used in the source identification process. In particular we found that in many cases the hardness ratio can be used to flag sources which are severely contaminated by an AGN. In this way we could also exclude previously identified X-ray clusters – as for example the two cases discussed in the previous section – from identification as cluster sources. One of the major points of concern of such catalogs of X-ray clusters is the contamination of the observed X-ray luminosity by AGN. As commented in Section 4.1 there are all combinations of clusters and optically identified AGN or radio galaxies: AGN which do not provide a significant contribution to the cluster X-ray flux, AGN that contribute partially, and AGN which completely outshine the clusters. These cases can only definitely be distinguished if high resolution X-ray images or X-ray spectra of high quality are available. E.g. ROSAT HRI observations provide the means for this distinction, but they are only available for a small fraction of the sources for survey sizes of the present survey. Lack of additional X-ray data creates a twofold danger. Since X-ray sources are often identified with the most plausible nearby optical counterpart, there is a significant risk that the cluster ID will be discarded when an AGN is found near or within the error circle. Conversely, an AGN which is the primary X-ray emitter may not be recognized and the emission falsely associated to the cluster. Also in the present case this problem of AGN cluster associations cannot be solved in each case and we still expect some hidden misclassifications which will be hard to find until the whole sample is probed more deeply in X-rays. However, we believe that the remaining uncertainties are small and not harmful for the application of our sample for cosmological statistical studies and we will substantiate this believe in the following. First of all we note that $`74.4\%`$ of the sources in the catalog in Table 1 are sources recognized as significantly extended. This requires a major contribution to the X-ray emission from diffuse cluster sources. Thus for this source population we have a very high reliability that these sources are true X-ray clusters. We can now - as shown in Fig. 22 - compare the spectral properties of these extended cluster sources to those cluster sources where no significant extent could be established in terms of the parameter $`P_{ext}`$. We note that the two distributions are hardly different. Also both distributions are almost symmetric around unity, as would be ideally expected. For some reason the larger sample of extended sources shows a broader distribution of the $`\mathrm{\Delta }HR`$ parameter. This seems to be caused by clusters with a larger number of photons for which a small systematic error in the $`HR`$ estimate leads to a significant additional scatter in $`\mathrm{\Delta }HR`$ in units of $`\sigma `$ and also the low temperature groups will add to this scatter. A comparison to the non-cluster sources shows a very large difference in the spectral parameter distribution, however. While the non-cluster sources have a median deviation of the parameter $`\mathrm{\Delta }HR`$ of $`4.6\sigma `$ the X-ray pointlike sources identified as clusters have a median $`\mathrm{\Delta }HR0.02\sigma `$. We can explore the contamination effect in somewhat more detail by considering the effect on a typical cluster source. The median number of source photons for the clusters listed in Table 1 is about 94 source photons. For such a source the typical uncertainty in the measured hardness ratio is about $`\delta HR0.12`$. An AGN with median properties easily shows a deviation of $`34\sigma `$ in the $`\mathrm{\Delta }HR`$ parameter. Even if the contamination fraction by AGN is only 50% or 25% the typical shift in the parameter of $`\mathrm{\Delta }HR1.7\sigma `$ and $`0.9\sigma `$ still constitutes a recognizable distortion of the $`\mathrm{\Delta }HR`$-distribution. Another test is shown in Fig. 22b. Here we compare the distribution of the $`\mathrm{\Delta }HR`$ values for the clusters with the sample of point like clusters artificially contaminated by removing 20% of the cluster sources and replacing them by a statistical sample of non-cluster sources. There is a clearly visible change in the distribution of the spectral parameters. A KS test shows that the two distributions are still only distinguishable at the $`90\%`$ level. A test with a similar contamination of 40%, however, leads to a clear difference with a KS probability for the two samples being statistically the same of $`10^4`$. Since a 20% contamination of the point like clusters corresponds to a contamination of about 5% of the total sample we expect that the real misclassification fraction is not larger than this percentage. The present X-ray cluster catalog provides a wealth of new data. 98 new cluster sources are listed in Table 1 and in addition new X-ray luminous groups associated with known giant ellipticals are reported in Table 1 and 8. Test searches for additional X-ray clusters in the 9<sup>h</sup> \- 14<sup>h</sup> region have shown that both, the selection of extended sources by the RASS standard analysis and the finding of X-ray clusters by selecting extended sources is quite incomplete in terms of the compilation of flux-limited X-ray cluster catalogues - except for very high flux limits. We have shown that probably most of the missing clusters can be recovered by using a better search algorithm for extended sources and by screening the optically known clusters. This may not be sufficient to produce a high quality, highly complete catalogue. Further screening of RASS sources, including CCD imaging of sources not visibly extended in X-rays, is required. The completion of the NORAS sample thus requires additional imaging and spectroscopic observations. ###### Acknowledgements. We like to thank first of all the ROSAT Team who helped in providing the RASS data fields, in discussions about the RASS source properties, and by providing the EXSAS software which was essential for the studies. We like to thank in particular J. Engelhauser, A. Vogler, C. Rosso, and C. Izzo for help with the data and programs. We would also like to thank the remote observers at FLWO, Perry Berlind and Jim Peters, for observations made on the 1.5-m telescope, and Susan Tokraz for help in the data reduction. This work has made use of the SIMBAD data base operated at CDS, Strasbourg. In addition this research also made use of the NASA/IPAC Extragalactic Database (NED), which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with NASA. H.B. and P.S. acknowledge the support by the Verbundforschung under grant No. 50 OR 93065 and 50 OR 970835, respectively. JPH and JM were supported by the Smithonian Institution.
warning/0003/hep-ph0003130.html	ar5iv	text	# Scaling of the 𝐵 and 𝐷 meson spectrum in lattice QCD ## I Introduction Mesonic bound states consisting of a single heavy quark, $`b`$ or $`c`$, and a light quark, $`u`$, $`d`$ or $`s`$, as well as gluons, provide an interesting laboratory to study strong interactions. The typical momentum within such states is much lower than the mass of the heavy quark. This leads to a situation where the heavy quark becomes non-relativistic and the properties of the bound state are essentially determined by the light quark and the glue. At leading order the splittings within the spectrum become independent of the properties of the heavy quark, such as its mass $`m_Q`$ and spin $`s_Q`$, so that orbital and radial excitation energies are expected to match between the $`B`$-system and the $`D`$-system. The resulting approximate SU($`2N_h`$) symmetry, with $`N_h`$ denoting the number of heavy flavours, is usually referred to as heavy quark symmetry, see and the references therein. At the next order, $`1/m_Q`$ effects give rise to fine structure in the spectrum, several times larger in the $`D`$-system than for the $`B`$, see e.g. for a review. The spectrum of $`B`$ and $`D`$ states is not yet well established experimentally although several new results have been reported recently . Here we study the spectrum theoretically and from first principles using lattice QCD. This will aid the experimental search for new states. In the case of well-established states it will provide a test for the theory and/or the systematic errors in our calculation. Of key interest are decay matrix elements for $`B`$-factory experiments. Knowing how well the spectrum has been obtained gives confidence that we understand how to simulate $`B`$ and $`D`$ mesons reliably. This is important for the analysis of systematic errors in matrix element determinations. To formulate heavy $`b`$ and $`c`$ quarks on the lattice, a naïve discretisation is inappropriate since the lattice spacings currently available are not small compared to the Compton wave length of those quarks ($`m_Qa>1`$). Presently there are two different formulations available to simulate heavy quarks, non-relativistic QCD (NRQCD) and the heavy Wilson approach . For the $`b`$-quark on present lattices both approaches become essentially the same. However, in this regime, NRQCD is to be preferred since the inclusion of higher order correction terms is easily implemented. In this publication we report on our calculations of the $`B`$-meson spectrum for two different values of the lattice spacing $`a`$. Together with the results of , which were obtained with the same methods at another value of the lattice spacing, we can investigate the dependence on $`a`$ of our results. Physical results must be independent of $`a`$ and hence we can perform a test of systematic errors inherent in our calculation. We find no such errors at a significant level. In addition, on our coarsest lattice, we were able to simulate the $`D`$-meson spectrum and compare to results using heavy Wilson methods on finer lattices (where NRQCD does not work well since $`am_c<1`$). Early results on our coarse lattice have already been published in . Section II gives details of the simulations we performed and section III gives details of our fit procedure. Section IV gives our determination of the bare $`b`$ and $`c`$ quark masses. Section V discusses the behaviour of the splittings in the spectrum that we obtain. This includes fits to the dependence of the splittings on the mass of the heavy quark. Section VI compares the results in physical units at different values of the lattice spacing and with previous results as well as with experiment. Readers interested in our results for the physical meson spectrum could jump directly to this section. Section VII contains our conclusions and our best estimate for the $`B`$-spectrum, based on the combined input from three different values of the lattice spacing. ## II Simulation details ### A Gauge field action Our calculation was performed on two sets of gauge field configurations, which were generated using the Wilson gauge action $$S_G=\beta \underset{x,\mu <\nu }{}[1\frac{1}{3}\mathrm{Re}\mathrm{Tr}(U_{x,\nu }U_{x+\widehat{\nu },\mu }U_{x+\widehat{\mu },\nu }^+U_{x,\mu }^+)].$$ (1) This action has lattice artifacts of $`𝒪(a^2)`$. For the bare gauge coupling $`\beta `$, we used $`5.7`$ and $`6.2`$. The lattice volumes and the number of configurations are given in table I. We will refer to these configurations by their respective $`\beta `$-values. ### B Light quark propagators The light quark propagators have been generated with the use of the Sheikholeslami-Wohlert action, also known as the clover action , $`S_L`$ $`=`$ $`a^4{\displaystyle \underset{x}{}}[\overline{\psi }_x\psi _x+\kappa {\displaystyle \underset{\mu }{}}[\overline{\psi }_{x\widehat{\mu }}(\gamma _\mu 1)U_{x\widehat{\mu },\mu }\psi _x\overline{\psi }_{x+\widehat{\mu }}(\gamma _\mu +1)U_{x,\mu }^+\psi _x]`$ (3) $`a\frac{1}{2}\mathrm{i}c_{\mathrm{sw}}\kappa {\displaystyle \underset{\nu ,\rho }{}}\overline{\psi }_xF_{\nu \rho ,x}\sigma _{\nu \rho }\psi _x].`$ On the configuration set with $`\beta =5.7`$ the clover coefficient $`c_{\mathrm{sw}}`$ is set to its tadpole-improved tree level value $`c_{\mathrm{sw}}=1.5667`$, as determined from the 4th root of the plaquette . This reduces the lattice spacing artifacts in the light quark propagators to $`𝒪(\alpha _sa,a^2)`$. At $`\beta =6.2`$ we used the non-perturbative determined value of $`c_{sw}=1.6138`$, which removes the $`𝒪(\alpha _sa)`$ artifacts from the light quark propagator as well. In reference the light hadron spectrum at $`\beta =6.2`$ has been calculated using the non-perturbative as well as the tadpole-improved tree level value for $`c_{\mathrm{sw}}`$. No significant differences in the meson and baryon spectrum could be resolved between the two values of $`c_{\mathrm{sw}}`$. From this we expect the difference between tadpole and non-perturbatively improved light quarks at $`\beta =6.2`$ to be well covered by the size of the statistical errors in our case as well. This allows us to compare our $`\beta =6.2`$ results to the tadpole-improved results at $`\beta =5.7`$ and in reference . For each value of $`\beta `$ we used 3 different values for the hopping parameter $`\kappa `$. The actual values are detailed in table II. The table also contains the values of $`\kappa _c`$ and $`\kappa _s`$ from the UKQCD collaboration used in our calculation. The use of these values is appropriate for the analysis in terms of chiral extrapolations and scale setting that we have done. We also carefully include systematic errors from different chiral extrapolations and associated uncertainties in setting the scale. A recent re-analysis by UKQCD of their light hadron spectrum gives somewhat different values for $`\kappa _c`$ and $`\kappa _s`$. Our errors encompass any changes this would produce in our physical results. ### C Heavy quark propagators The typical momentum scale inside a heavy light meson such as a $`B`$ or $`D`$ meson is of the $`𝒪(\mathrm{\Lambda }_{\mathrm{QCD}})`$, which is small compared to the mass of the heavy quark. Therefore the mass of the heavy quark $`m_Q`$ represents an irrelevant scale for the dynamics of the mesonic bound state and it is possible to simulate these states on lattices with a lattice spacing larger than the Compton wavelength of the heavy quark. In our simulation we use a non-relativistic expansion of the heavy quark Hamiltonian, which is known as NRQCD . $`H`$ $`=`$ $`H_0+\delta H,`$ (5) $`H_0`$ $`:=`$ $`{\displaystyle \frac{𝐃^2}{2m_Q}},`$ (6) $`\delta H`$ $`:=`$ $`c_4{\displaystyle \frac{g}{2m_Q}}𝝈𝑩+c_2{\displaystyle \frac{\mathrm{i}g}{8m_Q^2}}(𝐃𝑬𝑬𝐃)c_3{\displaystyle \frac{g}{8m_Q^2}}𝝈(𝐃\times 𝑬𝑬\times 𝐃)`$ (8) $`c_1{\displaystyle \frac{(𝐃^2)^2}{8m_Q^3}}+c_5a^2{\displaystyle \frac{𝐃^{(4)}}{24m_Q}}c_6a{\displaystyle \frac{(𝐃^2)^2}{16nm_Q^2}}.`$ Please note that the rest mass term of $`H`$ has been omitted, resulting in a shift of the Hamiltonian, which is discussed in section IV. In the case of a heavy-light meson the NRQCD expansion has to be organised in powers of $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q`$ . Here this expansion is used up to $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^2)`$. We also include the $`𝒑^4`$ term, which is believed to be the leading term in $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^3)`$. The last two terms correct for discretisation errors from finite lattice spacing in respectively the spatial and temporal derivatives. $`n`$ is a stability parameter used in the evolution equation (II C). The matching coefficients $`c_1`$, …, $`c_6`$ are set to their tadpole-improved tree level values . With the Hamiltonian $`H`$ and $`\delta H`$ the propagator of the heavy quark can be obtained from a Schrödinger-type evolution equation $`G_{t+1}`$ $`=`$ $`\left(1a\frac{1}{2}\delta H\right)\left(1a\frac{1}{2n}H_0\right)^nU_4^+\left(1a\frac{1}{2n}H_0\right)^n\left(1a\frac{1}{2}\delta H\right)G_t\text{ for }t>1,`$ (10) $`G_1`$ $`=`$ $`\left(1a\frac{1}{2}\delta H\right)\left(1a\frac{1}{2n}H_0\right)^nU_4^+\left(1a\frac{1}{2n}H_0\right)^n\left(1a\frac{1}{2}\delta H\right)\varphi _x.`$ (11) With $`\varphi _x`$ we denote the source smearing function used on the initial time slice. At $`\beta =5.7`$ we use 20 different values for $`m_Q`$ in the range $`0.6am_Q20.0`$ and at $`\beta =6.2`$ we use 10 values in the range $`1.1am_Q6.0`$. Details, including the $`n`$ values, are given in the tables III and IV. For each value of $`\beta `$ we performed 3 different runs. At $`\beta =5.7`$ we label them A, C and S; for $`\beta =6.2`$ they are labeled H, N and P. For the $`S`$-wave mesons at $`\beta =5.7`$ we used up to three different smearing functions, $`\varphi _{\mathrm{G},0}`$, $`\varphi _{\mathrm{G},1}`$ and $`\varphi _{\mathrm{G},2}`$, in the different runs. These are convolutions of Gaussian functions for the light and the heavy quark with radii as detailed in table V. The configurations were fixed to Coulomb gauge. A local sink will be denoted with $`\varphi _\mathrm{L}`$. In most cases our final $`\beta =5.7`$ results were obtained with both sink and source smearing. For $`\beta =6.2`$ we use smearing for the heavy quark propagators only. In run H at $`\beta =6.2`$ we applied a hybrid procedure of Jacobi smearing and fuzzing . For runs N and P we fixed the configurations to Coulomb gauge. We used hydrogenic wave functions $`\varphi _{\mathrm{Hg},1}`$, $`\varphi _{\mathrm{Hg},2}`$ and $`\varphi _{\mathrm{He},1}`$ for run N. The indices ‘g’ and ‘e’ denote wave functions of the ground and first excited state. The details are given in table VI. In the P run we used Gaussian smearing with two different radii, $`ar_Q=2.5`$ and $`5.0`$. The spin operators applied to construct mesonic states with the correct quantum numbers are detailed in table 1 of reference . ### D Lattice spacing In the quenched approximation one obtains different values for the lattice spacing, depending on the quantity it is determined from. This is expected to be caused by the strong coupling $`\alpha _s`$ running differently in the real world and the quenched theory. We use the physical mass of the $`\rho `$-meson to fix the lattice spacing. This procedure is justified from the typical gluon momentum in a $`B`$ or $`D`$ meson being of similar size to the momentum in a light meson such as the $`\pi `$ and $`\rho `$. Since heavyonium states probe a higher physical scale these are not appropriate to fix the scale for a heavy-light system in the quenched approximation. Using the $`\rho `$-scale should take care of most of the quenching effects. The determination of $`m_\rho `$ is complicated by the chiral extrapolation required, see reference for a review. At $`\beta =5.7`$ we use the result of . The result of the linear extrapolation in the light quark mass $`m_q`$ is quoted as the central value and the deviation of the quadratic fit is treated as a systematic uncertainty. At $`\beta =6.2`$ a linear extrapolation is reported in reference . We treat the difference to the 3rd order extrapolation from as a systematic uncertainty. The numbers are compiled in table VII. We use: $`\beta =5.7:a^1=1.116(12)(_0^{+56})\text{ GeV},a=0.1768(19)(_{88}^{+0})\text{ fm},`$ (13) $`\beta =6.2:a^1=2.59(_{10}^{+6})(_0^{+9})\text{ GeV},a=0.0762(_{18}^{+29})(_0^{26})\text{ fm}.`$ (14) For comparison, table VII also shows the lattice spacing as obtained from the string tension $`\sigma `$ and the bottomonium splitting $`\overline{\chi }_b\mathrm{{\rm Y}}`$. As a physical value for $`\sigma `$ we choose a result obtained from a potential model fit to the charmonium spectrum . The lattice numbers originate from . These results are in agreement with the outcome of the $`m_\rho `$ analysis. As explained above, the bottomonium system probes a different scale and the values obtained using it do not agree with the result from light spectroscopy . ## III Fitting techniques ### A Parametrisations At $`\beta =5.7`$ we used several different smearings at source and sink. For hadron correlators with a local sink, we applied simultaneous vector fits, requiring the fitted mass(es) $`m_k`$ to agree for all propagators: $`\varphi _\mathrm{L}(t)\|\varphi _i(0)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}A_{i,k}\mathrm{exp}(m_kt),1im,`$ (15) $`A_{i,k}`$ $`=`$ $`\varphi _\mathrm{L}\|\psi _k\psi _k\|\varphi _i.`$ (16) In the case of sink and source smearing, we used simultaneous matrix fits. In matrix fits, the fitted amplitudes are constrained in their relationship with each other as well: $`\varphi _j(t)\|\varphi _i(0)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}B_{j,k}^{}B_{i,k}\mathrm{exp}(m_kt),1j,im,`$ (17) $`B_{i,k}`$ $`=`$ $`\psi _k\|\varphi _i.`$ (18) The fitting techniques are described in more detail in reference . We found matrix fits to be more precise with respect to statistical errors. Due to the omission of the rest mass in eq. (II C) the fitted mass is shifted with respect to the bound state mass. We denote the result of the fit as the simulation mass, $`m_{\mathrm{sim}}`$. The determination of the shift will be discussed in section IV. To extract mass splittings we applied two different procedures. One is to fit the masses as above, take their difference, and then calculate the error from the bootstrap or jackknife samples of the difference. With this procedure one can easily take advantage of using different smearings. In the case of a single smearing function, a ratio-fit provides an alternative . For this one divides the bootstrap or jackknife samples of the two propagators and fits the outcome with an exponential ansatz. The mass shift cancels out of the difference in both procedures. ### B Pseudo-scalar and vector meson On the $`\beta =5.7`$ configurations the simulation masses for pseudo-scalar and vector mesons have been determined most accurately in run A. In this run we only used the smearing functions $`\varphi _{\mathrm{G},1}`$ and $`\varphi _{\mathrm{G},2}`$. We found the double exponential matrix fit with sink and source smearing to deliver the most precise result. For the fit range we choose the initial time slice $`t_{\mathrm{min}}`$ two time slices larger than the first time slice delivering a reasonable $`\chi ^2`$. In general we choose the number of dropped time slices multiplied by the number of propagators used for the fit to be larger than or equal to the number of fit parameters. The reason for this procedure is as follows. The first reasonable value of $`\chi ^2`$ is observed once the residual excitations are just masked by the statistical uncertainties, which allows for them to be still of similar size. Each excited data point can eat up one fit parameter. Dropping as many data points as fit parameters delivers a fit which is entirely dominated by statistical fluctuations. The residual fit range dependence of those fits becomes negligible against the statistical uncertainties. We judge $`\chi ^2`$ values resulting in a $`Q0.1`$ as reasonable, where $`Q`$ denotes the probability of a fit having an even higher value of $`\chi ^2`$. The final result is given in table VIII. In run H at $`\beta =6.2`$ we only had one smearing function available. We extracted the final results from single exponential fits to the propagators with source and sink smearing. Their fit results turned out to be more precise than the ones from using a local sink. The final fit range was determined such that we observed a reasonable $`\chi ^2`$ and achieved stability of the fitted result against variation of the fit range. The results are displayed in table IX. In the run N we used hydrogenic wave functions of different radii. We generated smeared local and smeared smeared meson propagators. However no cross correlators, e.g. $`\varphi _{\mathrm{Hg},1}`$ at sink and $`\varphi _{\mathrm{He},1}`$ at source, were calculated. Hence eq. (17) was inapplicable and we had to use vector fits in the case of smearing at sink and source as well. In double exponential vector fits to two smearing functions, we observed extremely low values of $`Q`$. We observed that this is connected to unfortunate statistical fluctuations on certain time slices. However the fit parameters turned out to be stable with respect to variations of the fit range. These fits will be discussed in subsection V B in more detail. To obtain a more precise result for the $`S`$-wave ground states, we resorted to single exponential fits to single propagators and compared the outcome for the different smearing functions. This is shown in figure 1 for the pseudo-scalar propagator at $`am_Q=2.5`$. The octagons indicate the final result for each propagator, as determined from the $`Q`$-value after dropping two time slices. Within statistical errors all results are in reasonable agreement with each other. For the final result, which is also included in table IX, we choose the smeared-smeared $`\varphi _{\mathrm{Hg},1}`$ propagator. In the end these deliver the more accurate hyperfine splitting, due to superior noise cancellation between the pseudo-scalar and the vector meson state. In this context it is interesting to note that the propagators with local sink and $`\varphi _{\mathrm{Hg},2}`$ source smearing plateau much later than the others, but the results are in agreement with those from other propagators. For the $`S`$ wave states in the run P we only had the Gaussian smearing at the source with radius $`ar_0=2.5`$ and local sink. Since these propagators plateau quite late, we used $`t_{\mathrm{min}}=16`$ for the final result, the error bars are not competitive with those above. Since they are needed for the later analysis of the $`P`$ states we include them as well in table IX. To describe physical bound states involving light $`u`$ and $`d`$ quarks, the results of tables VIII and IX have to be extrapolated in the light quark hopping parameter. On both sets of configurations, the difference between the critical and normal hopping parameter is smaller than the uncertainty we assigned to $`\kappa _c`$ in table II and we use $`\kappa _c`$ in our extrapolations. The normal hopping parameter is the one for which the extrapolations deliver the physical $`m_\pi /m_\rho `$ ratio. Due to the high statistical accuracy we achieved at $`\beta =5.7`$ in the pseudo-scalar case, a linear ansatz in $`am_q:=\frac{1}{2}(\frac{1}{\kappa }\frac{1}{\kappa _c})`$ in a full covariant fit to all three data points, results in a fit with $`\chi ^2/\mathrm{d}.\mathrm{o}.\mathrm{f}.>8/1`$ for $`am_Q<10`$. This corresponds to $`Q<0.004`$. The resulting curves do not describe the data. We carefully checked whether this is caused by a residual fit range dependence and found all the fit parameters including the would-be strange to non-strange meson splitting to be stable against variation of the fit range. This was done for an initial time slice $`t_{\mathrm{min}}`$ in a range from $`3`$ to $`6`$. We therefore extracted our final result from a linear spline to the points with highest and lowest $`m_q`$ and use the deviation of a quadratic spline as a systematic uncertainty of the chiral extrapolation. An example for the extrapolation is given in figure 2. From the figure it is obvious that interpolations to extract the heavy-strange meson mass are insensitive to the different ansätze and we do not assign an uncertainty due to the different interpolations. However, in the case of the heavy strange meson, we are faced with the problem that $`\kappa _s`$ is highly sensitive to the quantity it is determined from. Our central value is interpolated to the $`\kappa `$ as determined from $`m_K/m_\rho `$, and the difference to the outcome for $`\kappa `$ corresponding to $`m_\varphi /m_\rho `$ is treated as an uncertainty of the quenched approximation. The results are presented in table X. For $`\beta =6.2`$ the statistical accuracy is not as high and our data are well described by linear extrapolations. The results are presented in table XI. ## IV Heavy quark masses ### A Mass shift from dispersion relation The omission of the rest mass term $`m_Q`$ in the Hamiltonian eq. (II C) causes most of the shift of the simulation mass $`m_{\mathrm{sim}}`$ with respect to the physical meson mass. The mass, $`m_{\mathrm{rel}}`$, of the meson can be determined from the relativistic dispersion relation of the meson $`E(\stackrel{}{p})=\sqrt{m_{\mathrm{rel}}^2+\stackrel{}{p}^2}`$, which gives $$m_{\mathrm{rel}}=\frac{\stackrel{}{p}^2[E(\stackrel{}{p})E(0)]^2}{2[E(\stackrel{}{p})E(0)]}.$$ (19) Here $`E(\stackrel{}{p})`$ denotes the total energy of the meson. The mass shift $`\mathrm{\Delta }_{\mathrm{rel}}`$ is defined as the difference $$\mathrm{\Delta }_{\mathrm{rel}}:=m_{\mathrm{rel}}m_{\mathrm{sim}}.$$ (20) This shift per heavy quark should be universal for all hadronic states simulated at the bare heavy quark mass $`m_Q`$. In our calculation at $`\beta =5.7`$ we determined the mass shift from the difference in energy of the pseudo-scalar meson propagators with $`a\|\stackrel{}{p}\|=0`$ and $`2\pi /12`$. This was done in run C at $`\kappa =0.1400`$ with source smearing $`\varphi _{\mathrm{G},1}`$ and a local sink. At large values of $`m_Q`$, we found a single exponential fit to the ratio of the correlators to plateau much later than the fits to the individual propagators. This is reflected in a large fit range dependence of the jackknife difference of the masses of the individual fits, for time slices in which no plateau was observed in the ratio-fit. For our final result we choose a minimal $`t`$-value two time slices larger than the first $`t`$-value for which we obtained a decent $`\chi ^2`$ in a fit to the ratio of propagators. The final result is presented in table XII and figure 3. We also tried simultaneous vector fits according to eq. (15) with two exponents. We used propagators with source smearing $`\varphi _{\mathrm{G},1}`$ and $`\varphi _{\mathrm{G},2}`$. The jackknifed difference of the fitted ground state mass is in agreement with the above procedure; however, the statistical uncertainties, especially for large values of $`m_Q`$, are larger. For $`\beta =6.2`$ we calculated the mass shift in heavy quarkonia, since the statistical precision for heavy-light correlators at finite momentum was not sufficient. In the following, mass shifts from heavy quarkonia will be denoted by $`\mathrm{\Delta }_H`$. We simulated the vector-meson for $`a\|\stackrel{}{p}\|2\frac{2\pi }{24}`$. The kinetic mass $`m_1`$ was obtained from fits to the dispersion relations: $`E_{\mathrm{sim}}(\stackrel{}{p})`$ $`=`$ $`m_0+{\displaystyle \frac{\stackrel{}{p}^2}{2m_1}}{\displaystyle \frac{\stackrel{}{p}^4}{8m_2^3}},`$ (22) $`E_{\mathrm{sim}}(\stackrel{}{p})`$ $`=`$ $`m_0+{\displaystyle \frac{\stackrel{}{p}^2}{2m_1}}{\displaystyle \frac{\stackrel{}{p}^4}{8m_1^3}},`$ (23) $`E_{\mathrm{sim}}(\stackrel{}{p})`$ $`=`$ $`m_0+{\displaystyle \frac{\stackrel{}{p}^2}{2m_1}},`$ (24) with parameters $`m_0`$, $`m_1`$ and $`m_2`$. $`E_{\mathrm{sim}}(\stackrel{}{p})`$ denotes the simulation energy as determined from the propagator falloff. In the case of $`am_Q1.3`$ we used the ansätze (22) and (23); for the three heavier $`m_Q`$-values, (23) and (24). All fits gave fit parameters which were consistent within half of the statistical error. To obtain the shifts required for heavy-light spectroscopy we subtracted the simulation mass of the quarkonium vector-meson and divided by two. The final results are displayed in table XIII and figure 3. It is interesting to compare to the result from reference $`a\mathrm{\Delta }_H=1.29(2)`$ obtained at $`am_Q=1.22`$. Due to higher statistics, this result is much more precise. This value is included as a square into figure 3 and agrees well with the newer results. ### B Mass shift in perturbation theory The mass shift $`\mathrm{\Delta }`$ can also be calculated in lattice perturbation theory : $$\mathrm{\Delta }_{\mathrm{pert}}=Z_mm_QE_0.$$ (25) Here $`Z_m`$ denotes the renormalisation constant connecting the bare lattice mass $`m_Q`$ with the pole mass and $`E_0`$ denotes the heavy quark self energy constant. In the perturbative expansion the 1-loop contributions from $`Z_m`$ and $`E_0`$ cancel each other to a large extent and the direct perturbative expansion of $`\mathrm{\Delta }_{\mathrm{pert}}`$ is much better behaved than either perturbative series on its own. The Lepage Mackenzie scale $`aq^{}`$ has been determined separately for $`\mathrm{\Delta }`$ and it is larger than for $`Z_m`$ or for $`E_0`$. The coefficients for $$\mathrm{\Delta }_{\mathrm{pert}}=m_Q[1+\alpha _s(aq^{})\mathrm{\Delta }^{(1)}]$$ (26) can be found in table XIV. We use the $`\alpha _P(aq=3.4)`$ values as determined from the $`1\times 1`$ Wilson loop with 2-loop running in order to evolve to the respective $`aq^{}`$. For the final mass shift we assign a relative uncertainty of $`\alpha _s^2(aq^{})`$. Since $`\mathrm{\Delta }^{(1)}`$ is small, this is more conservative than the squared 1-loop contribution. The final results are displayed in table XII for $`\beta =5.7`$ and table XIII for $`\beta =6.2`$. For values of $`m_Q`$ not included in table XIV we interpolated linearly between the nearby values, which is completely sufficient within the claimed accuracy. The results for $`a\mathrm{\Delta }`$ from perturbation theory and the lattice simulation are compared in figure 3. Apart from possibly the low $`m_Q`$ region at $`\beta =5.7`$, the figure shows excellent agreement between the two ways of calculating the mass shift. For $`\beta =5.7`$ the stability parameter $`n`$ differs in some cases between the perturbative results and the simulation. However for $`am_Q=4`$, where perturbative results exist for $`n=1`$ and $`2`$, the effect of $`n`$ is completely negligible: we obtain $`\mathrm{\Delta }_{\mathrm{pert}}=3.88(22)`$ vs $`3.89(24)`$. From a comparison of the simulation result of the runs A and C at $`\beta =5.7`$ we can also obtain evidence of the effect of the different $`n`$ on the simulation mass $`m_{\mathrm{sim}}`$. For $`am_Q=1.0`$ and $`\kappa =0.1400`$ we measure $`am_{\mathrm{sim},\mathrm{ps}}=0.6265(21)`$ and $`0.6248(21)`$ for $`n=5`$ and $`6`$ respectively. This difference is again completely negligible against the uncertainty we assign to $`a\mathrm{\Delta }`$, even if we enlarge it by a factor of $`3`$ to allow for a larger effect between $`n=4`$ and $`5`$. The former $`n`$ was used in the perturbation theory. Note also that this difference tends to be in the opposite direction to that in $`a\mathrm{\Delta }`$ implying that the effect of $`n`$ on the physical mass is reduced when compared to the shift. Here it is interesting to note that physical mass differences like the hyperfine splitting $`m_{\mathrm{hpf}}=m_{\mathrm{sim},\mathrm{v}}m_{\mathrm{sim},\mathrm{ps}}`$ are even less sensitive to $`n`$. At the above mass parameter of $`am_Q=1.0`$ we measure $`am_{\mathrm{hpf}}=0.0833(20)`$ for $`n=4`$ and $`0.0835(20)`$ for $`n=5`$. In summary the differences in $`n`$ between the different runs as well as the perturbative shifts can be neglected safely even at the high level of accuracy we achieved here. This leaves us with a discrepancy between $`\mathrm{\Delta }_{\mathrm{pert}}`$ and $`\mathrm{\Delta }_{\mathrm{rel}}`$ for our lowest $`m_Q`$-values, which is roughly twice as large as the uncertainty we assign to the perturbative result. On the other hand for $`\beta =6.2`$ we observe excellent agreement between the precise result of with the perturbative calculation at the relatively low value $`am_Q=1.22`$. Given a value for the shift $`\mathrm{\Delta }`$ and the simulation mass $`m_{\mathrm{sim}}`$ from tables VIII, IX, XII and XIII, we can now calculate absolute masses for all the states. We do this for the ground state vector and pseudo-scalar mesons, both to fix the quark mass, as described in the next subsection. Moreover, we use the meson mass rather than the quark mass to discuss the $`m_Q`$ dependence, since it is more directly comparable to experiment. We frequently plot results against $`1/m_{\mathrm{sav}}`$, where $`m_{\mathrm{sav}}`$ is the spin-average of the ground state vector and pseudo-scalar mesons $$m_{\mathrm{sav}}=\frac{1}{4}(3m_\mathrm{v}+m_{\mathrm{ps}}).$$ (27) This is preferable to $`m_{\mathrm{ps}}`$ alone since the spin-averaging reduces the dependence on sub-leading spin-dependent terms. ### C Bare heavy quark mass To determine the bare quark mass $`m_Q`$ corresponding to the $`b`$ and $`c`$-quark, we compared the mass of the spin-averaged $`S`$-wave meson denoted with an overbar, with the experimental result. We used 5313 MeV for the $`\overline{B}`$, 5405 MeV for the $`\overline{B}_s`$, 1973 MeV for the $`\overline{D}`$ and 2076 MeV for the $`\overline{D}_s`$ . For the interpolations we used spline-fits to three neighbouring points. The fits were done quadratically in $`m_Q`$ and $`1/m_Q`$ and no significant difference was observed between the two. From the strange and non-strange meson we obtained identical results for the quark masses: $`am_b`$ $`=`$ $`1.64(5)(_5^{+8})(_7^{+0}),\beta =6.2`$ (29) $`am_b`$ $`=`$ $`4.20(25)(5)(_{24}^{+0}),\beta =5.7`$ (30) $`am_c`$ $`=`$ $`0.87(6)(3)(_{13}^{+0}),\beta =5.7`$ (31) The errors as indicated in the parentheses give the uncertainty arising from the mass shift and the statistical and systematic uncertainty of the $`a`$ determination. The uncertainties associated with the simulation mass $`m_{\mathrm{sim}}`$ are completely negligible here. For $`m_b`$ we used the perturbative shifts $`\mathrm{\Delta }_{\mathrm{pert}}`$. Using $`\mathrm{\Delta }_H`$ at $`\beta =6.2`$ delivers $`am_b=1.59(_5^{+14})(_3^{+6})(_5^{+0})`$ and using $`\mathrm{\Delta }_{\mathrm{rel}}`$ at $`\beta =5.7`$ gives $`am_b=4.16(_{31}^{+53})(7)(_{30}^{+0})`$, which is agreement with $`\mathrm{\Delta }_{\mathrm{pert}}`$ but with larger error. For $`m_c`$ we used the simulation result $`\mathrm{\Delta }_{\mathrm{rel}}`$. Here $`\mathrm{\Delta }_{\mathrm{pert}}`$ would give $`am_c=1.02(8)(2)(_{10}^{+0})`$. The deviation from the result eq. (31) reflects the difference between the $`\mathrm{\Delta }`$-values at low $`am_Q`$ discussed above. This careful analysis to fix the bare heavy quark mass is particularly necessary for fine structure splittings in the spectrum to be discussed in the next section. These are very sensitive to the heavy quark mass, generally as $`1/m_Q`$. In addition, any errors in the heavy quark mass must be fed into errors in the fine structure splittings in order to avoid underestimating those errors. The bare masses do not scale with the lattice spacing as expected, because they are unphysical. A better quantity to consider would be the mass in the $`\overline{\mathrm{MS}}`$-scheme. This will be discussed in a future publication . ## V Mass dependence of level splittings In this section we describe how the results for the level splittings are extracted from our data. The dependence of the different splittings on the light and heavy quark mass is also discussed. ### A Flavour dependent splittings The mass difference between heavy-light states distinguished only by their strangeness survives into the static limit. Based on the ideas of heavy quark symmetry such splittings are expected to depend weakly on the mass of the heavy quark. If a spin-averaged combination is taken, the leading heavy quark mass dependence arises purely from the kinetic term in eq. (6). The size of the slope in $`1/m_{\mathrm{sav}}`$ then gives information on the difference in $`p_b^2`$ for the $`b`$-quark in the strange and non-strange states. We calculated this splitting from the ground state S-wave results for $`\kappa _s`$ and $`\kappa _c`$. The result is highly sensitive to the reported uncertainties in the chiral extrapolation and the determination of the strange hopping parameter $`\kappa _s`$. At $`\beta =5.7`$ we determine the statistical uncertainties in a jackknife procedure applied to the difference of the individual results and at $`\beta =6.2`$ we use the bootstrap. For $`\beta =5.7`$ our statistical errors are very small and we consider additional systematic uncertainties for the chiral extrapolation and $`\kappa _s`$. For $`\beta =6.2`$ the quality of our data is not as good and we give statistical errors only. Using the $`\kappa _s`$-value determined from the $`\varphi `$ would lead to a 9% increase of the result, which is small compared to our statistical errors. The results are displayed in tables XV and XVI. In figure 4 we plot the result for the spin-averaged splitting at $`\beta =5.7`$ versus the inverse of $`m_{\mathrm{sav}}`$, in order to display its heavy quark dependence. The figure displays a clear increase of the splitting with decreasing heavy quark mass. To quantify the slope of this dependence we perform a linear fit of the splitting result versus $`1/m_{\mathrm{sav}}`$. The result, converted into physical units, is detailed in table XXV. The slope corresponds to a $`p_b^2`$ difference of $`[0.25(3)\text{ GeV}]^2`$, which is of the size of $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$, as expected. Because of the larger statistical uncertainties, we do not observe a significant slope at $`\beta =6.2`$. The data can be fitted nicely to a constant. A comparison of our results with the ones of for the pseudo-scalar case is plotted in figure 5. In this plot we show the result for the strange quark as determined from the $`K/\rho `$ mass ratio only. Due to the large error bars we do not include the results obtained at the larger values of $`m_Q`$ for $`\beta =6.2`$. Within the accuracy of around 12% in the case of $`\beta =6.2`$ or better, no sign of scaling violations shows up in the plot. We also observe excellent agreement with the experimental result. ### B Radial excitations In order to obtain a reasonably stable and long plateau for the radially excited $`S`$-states on the coarse lattice at $`\beta =5.7`$ we applied triple exponential matrix fits to the three smearing functions $`\varphi _{\mathrm{G},0}`$, $`\varphi _{\mathrm{G},1}`$ and $`\varphi _{\mathrm{G},2}`$. This was done for the run S for a single $`\kappa `$ of $`0.1400`$ only, which is approximately equal to the strange as determined from the $`K`$-meson. Since in reference the dependence of the $`2^1S_01^1S_0`$ splitting on the light quark mass was found to be very small, a variation of less than 2% when fixing $`\kappa _s`$ from the $`K`$ or $`K^{}`$-meson, we can ignore any mismatch in our $`\kappa `$ vs $`\kappa _s`$ compared to the statistical uncertainties. We therefore treat our result as the answer for this splitting with $`\kappa _s`$ as determined from the $`K`$. In figure 6 we show a typical example for the excellent stability of the simulation masses $`am_{\mathrm{sim}}`$ against variation of the starting point $`t_{\mathrm{min}}`$ of the fit range. The extent in $`t_{\mathrm{min}}`$ for which we can resolve the excited state is 5 time slices or $`0.28`$ GeV<sup>-1</sup>. The rate of its disappearance is set by the $`2S1S`$ splitting of $`600`$ MeV. In the figure the first good value of $`Q`$ is observed for $`t_{\mathrm{min}}=2`$. To be safe with respect to residual excitations we quote the final result for a fit range starting at time slice $`4`$, which is the procedure described in subsection III B. The peak in $`Q`$ at $`t_{\mathrm{min}}=5`$ results from the fit becoming insensitive to the third exponential at this point. The results for all 6 heavy quark masses are given in table XVII. In figure 7 we plot the heavy quark mass dependence of the spin-averaged splitting. The result shows a clear increase as the heavy quark mass is reduced. In table XXV we detail fit results for the offset and slope of this splitting with respect to $`1/m_{\mathrm{sav}}`$. From the assumption that the increase of the splitting with $`1/m_{\mathrm{sav}}`$ is caused by the difference in the kinetic energy $`p^2/2m_Q`$ between ground and radially excited states, the fitted slope gives a $`p_b^2`$ difference of $`[0.95(15)\text{ GeV}]^2`$. This is of the size of a few times $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ as would be expected. On our fine lattice, since no cross-correlators between the different smearings had been calculated, we used simultaneous vector fits in all cases. The differences between the smearings $`\varphi _{\mathrm{Hg},1}`$, $`\varphi _{\mathrm{He},1}`$ and $`\varphi _{\mathrm{Hg},1}`$ turned out to be too small for simultaneous fits with three exponentials and we had to resort to fits using 2 exponentials. Again we choose the starting point $`t_{\mathrm{min}}`$ of the fit range as described in subsection III B. With this procedure it is possible to extract reliable information on the excited state, as can be verified from the tables of reference for the case of $`\mathrm{{\rm Y}}`$-spectroscopy. Using propagators with sink and source smearing in vector fits, leads to a suppression of the excited state contamination, which made it impossible to extract a signal for the excited state. Therefore we used propagators with smearing at the source and local sinks to extract the radially excited states. As denoted earlier, the fits to these propagators are plagued by statistical fluctuations, which lead to quite large $`\chi ^2`$ and low $`Q`$ values. However, the fits describe the data reasonably and the fitted parameters are stable against variation of $`t_{\mathrm{min}}`$. This is shown in figure 8. For a fit range starting point $`t_{\mathrm{min}}5`$ we obtain $`Q>1\%`$, which is low compared to what we obtained in the other fits. However it is not that low, that the fit could be ruled out on statistical grounds. Together with the good stability of the fitted masses against variations of $`t_{\mathrm{min}}`$ as displayed on the right hand side of the figure, we believe that our signal is genuine. In this example we extract the final result from $`t_{\mathrm{min}}=8`$. The results for the radially excited $`S`$-wave for this lattice spacing are compiled in table XVIII. ### C Orbital excitations Orbitally excited $`P`$-state mesons have been investigated at both of our lattice spacings. The possible states consist of four non-degenerate energy levels; total angular momentum $`J=0`$ and $`2`$ as well as two $`J=1`$ states. In the heavy quark symmetry picture a $`jj`$ coupled basis is appropriate. In the vicinity of the static limit, the $`J=2`$ and the higher of the $`J=1`$ states are close and separated from the $`J=0`$ and lower $`J=1`$. The former correspond to a light quark total angular momentum of $`j_l=\frac{3}{2}`$, the latter of $`j_l=\frac{1}{2}`$. We use an $`LS`$ coupled basis to study the states, but expect our $`{}_{}{}^{1}P_{1}^{}`$ and $`{}_{}{}^{3}P_{1}^{}`$ channels to mix, leading to the observation of the lower $`J^P=1^+`$ state with both operators. We will denote the state corresponding to $`j_l=\frac{3}{2}`$ with a prime. At $`\beta =5.7`$ again we use one light hopping parameter, $`\kappa =0.1400`$. As in the case of the radial excitations, we treat this as the value corresponding to the strange quark as determined from the $`K`$-meson and the simulations have been performed in the run S. We used the derivatives of the smearings $`\varphi _{\mathrm{G},0}`$ and $`\varphi _{\mathrm{G},1}`$ at source and sink. The final results were obtained from double exponential simultaneous matrix fits to both smearings and are listed in table XIX. The selection of the fit range proved to be very delicate for this $`a`$-value. The statistical error grows rapidly when increasing $`t_{\mathrm{min}}`$, since the signal to noise is exponentially related to the $`PS`$ splitting . We give an example in section V F, where the fine structure is discussed. For the lightest values of $`am_Q`$ and correlators with $`{}_{}{}^{3}P_{2}^{}`$ wave operators in the $`T`$-representation we obtained small values of $`Q`$ of a few permille. We include the corresponding mass values into the table for the sake of completeness and mark them in italics. However we disregard them in the further evaluation. The results are always in agreement with the ones obtained in the $`E`$-representation and we do not believe there is a serious problem with this, simply statistical fluctuations. In the second to last column of table XIX we give the spin-averaged $`P`$-state result which we calculate as $$m(P_{\mathrm{sav}})=\frac{1}{12}[1m(^3P_0)+3m(^1P_1)+3m(^3P_1)+5m(^3P_2E)].$$ (32) The result is also shown in figure 9. For comparison we include the experimental result for the $`B_{sJ}^{}(5850)`$ resonance and the spin-average of the $`D_{s1}`$ and $`D_{s2}^{}`$ . The figure displays at most a mild heavy quark mass dependence. To quantify this, we report in table XXV on the offset and slope of this splitting in physical units. The slope is consistent with a $`p_b^2`$ difference of $`𝒪(\mathrm{\Lambda }_{\mathrm{QCD}}^2)`$, but is also consistent with zero. $`P`$-states were also investigated in the run P for $`\beta =6.2`$. We choose $`am_Q=1.6`$, directly corresponding to the $`b`$-quark in eq. (29). For the light quarks we use the strange hopping parameter $`\kappa =0.1346`$. Since the results on our coarse lattice depend only weakly on $`m_Q`$ we take the outcome as the final answer for $`B_s`$. In the simulations we used two different smeared sources together with local sinks. Again we used derivatives of Gaussian smearing functions. The masses were extracted from double exponential vector-fits to both propagators. We observe reasonable $`Q`$-values for all applied operators and include all channels into the spin-average $$m(P_{\mathrm{sav}})=\frac{1}{12}[1m(^3P_0)+3m(^1P_1)+3m(^3P_1)+2m(^3P_2E)+3m(^3P_2T)].$$ (33) The results for the fitted masses are displayed in table XX. The splitting between the spin-averaged $`P`$ and $`S`$ waves is given in tables XXI and XXII. ### D Radially excited $`P`$-states Having available 2 different smearing functions at both of our lattice spacings, it is possible to obtain information on the radially excited $`P`$-states as well. In figures 10 we show the dependence of the fitted masses of the spin-averages of the $`1P`$ and the $`2P`$ on the starting point $`t_{\mathrm{min}}`$ of the fit range for two different values of $`m_Q`$. The figure displays a clear signal for an excited state and reasonable stability with respect to variations of $`t_{\mathrm{min}}`$. However the error grows rapidly with increasing $`t_{\mathrm{min}}`$. The $`Q`$-values of the $`{}_{}{}^{3}P_{2}^{}E`$-fit, which is the last of the individual states included in the spin-average to reach a plateau, are included in figure 17. Decent $`Q`$-values are observed for $`t_{\mathrm{min}}=3`$. Since this is a 6 parameter matrix fit to four propagators, we take our final result from $`t_{\mathrm{min}}=5`$. The results for the spin-averaged $`2P`$-state are summarised in table XIX. We give the splitting to the spin-averaged $`1S`$ and $`1P`$-states in table XXI. We do not observe a significant slope for the splitting with respect to $`1/m_{\mathrm{sav}}`$. For $`\beta =6.2`$ we show the plateau in figure 11. Due to the finer lattice, the growth in the error with increasing $`t_{\mathrm{min}}`$ is smaller than before. For the mass parameters used here, an example of a $`Q`$-plot will be given in figure 18 below. Here the first decent $`Q`$ is observed for $`t_{\mathrm{min}}=2`$. Since this is a vector fit to 2 propagators, we take our final results from $`t_{\mathrm{min}}=6`$. The result and splittings are included in tables XX and XXII. As noted when discussing radially excited $`S`$-states, due to our conservative selection of the fit range, we expect residual excitations to be negligible within the quoted statistical errors. ### E Hyperfine splittings The mass difference between a pseudo-scalar and a vector $`S`$-wave meson is caused by the spin of the heavy quark. This hyperfine splitting is expected to vanish in the limit of infinitely heavy quark mass. On our $`\beta =5.7`$ lattice we determined the hyperfine splitting $`m_{\mathrm{hpf}}`$ from the difference of the results in table VIII. A crucial ingredient in obtaining a small statistical error is to choose identical fitting ranges for both correlators. If one of them has a plateau at a larger value of $`t_{\mathrm{min}}`$ than the other, we took this larger value to obtain $`m_{\mathrm{hpf}}`$ from the difference of the fitted masses. The error in this procedure is estimated with a jackknife. The results are displayed in table XXIII. The chiral extrapolation of the hyperfine splitting turned out to be less difficult than that for the simulation mass of the pseudo-scalar and vector mesons. The curvature seems to cancel out between them and we have been able to perform linear fits to extrapolate to the chiral limit. However, in order to be consistent, we assign a systematic uncertainty to the result. This uncertainty is obtained from the difference to the outcome of first extrapolating the individual mesons to the chiral limit and then determining the hyperfine splitting. In this case we use the quadratic extrapolations to the chiral limit to take a possible curvature into account. In our chiral extrapolations of the hyperfine splitting we observe a negative slope with respect to the mass $`am_q`$ of the light quark, which is illustrated in figure 12. The left hand side gives an example of our chiral extrapolations and the right hand side shows the slopes measured at each value of $`am_Q`$. In order to construct a physically meaningful quantity, the latter has been multiplied by the strange quark mass, such that we can compare with experimental results for the difference between the strange and non-strange hyperfine splittings. For $`B`$-mesons the light quark dependence of the hyperfine splitting is not well resolved experimentally, because of large uncertainties in the $`B_s`$ hyperfine splitting. For $`D`$-mesons the situation is much clearer and one observes an increase with the light quark mass. However the magnitude of the slope is largely dependent on whether you compare the $`D_s`$ hyperfine splitting with the hyperfine splitting of the $`D^+`$ or the $`D^0`$. We expect this difference in the experimental results to be mainly due to QED-effects, since these come in with opposite signs in the $`D^+`$ and the $`D^0`$. Since the $`D_s`$ is positively charged as well, QED effects should largely cancel when comparing the hyperfine splittings of the $`D_s`$ and the $`D^+`$ and one obtains a positive slope for the $`D`$-meson from the experiment. Comparing our data to the experimental results, one observes our hyperfine splittings to be too small. This will be discussed in more detail in section VI. With respect to the slope, the result at the $`D`$ has clearly the wrong sign and its magnitude is approximately twice as large as that from the experiment. We did not observe this effect in our $`\beta =6.2`$ results, neither was it observed in . Both of these results did not achieve the high statistical accuracy we have at $`\beta =5.7`$ and also use values of the heavy quark mass at around the $`B`$ or heavier. For those values of $`m_Q`$ the light quark mass dependence at $`\beta =5.7`$ is also not that significant. A slope of similar sign and size has been observed in the calculations of , although the authors did not comment on this. Reference used a highly improved gluonic action with NRQCD heavy quarks on even coarser lattices and a heavy clover action for the heavy quarks on a finer lattice $`\beta =6.0`$. A detailed comparison with these results will be given below in section VI. In this context it is interesting to note, the slope of the hyperfine splitting as a function of the quark mass turns out to be to small in light hadron spectroscopy as well . The calculations listed above are performed in the quenched approximation, which could be a factor contributing to the wrong slope. In potential model language, which is not necessarily appropriate here, the hyperfine splitting is related to the square of the wave-function at the origin. This in turn depends on the light quark mass and is independent of $`m_Q`$ as $`m_Q\mathrm{}`$. The wrong slope could then reflect the fact that the wave-function at the origin is not increasing rapidly enough as the light quark mass increases. This is natural in the quenched approximation as the potential at the origin is weakened by the coupling constant running too quickly to zero at short distance. An alternative scenario is one in which the coefficients of the relevant terms in the action, here $`c_{\mathrm{sw}}`$ in the light quark action, effectively carry some quark mass dependence that has not been included, leading to an underestimate of the hyperfine splitting at large $`am_q`$. In this case the effect would disappear as $`a`$ is reduced and this seems to be contradicted by results on finer lattices . Another cause could be a problem in the chiral extrapolation itself. The experimental result for hyperfine splitting $`J/\psi \eta _c`$ in the charmonium system is smaller than the hyperfine splitting for the $`D`$ and $`D_s`$-mesons. If one considers charmonium as a $`D_c`$-meson, one has to conclude that there is a maximum of the hyperfine splitting as a function of the light quark mass for $`m_q<m_c`$. If this maximum is attained for $`m_q<m_s`$, our observation of a negative slope of the hyperfine for $`m_qm_s`$ would be in agreement with nature. In this case extrapolations from the strange region to lighter $`u`$ and $`d`$-quarks as well as the chiral limit would be impossible. Our final results for the hyperfine splitting at $`\kappa _c`$ and $`\kappa _s`$ are given in table XXIII. In figure 13 we display the dependence of the hyperfine splitting on the spin-averaged heavy-light meson mass. A linear fit in $`m_{\mathrm{sav}}^1`$ for the five heaviest $`m_{\mathrm{sav}}`$ values gives reasonable values of Q. In table XXV we give the numerical outcome of this fit for the strange and non-strange hyperfine splittings. As expected from HQET the intercept always turns out to be zero within statistical errors. In the H run for $`\beta =6.2`$ we determine the hyperfine splitting from ratio-fits. In order to determine this without excited state contamination we use a fit interval for which both of the individual correlators have reached a plateau. As noted above, no significant dependence on the light quark mass was observed and we were able to fit the results to a constant with reasonable values of $`Q`$. The result is given in table XXIV. In case of the N and P run we determined the hyperfine splitting from the jackknife difference of masses obtained from the pseudo-scalar and vector meson propagator. In the N run we compared the outcome for the different smearings available, for different values of the starting point $`t_{\mathrm{min}}`$ of the fit range. An example is shown in figure 15. The different smearing functions lead to compatible answers for the hyperfine splitting. We use the outcome from the propagators with sink and source smearing $`\varphi _{\mathrm{Hg},1}`$ for our final result. In figure 16 we compare the outcome of the different runs at $`\beta =6.2`$. Clearly the outcome from the run N is the most precise. The result for the physical $`B_s`$ hyperfine splitting will be extracted from this results. Having observed clear signals for the radially excited $`S`$-wave states on our coarse lattice, we also studied their hyperfine splittings. Unfortunately the statistical noise grows rapidly and we observe no clear signal for a non-zero splitting. Our results are given in figure 14, comparing the radially excited state hyperfine splitting to that of the ground state. Although we cannot give a value for the hyperfine splitting of the radially excited $`S`$-state, our results support the expectation that it should be equal to or smaller than the ground state splitting. ### F $`P`$-state fine structure To extract the $`P`$-state fine structure we investigate the jackknife difference of the masses of the individual channels reported in tables XIX and XX. Because the statistical noise grows rapidly as $`t_{\mathrm{min}}`$ increased, this proved to be delicate. For $`\beta =5.7`$ this is illustrated in figure 17. For the matrix-fit to the $`{}_{}{}^{3}P_{2}^{}E`$ propagators we observe a jump in $`Q`$ for $`t_{\mathrm{min}}=3`$. However in the plot of the fit range dependence of the $`{}_{}{}^{3}P_{2}^{}E{}_{}{}^{3}P_{0}^{}`$ splitting the statistical uncertainty doubles between $`t_{\mathrm{min}}=3`$ and $`5`$. Therefore it is hard to tell whether there is a plateau or not. We quote final results for $`t_{\mathrm{min}}=5`$, which corresponds to dropping 2 timeslices from the first reasonable $`Q`$-value. With this procedure we obtain a large statistical error and no significant splitting can be resolved. More aggressive fitting would have led to a result compatible with zero but with a statistical error of $`30`$ MeV. We give our final numbers in table XXI. For the splitting we used the same fit range for both channels, which leads to slight deviations from the direct difference of the results in table XIX. For $`\beta =6.2`$ the situation is easier, as shown in figure 18. The noise on the splitting does not grow as fast as on the coarse lattice, because the $`PS`$ splitting is smaller in lattice units. We observe the first reasonable $`Q`$ values for $`t_{\mathrm{min}}=3`$. Since this is a 6 parameter fit to two propagators, we drop 3 time slices and quote the final result for $`t_{\mathrm{min}}=6`$. The results are given in table XXII. Here we also quote results for the splitting of the $`J=1`$ channels to the $`J=0`$ state. We reiterate that no significance should be attached to any difference in the results between the $`{}_{}{}^{3}P_{1}^{}`$ and the $`{}_{}{}^{1}P_{1}^{}`$ operators. ## VI The physical meson spectrum In this section we determine the physical $`B`$ and $`D`$-meson spectrum and investigate scaling by comparing results at different values of the lattice spacing. We also compare with experimental results and other lattice calculations. ### A $`B`$-meson spectrum At both of our lattice spacings we can simulate the $`b`$-quark directly. Here we discuss our results for the physical $`B`$-spectrum. Together with the findings of we want to investigate the dependence of the individual splittings on the lattice spacing. The findings are compared both to the existing experimental results and lattice investigations performed by other groups within a similar framework using NRQCD . The results of reference are useful in that they work at a larger lattice spacing than we do here. There are a number of problems, however that make their results not directly comparable. For example, they do not use either smeared correlators or standard fitting techniques, and this will give rise to an unknown systematic error. In addition they do not see a difference between fixing the lattice spacing from $`m_\rho `$ using the clover action and from the charmonium $`1P1S`$ splitting. This is clearly seen on finer lattices . If this arises from overestimating $`a^1`$ from $`m_\rho `$ because of discretisation errors, then this is another source of systematic error. In particular, this feeds into the fixing of the bare $`b`$ or $`c`$ quark mass and into hyperfine splittings. Their final result for the splitting does not take into account the effect of any of the uncertainties in the bare quark mass determination. This is particularly important for the hyperfine splitting and causes their errors to be heavily underestimated. The results of overlap with, but are not as complete as ours. Unfortunately there are no results for heavy clover fermions available that we can use. References quote numbers for the $`B`$ spectrum. The first reference still uses extrapolations from the lighter quark masses into the $`b`$-region. Reference determines the bare $`b`$-quark mass from heavyonium, which is not suitable for the heavy clover approach at the lattice spacings used . However we will later compare to their findings for the $`D`$-spectrum, since this problem is not so severe for charmonium at the lattice spacings used. Results from taking the $`b`$-quark as a static source also exist for spin-independent and flavour splittings, which survive in the infinite mass limit, see for example . However we restrict the discussion here to a comparison with results simulated at the physical $`b`$-quark mass directly. In the following we denote spin-averaged states by an overline. We summarise our results for the $`B`$-spectrum in the tables XXVI and XXVII. As an example of the splitting between a strange and a non-strange meson we discuss the difference of the pseudo-scalar $`B_sB`$ splitting in figure 19. We observe no scaling violations between our results and the agreement with the experimental value is excellent. The results from given in figure 19 include the statistical errors only and are taken from their figure 15. In additional uncertainties for the average of the results at the two finest lattices are mentioned in the text. The overall agreement with our results is good. They notice an upward jump, however, for their result on their finest lattice. Unfortunately our result from our finest lattice comes with large uncertainties so that we are unable to clarify whether there is any real effect here. Given the lack of scaling violations in the rest of the results, it seems unlikely to us. Table VII confirms that the scale from $`m_\rho `$ used by us and the scale from $`\sigma `$ used in are very close. The mismatch of scales of $`3\%`$ can be neglected safely. The results from are also in agreement with ours. They use the $`K^{}/K`$ ratio to fix the strange quark mass. This reduces their results compared to that using the $`K/\rho `$ ratio. Assuming a shift of $`10`$ to $`20`$ MeV from this would increase the agreement. This is the size of the effect we observe on our coarse lattice from fixing the strange from $`\varphi /\rho `$ instead of the $`K`$. Figure 20 shows the scaling of the radially excited $`B_s`$-meson. As discussed in subsection V B already, the extraction of a result for our finest value of $`a`$ turned out to be more problematic than anticipated, and we are left with quite large statistical uncertainties. However the final result is in good agreement with the result from our coarser lattice as well as the result of . We also included a preliminary experimental result from the DELPHI collaboration for an admixture of the non-strange $`B^{}B`$ and $`B^{}{}_{}{}^{}B^{}`$ splitting . Assuming the hyperfine splitting of the two states to be of similar size, which our findings support, we observe reasonable agreement here. Table XXVI also contains results for the radial excitation energy of the vector state and the spin-averaged $`S`$-wave. The orbital excitations are compared in figure 21. The lattice results for the spin-averaged strange $`P`$ state scale very well. The magnitude agrees nicely with the $`B_{sJ}^{}(5850)`$ resonance, which is expected to be an admixture of the two $`j_l=\frac{3}{2}`$ states. Our results for the radially excited $`P`$-states are compared in figure 22. This is the first observation of a signal for these states in a lattice calculation. As the figure shows we get consistent results from the two different lattice spacings investigated. To the best of our knowledge radially excited $`P`$-states have not been observed yet experimentally. The splittings discussed above are all essentially light quark quantities, which survive into the static limit. Their scaling or non-scaling says more about the light quark action than the heavy quark sector. The hyperfine splitting is of a different nature and from its scaling behaviour one can learn about how well the heavy quarks are being described on the lattice. We display this in figure 23. Our result for the strange hyperfine splitting together with the findings from shows good scaling. However the result is much smaller than the experimental value. Since the leading term in the hyperfine splitting arises from the $`𝝈𝑩`$ term in the action, eq. (8), the result for the splitting is sensitive to the coefficient $`c_4`$ and the inclusion of radiative corrections beyond tadpole improvement is required. Preliminary calculations indicate that the inclusion of the 1-loop corrections would increase the hyperfine splittings on the order of $`10\%`$ for the lattice spacings used. The quenched approximation might also play a rôle here, since in light spectroscopy the hyperfine splittings turn out to be too small as well, see for a review. This effect increases with increasing quark mass. Unfortunately reference , which investigates the effect of the inclusion of two flavours of dynamical quarks on the $`B`$-meson, does not give any evidence for an increase of the $`B^{}B`$ splitting due to sea-quark vacuum polarisation effects. From the experience in $`\mathrm{{\rm Y}}`$-spectroscopy using NRQCD, one could have expected to observe scaling violations in the hyperfine splitting. Using the same heavy quark Hamiltonian as we do, reports an increase of $`50\%`$ for the $`\mathrm{{\rm Y}}\eta _b`$ splitting, within the range from $`\beta =5.7`$ up to $`\beta =6.2`$. The leading discretisation correction for the hyperfine splitting is $`𝒪((ap_{\mathrm{gluon}})^2)`$ . Typical gluon momenta for the $`\mathrm{{\rm Y}}`$-system are $`1`$ GeV, while for the $`B`$-system they are $`𝒪(\mathrm{\Lambda }_{\mathrm{QCD}})`$. From this one expects reduced scaling violations of $`10\%`$ in the $`B`$-system for our range of lattice spacings. This is the same size as our uncertainties on the hyperfine splitting and therefore consistent with the fact that no scaling violations show up in figure 23. Results from are for the chirally extrapolated splitting $`B^{}B`$ with statistical errors only, taken from their figure 17. In the text, the authors quote a result for the strange hyperfine splitting from the average of the two finest lattices, which is $`3`$ MeV higher than the same average for the non-strange hyperfine splitting. An upwards shift of $`3`$ MeV increases the already excellent agreement even further. The results of , on the other hand, exhibit a clear disagreement to our findings as well as the findings of . We believe that this is because they have determined the bare $`b`$-quark mass $`am_b`$ using heavyonium. The fine structure of the $`P`$-states is the last topic of this section. Unfortunately we have not been able to resolve this clearly on our coarsest lattice. The situation is displayed in figure 24, for the three sublevels which were resolved at $`\beta =6.0`$ and $`6.2`$. To investigate whether there is evidence for scaling violations in the fine structure we calculate the jackknifed difference of the highest and lowest state. This is shown in figure 25. The error bars turn out to be large and the figure is inconclusive. A more aggressive fit on the coarser lattice, as discussed in subsection V E would lead to the conclusion that scaling violations were seen, but we believe that further work is needed to resolve this question. Our results for the $`B`$-meson spectrum, together with those of do not show signs of residual lattice spacing dependence within the achieved accuracy. Therefore we can average the results for the different values of the lattice spacing $`a`$ to obtain our final results on the quenched $`B`$-meson excitation spectrum. The averages were obtained in the following way, for each value of $`a`$ we add the different uncertainties in quadrature to obtain a single value. Here we omitted those sources of uncertainty which are associated with the quenched approximation. These are the uncertainty arising from fixing the strange quark mass from different physical quantities and in the case of the results of the additional uncertainty of the lattice spacing $`a`$ associated with the physical quantity used to fix $`a`$. At this step we also symmetrised with respect to unsymmetric uncertainties. The central values have been obtained from fitting the results to a constant in an uncorrelated fit with the above described uncertainties. This puts more weight on the more precise results than a simple average. Our analysis at the individual values of $`a`$ does not include an uncertainty for the residual effect of the lattice spacing. Our actions are improved to $`𝒪(\alpha _sa,a^2)`$. Therefore for each value of $`a`$ we add the maximum of $`\alpha _sa\mathrm{\Lambda }_{\mathrm{QCD}}`$ and $`a^2\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ in quadrature to the uncertainty used in the fit. We quote the smallest of these three so obtained uncertainties as our final uncertainty for the quenched $`B`$-meson spectrum. This way we quote an accuracy which is of the same size as the one we checked for scaling violations. Determining the final uncertainty from the $`\chi ^2`$ of the fit would reduce the uncertainty beyond this level. This procedure also ensures that residual lattice spacing artifacts are properly included if the achieved accuracy differs over the three individual results and the average is largely determined by the coarser lattices. Our final result on the $`B`$-meson splitting spectrum in the quenched approximation is given in table XXVIII and figure 26. The question of the effect of quenching on the spectrum goes beyond the scope of this paper. We refer the reader to . There the effects of the inclusion of 2 flavours of dynamical quarks on the spectrum in NRQCD have been investigated and compared to the findings of . No significant difference between the quenched and $`n_f=2`$ results for the $`1S`$, $`2S`$ and $`1P`$-states was found. In particular, as mentioned earlier, no significant sea quark effects were seen in the $`1S`$-hyperfine splitting. Since our investigations confirm scaling in the quenched heavy-light spectrum the conclusions of are unchanged. ### B $`D`$-meson spectrum In this section we discuss the $`D`$-meson spectrum and compare our result to existing lattice results as well as to experiment. The convergence of the NRQCD expansion is particularly important in the $`D`$-range, where the expansion parameter $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q\frac{1}{4}`$. Useful results on the question of the convergence are contained in . There the authors study the complete NRQCD action to $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^3)`$. Here we include all terms up to $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^2)`$ and the relativistic correction to the kinetic energy in $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^3)`$. The authors of calculate the heavy light kinetic masses using eq. (24) and show that the difference arising from $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^3)`$-terms is consistent with the expectation that they are sub-sub-leading in a $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q`$ expansion. The changes to the spin-averaged meson mass that we use to fix the quark mass are dominated by the $`\frac{p^4}{m_Q^3}`$ relativistic correction that we include. From this we conclude that remaining $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^3)`$ and higher order terms in the NRQCD expansion would only change the physical masses by at most a few percent. This allows us to use the results of to estimate the changes in the hyperfine splitting which would be produced by these additional terms at fixed bare quark mass. The authors of find that the $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^2)`$ terms produce an effect somewhat smaller than a $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q`$ expansion might suggest, since they affect the hyperfine splitting indirectly. The only spin-dependent term at $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^2)`$ is a spin-orbit type interaction. At $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^3)`$ most terms produce a change of a few percent, but the one which is directly related to the $`𝝈𝑩`$-term: $`\{𝐃^2,𝝈𝑩\}`$ reduces the hyperfine splitting at the charm by $`20\%`$. However when including the other operators of $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^3)`$, the second largest effect comes from the $`𝝈(𝑬\times 𝑬+𝑩\times 𝑩)`$ operator, which is also spin-dependent and works in the opposite direction to the other one. The total effect of $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^3)`$ is below $`10\%`$. This is of the size of the naïve expectation for the suppression with respect to the leading term and not at all inconsistent with good convergence of the NRQCD expansion. Since we do not include these terms, we conclude, that we may be overestimating the quenched lattice hyperfine splitting of the $`D`$-meson by $`10\%`$. Our results on the $`D`$-meson spectrum are summarised in table XXIX and in figure 27. The overall agreement to the experimentally observed spectrum is good. We will now discuss the individual splittings in more detail. We also compare our results to the lattice studies of . All of these results use the quenched approximation as well. The publications apply the heavy clover approach , which has quite different systematic uncertainties from NRQCD for charm quarks. The flavour dependent $`D_s^{()}D^{()}`$ splittings are in good agreement with the experimental results. Here it is interesting to note that our results reflect the increase of $`10`$ MeV from the $`B`$ to the $`D`$-meson system, which can already be expected from the good agreement of the slope with the experimental outcome in table XXV. In figure 28 we compare our result for the strange to non-strange spin-averaged splitting to other lattice calculations. In order not to disguise other possible effects, we excluded the uncertainty of the strange quark mass from the plot. The results from use the $`K^{}/K`$ ratio to define the strange quark. This should shift the results downwards, compared to fixing $`\kappa _s`$ from the $`K/\rho `$ ratio as used for the other results. The implications have already been discussed in the previous sub-section VI A. The uncertainties again allow for an upwards shift of these results by $`10`$ to $`20`$ MeV. It should be noted that we combined the results for the hyperfine splittings and the pseudo-scalar $`D_sD`$ splitting from to obtain the spin-averaged splitting. Because of different systematic uncertainties, it is particularly interesting to compare to the heavy clover results of . The results obtained with the use of the $`m_\rho `$-scale, which is the same as what we use, agree very well with ours. This agreement is expected since uses the same light quark and gauge field action and this quantity is essentially determined by the light quarks and the gluon field. We conclude that the results for the flavour dependent $`D_sD`$ splitting agree well between the different approaches and, within the accuracy achieved, agree well with the experimental result. For radially excited $`D_s^{()}^{}`$-mesons no experimental results are known to us. However the DELPHI collaboration reports on the observation of the non-strange $`D^{}^{}`$ . This result is still awaiting confirmation by the OPAL and the CLEO collaboration and its interpretation is disputed on the ground of its small experimentally observed width . The splitting between the DELPHI result and the $`D^{}`$ has a similar size to our $`D_s^{}{}_{}{}^{}D_s^{}`$ splitting. Reference reports lattice results from the heavy-clover approach. From their plot we read $`D_s^{}D_s840(160)`$ MeV, which is in agreement with our findings. However this includes what the authors call “continuum” extrapolation out of a regime where the expansion parameter $`am_Q=𝒪(1)`$ is not small. We would prefer to compare to the unextrapolated results at the individual values of $`a`$. Experimentally the only well established charmed $`P`$-states in the particle data book are those which are expected to correspond to the states of total light angular momentum $`j_l=\frac{3}{2}`$. Recently the CLEO collaboration claimed the observation of the $`D_1`$ state corresponding to $`j_l=\frac{1}{2}`$. CLEO gives a preliminary result of $`D_1=2461(_{34}^{+41})(10)(32)`$ MeV, which is slightly heavier but compatible in error bars to the $`D_1^{}=2425`$ MeV <sup></sup><sup></sup>*We quote the charge-average.. Our lattice calculation delivers the mass of the lighter of the two states. We did not observe a signal for an excited state slightly heavier than this. In table XXIX our result for the $`\overline{D}_s(1P)\overline{D}_s`$ splitting is compared to the spin-average of the $`D_{s1}^{}`$ and $`D_{s2}^{}`$, the $`j_l=\frac{3}{2}`$ states. The agreement is reasonable. Reference reports on the $`D_{s1}\overline{D}_s`$ splitting from a lattice study with the heavy clover approach. A comparison to our result for this splitting is given in figure 29. When using the same scale obtained from $`m_\rho `$ both lattice results agree very well with each other. The agreement with experiment is also good. It should be noted however, that the experimental result included in figure 29 is not necessarily the same as ours. The experimental $`P`$-state corresponds to $`j_l=\frac{3}{2}`$. If the CLEO trend is confirmed and the $`D_1`$ is indeed heavier than the $`D_1^{}`$ and the same holds for the $`D_{s1}`$ states, then the lattice result also corresponds to $`j_l=\frac{3}{2}`$. If not the lattice result will correspond to $`j_l=\frac{1}{2}`$, but the two states will be so close, that any mismatch is well covered by the error bars. Table XXIX contains our final result for the radially excited $`2P`$-state. This is the first result for this state from a lattice simulation. As in the $`B`$-system the hyperfine splittings are too small when compared to the experimental result. Whereas in the $`B`$-system they were too small by $`40\%`$, here they are low by $`25\%`$. This could reflect a more severe quenching error for $`B`$-mesons. $`B`$-mesons are somewhat smaller states than $`D`$ mesons and probe slightly different scales. This is a sub-leading effect in a heavy quark symmetry picture, however. Alternatively, if the error comes from radiative corrections to the $`c_4`$ coefficient, that would need to increase with $`m_Q`$. That is seen by the authors of . The large uncertainty of $`20`$ MeV on our result arises from the chiral extrapolations used in the lattice spacing determination and the way in which this feeds into the fixing of the bare quark mass. Naively we expect a doubling of the relative error, because a larger value of $`a`$ requires a smaller value of $`am_Q`$ to deliver the same physical $`m_{\mathrm{sav}}`$. This smaller value of $`am_Q`$ gives a larger hyperfine splitting $`am_{hpf}`$. When converting to physical units it picks up the uncertainty of $`a`$ for the second time. In fact a factor of four is seen because of the flattening of the relation between the mass shift $`\mathrm{\Delta }`$ and bare heavy quark mass $`m_Q`$, see figure 3, as well as the steepening up of the hyperfine splitting curve for large values of $`m_{\mathrm{sav}}`$ in figure 13. In figure 30 we compare our results to the results from obtained in NRQCD. We choose their result in $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^2)`$ as most relevant for this comparison. The good agreement with our results is in fact misleading. They fix their $`c`$-quark mass from charmonium instead of the $`D`$. Fixing from the $`D`$ would lead to a larger value of $`am_c`$ and lower hyperfine splitting, see subsection IV C and reference . It is interesting to compare the result for the hyperfine splitting between NRQCD and heavy clover quarks. This is also done in figure 30 for the $`D_s`$ hyperfine splitting. The heavy clover results of appear to be higher than our result. However this is a result of their higher choice of scale coming from $`J/\psi `$ instead of from $`m_\rho `$. This is confirmed by the findings of . Using a scale from $`m_\rho `$ gives a result which agrees with ours, using a scale from $`J/\psi `$ agrees with reference . It should be noted that in the bare quark mass is determined from the $`D`$, where as in it is determined from charmonium. For the heavy clover approach at $`\beta =6.0`$ these differences are negligible within statistical errors . As discussed at the beginning of this subsection, the inclusion of the terms $`𝒪((\mathrm{\Lambda }_{\mathrm{QCD}}/m_Q)^3)`$ contributing to the hyperfine splitting would decrease our result by $`10\%`$. However the heavy clover approach requires similar correction terms in the Hamiltonian to achieve this level of accuracy. The agreement of the NRQCD and the heavy clover result for the spin-dependent hyperfine splitting is encouraging, since the systematic uncertainties are quite different. The light quark content plays only a minor rôle for the hyperfine splitting, which depends essentially only on the heavy quark Hamiltonian. In NRQCD the leading contribution to the hyperfine splitting comes from the $`𝝈𝑩`$ term in the action, whereas for heavy clover this is split between the kinetic hopping term and the clover term $`\sigma _{\nu \rho }F_{\nu \rho }`$, with the latter becoming more important as the lattice spacing becomes coarser. Both these actions give rise to systematic errors in the hyperfine splitting from mass-dependent radiative corrections to coefficients and neglected higher order terms, each at the $`10\%`$ level, so the differences could have been significantly larger than observed. ## VII Discussion We present an extensive study of the $`B`$ and $`D`$-meson spectrum using NRQCD heavy quarks and clover light quarks in the quenched approximation. Our results include the splitting between the strange and the non-strange meson, hyperfine splittings, radially and orbitally excited states. For the first time in a lattice calculation we obtained a result on radially excited $`P`$-wave states. For spin-independent splittings we observe good agreement with experimental results. However, our result for the spin-dependent hyperfine splitting turns out to be too low in comparison to experiment. This is a well known effect in quenched hadron spectroscopy. Furthermore, in the present calculation hyperfine splittings are also affected by the neglect of radiative corrections in the matching of lattice NRQCD to continuum QCD. Using two different values of the lattice spacing in the $`B`$-spectrum together with the results of allows for a detailed investigation of the residual lattice spacing dependence of our final results. No scaling violations are observed within the achieved accuracy. Of particular interest is the scaling of the $`B_s^{}B_s`$ splitting, which depends heavily on the properties of the heavy quark content of the theory. Here scaling violations could be ruled out with an accuracy of $`10\%`$. The $`P`$-fine structure has not been resolved for all values of the lattice spacings and further work is needed for this quantity. Our results on the $`B`$-meson spectrum are summarised in table XXVIII and figure 26 together with the findings of . In addition to the uncertainties considered in the analysis at the individual values of $`a`$, the quoted uncertainties also contain an estimate for the residual lattice spacing artifacts of $`𝒪(\alpha _sa,a^2)`$. The table gives our final results for the $`B`$-meson spectrum in the quenched approximation. Our final results on the $`D`$-meson spectrum are shown in table XXIX and figure 27 above. This is our final result for a lattice spacing of $`a^11.1`$ GeV and does not include an estimate of the residual lattice spacing artifacts of $`𝒪(\alpha _sa,a^2)`$. For the above value of $`a`$, this corresponds to $`13\%`$, which is of similar size to or smaller than the otherwise achieved accuracy. We compared our results to lattice results of other collaborations obtained with NRQCD or in the heavy clover framework. In general we observe good agreement. Discrepancies which appear at first sight could be traced to underestimated errors in these other results or the use of different scales when converting the lattice results into physical units. The excellent agreement of our results with the results obtained in the heavy clover approach is noteworthy because of the different systematics of these approaches. These results are the most complete lattice results on the $`B`$ and $`D`$-meson spectrum to date. ### Acknowledgements We would like to acknowledge useful discussions with Peter Boyle and Peter Lepage. J.H. was supported by a Marie Curie research fellowship by the European commission under ERB FMB ICT 961729, by PPARC and the National Science Foundation. S.C. acknowledges fellowships from the Royal Society of Edinburgh and the Alexander von Humboldt Stiftung. A.A.K. was supported by the Research for the Future Program of the Japanese Society for the Promotion of Science. J.Sh. would like to thank members of the theoretical physics group at the University of Glasgow for their hospitality during an extended visit. Support from an UK PPARC Visiting Fellowship PPA/V/S/1997/00666, is gratefully acknowledged. This work was supported by the DOE under DE-FG02-91ER40690, PPARC under GR/L56343 and NATO under CRG/94259. The gauge configurations at both values of $`\beta `$ and the light quark propagators at $`\beta =6.2`$ have been generously provided by the UKQCD-collaboration. The simulations at $`\beta =5.7`$ have been performed at NERSC supported by DOE and the ones at $`\beta =6.2`$ at EPCC in Edinburgh supported by PPARC.
warning/0003/hep-ph0003038.html	ar5iv	text	# 1 INTRODUCTION ## 1 INTRODUCTION Super-Kamiokande collaboration has confirmed a neutrino oscillation by their atmospheric neutrino experiments . In their two-flavor mixing analyses of the sub-GeV and the multi-GeV zenith angle distribution, it has been obtained that the $`\nu _\mu \nu _\tau `$ oscillation is preferred to the $`\nu _\mu \nu _e`$ oscillation and the range of mass parameter $`\mathrm{\Delta }m^2`$ is from $`10^3\mathrm{eV}^2`$ to $`10^2\mathrm{eV}^2`$. Recently, Super-Kamiokande collaboration has reported that the $`\nu _\mu \nu _\tau `$ oscillation is maximally $`\mathrm{sin}^22\theta =0.9\text{}1`$ and mass parameter $`\mathrm{\Delta }m^2`$ is from $`2.5\times 10^3\mathrm{eV}^2`$ to $`5\times 10^3\mathrm{eV}^2`$, using the sub-GeV, multi-GeV neutrino and upward muons zenith angle distribution experiments (830–920 live days) . However, these results are obtained from the two-flavor neutrino analyses. In order to account for the solar neutrino anomaly data together with atmospheric neutrino experiments, three flavor neutrinos are necessary at least. In three-flavor neutrinos scenario with a mass hierarchy $`m_1m_2m_3`$, there are necessary two mass parameters $`\mathrm{\Delta }m_{12}^2`$ and $`\mathrm{\Delta }m_{23}^2`$, and three mixing angles $`\theta _{12}`$, $`\theta _{13}`$ and $`\theta _{23}`$. Solar neutrino anomaly gives constraint on only three parameters $`\mathrm{\Delta }m_{12}^2`$, $`\theta _{12}`$ and $`\theta _{13}`$. The MSW solution for solar neutrinos predicts the large mixing angle solution ($`\mathrm{\Delta }m_{12}^2=4\times 10^6`$ $`7\times 10^5\mathrm{eV}^2`$, $`\mathrm{sin}^22\theta _{12}=0.6\text{}0.9`$) and the small mixing angle solution ($`\mathrm{\Delta }m_{12}^2=3\times 10^6`$$`1.2\times 10^5\mathrm{eV}^2`$, $`\mathrm{sin}^22\theta _{12}=0.003\text{}0.01`$) for $`\theta _{13}=0^{}\text{}20^{}`$, and these large and small mixing angle solutions are merged for $`25^{}\text{}50^{}`$ . The vacuum solution is also obtained as $`\mathrm{\Delta }m_{12}^210^{10}\mathrm{eV}^2`$. However, CHOOZ experiment which is a terrestrial experiment using reactor neutrinos gives a strong constraint $`\mathrm{sin}^22\theta _{13}<0.18`$ for large $`\mathrm{\Delta }m_{23}^2`$ and $`\mathrm{\Delta }m_{23}^2<0.9\times 10^3\mathrm{eV}^2`$ for $`\mathrm{sin}^22\theta _{13}1`$. From the recent many atmospheric neutrino analyses using either two- or three-flavor neutrinos framework, $`\mathrm{\Delta }m_{23}^2=10^3\text{}10^2\mathrm{eV}^2`$ is obtained. Therefore, it can be said that the mixing angle $`\theta _{13}`$ is small but not necessarily 0. In the three-flavor neutrino framework, if the mixing parameter $`\theta _{13}`$ is not zero, it is necessary to consider the interplay between two mass parameters $`\mathrm{\Delta }m_{12}^2`$, $`\mathrm{\Delta }m_{23}^2`$ and three mixing angles $`\theta _{12}`$, $`\theta _{13}`$, $`\theta _{23}`$ for atmospheric neutrino experimental analyses. The oscillation term $`\mathrm{sin}^21.27{\displaystyle \frac{\mathrm{\Delta }m_{12}^2}{E}}L`$ cannot be neglected though $`\mathrm{\Delta }m_{12}^2\mathrm{\Delta }m_{23}^2`$, because the term is not so small in the sub-GeV experiment ($`E=0.21.3\mathrm{GeV}`$) of atmospheric neutrino with zenith angle $`\theta 180^{}`$ at which $`L10000\mathrm{k}\mathrm{m}`$. In the multi-GeV experiment ($`E>1.3\mathrm{GeV}`$) of atmospheric neutrino and terrestrial short- and long-baseline experiments, the neutrino oscillation term $`\mathrm{sin}^21.27{\displaystyle \frac{\mathrm{\Delta }m_{12}^2}{E}}L`$ can be neglected. The angle $`\theta _{13}`$ has been seen to be small, then the $`\theta _{13}`$ is approximated to be $`0`$ in usual analyses. However, in the case considering the matter effects of the Earth, above approximation is not appropriate. The matter effect is expressed by the induced mass squared $`A=2\sqrt{2}EG_FN_e=7.59\rho E\times 10^5\mathrm{eV}^2`$, where $`N_e`$ is the number density of electrons, $`E`$ is the energy of neutrinos (measured in GeV) and $`\rho `$ is the matter density (measured in gm/cm<sup>3</sup>). The Earth density $`\rho `$ is $`3.513`$g/cm<sup>3</sup> and the multi-GeV experiment energy of neutrinos ranges from 1.3GeV to 100GeV. Then the value of $`A`$ go through the value of mass parameter $`\mathrm{\Delta }m_{23}^2`$ and the MSW effect occurs in the mixing angle $`\theta _{13}`$ (see section 3). In this paper, we analyze the Super-Kamiokande atmospheric neutrino experimental data(850–920 live days) of sub-GeV and multi-GeV neutrino, upward through-going and stopping muons, using the three-flavor neutrino framework with the hierarchy $`m_1m_2m_3`$. In previous paper , we analyzed the 535 live days data of Super-Kamiokande atmospheric neutrino experiment using similar framework to the present analysis, approximating the Earth density to be constant. In this present analysis, we will pursue a full calculation adopting the varying density of the Earth. We will also discuss the long-baseline K2K experiment in the three-flavor neutrino framework. ## 2 NEUTRINO OSCILLATION IN THREE-FLAVOR NEUTRINOS The unitary matrix $`U`$ which transforms the mass eigenstate neutrinos $`\nu _\alpha `$ to the flavor eigenstate neutrinos $`\nu _l`$ as the formula $$\nu _l=\underset{\alpha =1}{\overset{3}{}}U_{l\alpha }\nu _\alpha ,l=e,\mu ,\tau ,$$ (1) is parametrized as follows: $`U`$ $`=`$ $`\mathrm{exp}(i\theta _{23}\lambda _7)\mathrm{exp}(i\theta _{13}\lambda _5)\mathrm{exp}(i\theta _{12}\lambda _2)`$ $`=`$ $`\left(\begin{array}{ccc}c_{12}c_{13}& s_{12}c_{13}& s_{13}\\ s_{12}c_{23}c_{12}s_{23}s_{13}& c_{12}c_{23}s_{12}s_{23}s_{13}& s_{23}c_{13}\\ s_{12}s_{23}c_{12}c_{23}s_{13}& c_{12}s_{23}s_{12}c_{23}s_{13}& c_{23}c_{13}\end{array}\right),`$ $`c_{ij}=\mathrm{cos}\theta _{ij},s_{ij}=\mathrm{sin}\theta _{ij}.`$ in a case disregarding the CP violation. We assume the mass hierarchy $$m_1m_2m_3,$$ (3) then $`\mathrm{\Delta }m_{12}^2\mathrm{\Delta }m_{13}^2\mathrm{\Delta }m_{23}^2`$. In this mass hierarchy, the transition probabilities $`P(\nu _l\nu _l^{})`$ can be written as $`P(\nu _l\nu _l)`$ $`=`$ $`12(12U_{l3}^2U_{l1}^4U_{l2}^4+U_{l3}^4)S_{12}4U_{l3}^2(1U_{l3}^2)S_{23},`$ (4) $`P(\nu _l\nu _l^{})`$ $`=`$ $`P(\nu _l^{}\nu _l)=2(U_{l1}^2U_{l^{}1}^2+U_{l2}^2U_{l^{}2}^2U_{l3}^2U_{l^{}3}^2)S_{12}+4U_{l3}^2U_{l^{}3}^2S_{23},`$ (5) where $`S_{\alpha \beta }`$ is a term representing the neutrino oscillation defined as; $$S_{\alpha \beta }=\mathrm{sin}^21.27\frac{\mathrm{\Delta }m_{\alpha \beta }^2}{E}L.$$ (6) Here $`\mathrm{\Delta }m_{\alpha \beta }^2=\|m_\alpha ^2m_\beta ^2\|`$, $`E`$ and $`L`$ are measured in units eV<sup>2</sup>, GeV and km, respectively. For the mass parameter $`\mathrm{\Delta }m_{12}^2`$ and mixing angle $`\theta _{12}`$, we use values obtained in the solar neutrino experimental analyses: the large mixing angle solution $`\mathrm{\Delta }m_{12}^2=3\times 10^5\mathrm{eV}^2`$, $`\mathrm{sin}^22\theta _{12}=0.7`$, and small mixing angle solution $`\mathrm{\Delta }m_{12}^2=10^5\mathrm{eV}^2`$, $`\mathrm{sin}^22\theta _{12}=0.005`$. It should be noted that the oscillation term $`S_{12}=\mathrm{sin}^21.27{\displaystyle \frac{\mathrm{\Delta }m_{12}^2}{E}}L`$ cannot be neglected in the sub-GeV experiment of atmospheric neutrino though $`\mathrm{\Delta }m_{12}^2\mathrm{\Delta }m_{23}^2`$, because this oscillation term is not so small in sub-GeV neutrino energy ($`E=0.2\text{}1.3\mathrm{GeV}`$) with zenith angle $`\theta 180^{}`$($`L10000\mathrm{k}\mathrm{m}`$). In the atmospheric neutrino experiments, the matters of the Earth have a important effect. Matter effect is represented by a term $`A=2\sqrt{2}EG_FN_e=7.59\rho E\times 10^5\mathrm{eV}^2`$ induced by the Earth matter, where $`\rho `$ is the matter density. In the Earth, $`\rho =3.5\text{}13\mathrm{g}/\mathrm{cm}^3`$, and sub-GeV and multi-GeV neutrino energy ranges from 0.1GeV to 100GeV, then the value of $`A`$ goes through those of $`\mathrm{\Delta }m_{12}^2`$ and $`\mathrm{\Delta }m_{23}^2`$ and the resonance happens in mixing angles $`\theta _{12}`$ and $`\theta _{13}`$ as will be seen in Eqs.(8)–(11) and Eqs.(13)–(16). In the sub-GeV neutrino energy, $`A\mathrm{sin}2\theta _{13}/2\mathrm{\Delta }m_{23}^21`$, then we can approximate the mixing matrix $`U`$ as $$U_{\mathrm{sub}}^M=\mathrm{exp}(i\lambda _7\theta _{23})\mathrm{exp}(i\lambda _5\theta _{13})\mathrm{exp}(i\lambda _2\theta _{12}^M),$$ (7) where $`\mathrm{sin}2\theta _{12}^M={\displaystyle \frac{\mathrm{\Delta }m_{12}^2}{\mathrm{\Delta }m_{12}^{M2}}}\mathrm{sin}2\theta _{12},`$ (8) $`\mathrm{\Delta }m_{12}^{M2}=m_2^{M2}m_1^{M2},\mathrm{\Sigma }=m_1^2+m_2^2,`$ (9) $`m_{1,2}^{M2}={\displaystyle \frac{1}{2}}\{(\mathrm{\Sigma }+A\mathrm{cos}^2\theta _{13})`$ $`\sqrt{(A\mathrm{cos}^2\theta _{13}\mathrm{\Delta }m_{12}^2\mathrm{cos}2\theta _{12})^2+(\mathrm{\Delta }m_{12}^2\mathrm{sin}2\theta _{12})^2}\},`$ (10) $`m_3^{M2}=m_3^2+A\mathrm{sin}^2\theta _{13}.`$ (11) On the other hand, $`A\mathrm{sin}2\theta _{13}/2\mathrm{\Delta }m_{23}^2`$ is not so small in multi-GeV neutrino and upward muon energy. In this case, the mixing matrix can be expanded by powers of a quantity $`\mathrm{\Delta }m_{12}^2\mathrm{sin}2\theta _{12}/2m_3^2`$ which is very small. In the lowest order, the mixing matrix is expressed as $$U_{\mathrm{multi}}^M=\mathrm{exp}(i\lambda _7\theta _{23})\mathrm{exp}(i\lambda _5\theta _{13}^M),$$ (12) where $`\mathrm{sin}2\theta _{13}^M={\displaystyle \frac{\mathrm{\Delta }m_{13}^2}{\mathrm{\Delta }m_{13}^{M2}}}\mathrm{sin}2\theta _{13},`$ (13) $`\mathrm{\Delta }m_{13}^{M2}=m_3^{M2}m_1^{M2},\mathrm{\Lambda }=\mathrm{\Sigma }\mathrm{\Delta }m_{12}^2\mathrm{cos}2\theta _{12},\mathrm{\Sigma }=m_1^2+m_2^2,`$ (14) $`m_{1,3}^{M2}={\displaystyle \frac{1}{2}}\{(m_3^2+{\displaystyle \frac{\mathrm{\Lambda }}{2}}+A)`$ $`\sqrt{\left(A\mathrm{\Delta }m_{13}^2\mathrm{cos}2\theta _{13}\right)^2+\left(\mathrm{\Delta }m_{13}^2\mathrm{sin}2\theta _{13}\right)^2}\},`$ (15) $`m_2^{M2}={\displaystyle \frac{1}{2}}(\mathrm{\Sigma }+\mathrm{\Delta }m_{12}^2\mathrm{cos}2\theta _{12}).`$ (16) These expressions correspond to the neutrino case, and $`A`$ for the anti-neutrino case has an opposite sign of the neutrino case. From these expressions, we can recognize that a resonance exists at $`A\mathrm{\Delta }m_{13}^2\mathrm{cos}2\theta _{13}`$ of $`\mathrm{sin}2\theta _{13}^M`$. When $`\mathrm{\Delta }m_{13}^2=3\times 10^3\mathrm{eV}^2`$ and $`\mathrm{cos}2\theta _{13}1`$, the resonance occurs at $`E10\mathrm{GeV}`$. When the density of the Earth matter is constant, the transition probability $`P^M(\nu _l\nu _l^{})`$ with the matter effects can be expressed by the expression (5) with $`U`$ and $`\mathrm{\Delta }m_{\alpha \beta }^2`$ replaced by $`U^M`$ and $`\mathrm{\Delta }m_{\alpha \beta }^{M2}`$. Actual density of the Earth is not constant but the density in shells which compose the Earth is almost constant. The net transition probability $`P_{\nu _l\nu _l^{}}`$ of neutrinos going through the Earth can be given by connecting the transition amplitudes of each shell. The transition amplitude in $`k`$-th shell is expressed as $$T_{l^{}l}(M_k,x_k)=\underset{\alpha }{}(U^{M_k})_{l^{}\alpha }(U^{M_k})_{l\alpha }\mathrm{exp}\left(i\frac{m_\alpha ^{M_k2}}{2p}x_k\right),$$ (17) where $`U^{M_k}`$ represents the mixing matrix containing the $`k`$-th shell matter effects, $`m_\alpha ^{M_k}`$ the mass of $`\nu _\alpha `$ in the $`k`$-th shell and $`x`$ the length traveling in the $`k`$-th shell. Using this transition amplitude, we get the transition probability as follows: $`P_{\nu _l\nu _l^{}}`$ $`=`$ $`\|{\displaystyle \underset{l_1,l_2,\mathrm{}}{}}T_{l^{}l_{n1}}(M_n,x_n)T_{l_{n1}l_{n2}}(M_{n1},x_{n1})\mathrm{}`$ (18) $`\mathrm{}T_{l_2l_1}(M_2,x_2)T_{l_1l}(M_1,x_1)\|^2.`$ Calculation using this expression is very complicated and cannot be carry out analytically, then we calculate this numerically. ## 3 NUMERICAL ANALYSES OF ATMOSPHERIC NEUTRINOS In this work, we analyze the ratio $`N_{\mathrm{Exp}}(l)/N_{\mathrm{MC}}(l)`$ of experimentally observed events $`N_{\mathrm{Exp}}(l)`$ and expected events $`N_{\mathrm{MC}}(l)`$ without oscillation, where $`l`$ represents $`\mu `$-like and $`e`$-like events. The zenith angle $`\theta `$ dependent events $`dN_{\mathrm{Exp}}(l)/d\mathrm{cos}\theta `$ and $`dN_{\mathrm{MC}}(l)/d\mathrm{cos}\theta `$ are defined as $$\frac{dN_{\mathrm{Exp}}(l)}{d\mathrm{cos}\theta }=\underset{\nu _l^{}}{}ϵ_l(E_l)\sigma _{\nu _l}(E_{\nu _l^{}},E_l,\psi )F_{\nu _l^{}}(E_{\nu _l^{}},\theta \psi )P^M(\nu _l^{}\nu _l)𝑑E_{\nu _l^{}}𝑑E_ld\mathrm{cos}\psi ,$$ (19) $$\frac{dN_{\mathrm{MC}}(l)}{d\mathrm{cos}\theta }=ϵ_l(E_l)\sigma _{\nu _l}(E_{\nu _l},E_l,\psi )F_{\nu _l}(E_{\nu _l},\theta \psi )𝑑E_{\nu _l}𝑑E_ld\mathrm{cos}\psi ,$$ (20) where the summation $`_{\nu _l^{}}`$ are taken over $`\nu _\mu `$ and $`\nu _e`$. In these expressions, processes of $`\overline{\nu }_\mu `$ and $`\overline{\nu }_e`$ are contained. $`ϵ_l(E_l)`$ is the detection efficiency of the detector for $`l`$-type charged lepton with energy $`E_l`$. $`\sigma _{\nu _l}(E_{\nu _l^{}},E_l,\psi )`$ is the differential cross section of scattering $`l`$ with energy $`E_l`$ by incident $`\nu _l`$ with energy $`E_{\nu _l^{}}=E_{\nu _l}`$, where angle $`\psi `$ is the scattering angle between the directions of incident $`\nu _l`$ and scattered $`l`$. $`F_{\nu _l^{}}(E_{\nu _l^{}},\theta )`$ is the incident $`\nu _l^{}`$ flux with energy $`E_{\nu _l^{}}`$ produced at the atmosphere coming to the detector with zenith angle $`\theta `$. $`E_{\nu _l}`$ and $`\theta `$ dependences of $`F_{\nu _l}(E_{\nu _l},\theta )`$ for multi-GeV experiment ($`E_{\nu _l}>1.33\mathrm{GeV}`$) are given in Refs. . These dependences including the geomagnetic effects for sub-GeV case ($`0.2\mathrm{GeV}<E_{\nu _l}<1.33\mathrm{GeV}`$) are taken from Ref. . Other informations of $`ϵ_l(E_l)`$ and $`\sigma _{\nu _l}(E_{\nu _l},E_l,\psi )`$ are given in Ref. . The upward through-going muons (thru-muons) and stopping muons (stop-muons) fluxes are shown in Ref. . Typical energies of $`\nu _\mu `$ that produce thru- and stop-muons are 100GeV and 10GeV, respectively. Explicit calculation of Eqs. (19) and (20) is explained precisely in Appendix A of the second paper in Ref. . $`P^M(\nu _l^{}\nu _l)`$ is the transition probability with the matter effects expressed in Eq. (18). The Earth consists of 5 shells approximately ; the density of the most outside shell ($`\mathrm{radius}r=59916371\mathrm{k}\mathrm{m}`$) is $`3.5\mathrm{g}/\mathrm{cm}^3`$, the next outside shell ($`r=57195991\mathrm{k}\mathrm{m}`$) $`4\mathrm{g}/\mathrm{cm}^3`$, the middle shell ($`r=35195719\mathrm{k}\mathrm{m}`$) $`5\mathrm{g}/\mathrm{cm}^3`$, the outer core shell ($`r=12313519\mathrm{k}\mathrm{m}`$) $`11\mathrm{g}/\mathrm{cm}^3`$, the inner core ($`r=01231\mathrm{k}\mathrm{m}`$) $`13\mathrm{g}/\mathrm{cm}^3`$. A number of the shells through which neutrinos pass changes with a change of the zenith angle $`\theta `$. For example, the number of the shell is 0 for $`0<\theta <90^{}`$ and 9 for $`\theta =180^{}`$. The zenith angle distributions of sub-GeV, multi-GeV neutrino events and upward thru- and stop-muons fluxes are given by the SuperKamiokande 850-920 live days experiments . The data are tabulated in Table I. These values are taken from the experimental event data and Monte-Carlo simulations which are given graphically in Ref. . $`\mu `$-like events include the fully contained and partially contained events. Errors represent statistical ones only. Since $`P^M(\nu _l^{}\nu _l)`$ is a function of $`\mathrm{\Delta }m_{12}^2,\mathrm{\Delta }m_{23}^2,\theta _{12},\theta _{13}`$ and $`\theta _{23}`$, the ratio $`(dN_{\mathrm{Exp}}(l)`$ $`/d\mathrm{cos}\theta )/(dN_{\mathrm{MC}}(l)/d\mathrm{cos}\theta )`$ of the zenith angle distributions is a function of $`\mathrm{\Delta }m_{12}^2,\mathrm{\Delta }m_{23}^2,`$ $`\theta _{12},\theta _{13},\theta _{23}`$ and $`\theta `$. We analyze the atmospheric neutrino data fixing the values of parameters $`\mathrm{\Delta }m_{12}^2`$ and $`\mathrm{sin}^22\theta _{12}`$ determined from the solar neutrino experiments as $`\mathrm{\Delta }m_{12}^2=3\times 10^5\mathrm{eV}^2`$ and $`\mathrm{sin}^22\theta _{12}=0.7`$, which corresponds to the large angle solution, and $`\mathrm{\Delta }m_{12}^2=10^5\mathrm{eV}^2`$ and $`\mathrm{sin}^22\theta _{12}=0.005`$, which corresponds to the small angle solution. We treat the ratios of the zenith angle distributions of the experimental events and the ones of Monte-Calro simulation, $`(N_{\mathrm{Exp}}(l)/N_{\mathrm{MC}}(l))_i`$, where $`i`$ represents the region number of the bins of zenith angle $`\theta `$. We calculate numerically the $`\chi ^2`$ defined as $$\chi ^2=\underset{i,l}{}\frac{\left\{(N_{\mathrm{Exp}}(l)/N_{\mathrm{MC}}(l))_i^{\mathrm{cal}}(N_{\mathrm{Exp}}(l)/N_{\mathrm{MC}}(l))_i^{\mathrm{data}}\right\}^2}{(\sigma _{\mathrm{st}})_i^2+(\sigma _{\mathrm{sy}})_i^2}.$$ (21) For the sub-GeV, multi-GeV and upward thru-muon experiments, the summation on $`i`$ are from 1 to 10 of zenith angle range bins and for the upward stop-muon, from 1 to 5. The summation on $`l`$ are over $`\mu `$ and $`e`$ for the sub-GeV and multi-GeV experiments. $`\sigma _{\mathrm{st}}`$ represents the statistical error and $`\sigma _{\mathrm{sy}}`$ systematic one. We assumed that the value of $`\sigma _{\mathrm{sy}}`$ is $`5\%`$ of an amount of the $`N_{\mathrm{MC}}`$ for the sub-GeV and multi-GeV neutrino data, $`10\%`$ for the upward stop-muon data and $`20\%`$ for the upward thru-muon data. We adopted the rather large value of $`\sigma _{\mathrm{sy}}`$ for upward thru-muons, because there exists larger uncertainty of neutrino flux for high energy neutrinos. We estimated the values of $`\chi ^2`$ for the various values of $`\mathrm{\Delta }m_{23}^2,\theta _{13}`$ and $`\theta _{23}`$. In Fig. 1, we showed the contour plots of $`\chi ^2`$ in the $`\mathrm{tan}^2\theta _{13}`$$`\mathrm{tan}^2\theta _{23}`$ plane for various values of $`\mathrm{\Delta }m_{23}^2`$. These plots show the $`\chi ^2`$ of the sub-GeV neutrino plus multi-GeV neutrino zenith angle distributions for the large $`\theta _{12}`$ angle solution. We did not show the plots for the small angle solution, because $`\chi ^2`$ for the small angle solution is almost similar to that for large one. In these figures, broken thick, broken thin and dotted curves denote the regions allowed in 99%, 95% and 90% C.L., respectively. From these plots we can say that the values of allowed $`\mathrm{\Delta }m_{23}^2`$ is from $`2\times 10^3\mathrm{eV}^2`$ to $`10^2\mathrm{eV}^2`$, $`\theta _{13}`$ is $`<17^{}`$ and $`\theta _{23}`$ is from $`35^{}`$ to $`55^{}`$. Plots of the upward thru- plus stop-muons are shown in Fig. 2 for the large $`\theta _{12}`$ angle solution. We did not show the plots for small angle solution, because there is no difference between the large $`\theta _{12}`$ angle solution and small one. In these figures, sharp allowed regions similar to those of previous sub-GeV plus multi-GeV neutrinos data is not appeared. The reason is as follows: although stop-muons data show the sharp allowed regions similar to those of sub-GeV plus multi-GeV neutrino data, thru-muons data show the excluded regions near $`45^{}`$ of $`\theta _{23}`$. The latter situation is caused from small deficit (about 0.2) from 1 of the $`\mu `$ event ratio $`N_{\mathrm{Exp}}/N_{\mathrm{MC}}`$ as shown in Table 1. In Fig. 3, we showed the contour plots of $`\chi ^2`$ for the combination of the sub-GeV and multi-GeV neutrinos, the upward thru- and stop-muons zenith angle distributions for large $`\theta _{12}`$ solution. In these figures, broken thick, broken thin and dotted curves denote the regions allowed in 99%, 95% and 90% C.L., respectively. From these plots, we can get the following results for the neutrino mass and mixing parameters: (1) As shown in Fig. 1, the allowed region for $`\mathrm{\Delta }m_{23}^2`$ obtained from the sub-GeV and multi-GeV neutrino experiments is from $`2\times 10^3\mathrm{eV}^2`$ to $`10^2\mathrm{eV}^2`$ at 90%C.L., and the allowed region for $`\theta _{23}`$ angle is from $`35^{}`$ to $`55^{}`$ and for $`\theta _{13}`$ angle is less than $`17^{}`$. The minimum $`\chi ^2`$ is obtained as $`29`$ for 39 degrees of freedom (DOF) at $`\mathrm{\Delta }m_{23}^2=5\times 10^3\mathrm{eV}^2`$, $`\theta _{13}=10^{}`$ and $`\theta _{23}=45^{}`$. The minimum $`\chi ^2`$ in the restriction $`\theta _{13}=0`$ which corresponds to the two flavor mixing $`\nu _\mu \nu _\tau `$ approximation, is $`30`$ for 39 DOF at $`\mathrm{\Delta }m_{23}^2=4\times 10^3\mathrm{eV}^2`$. (2) As shown in Fig. 2, the allowed region for $`\theta _{13}`$ and $`\theta _{23}`$ is excluded at small $`\theta _{13}`$ and the maximal mixing of $`\theta _{23}`$ for large values of $`\mathrm{\Delta }m_{23}^2`$. This is because of the fact that the data in upward thru-muons has not so large deficit from the no-oscillation case as seen in the Table 1. (3) There is no significant difference between the large $`\theta _{12}`$ solution and the small solution. This fact is not seen from these figures, but we can see from our numerical calculation that the difference between the large angle solution and the small one is less than 1% for multi-GeV neutrinos and upward muons and less than 10 % for sub-GeV neutrinos. (4) The allowed region for $`\mathrm{\Delta }m_{23}^2`$ obtained from the sub- and multi-GeV neutrinos, and upward thru- and stop-muons is given as $`2\times 10^3\mathrm{eV}^2`$ to $`\times 10^2\mathrm{eV}^2`$, and for $`\theta _{23}`$ as $`35^{}`$ to $`55^{}`$ and for $`\theta _{13}`$ as less than $`13^{}`$. The minimum $`\chi ^2`$ is obtained as $`55`$ for 54 DOF at $`\mathrm{\Delta }m_{23}^2=4\times 10^3\mathrm{eV}^2`$, $`\theta _{13}=10^{}`$ and $`\theta _{23}=45^{}`$. The minimum $`\chi ^2`$ in the restriction $`\theta _{13}=0^{}`$ which corresponds to the two flavor mixing $`\nu _\mu \nu _\tau `$ approximation, is $`61`$ for 54 DOF at $`\mathrm{\Delta }m_{23}^2=3\times 10^3\mathrm{eV}^2`$ and $`\theta _{23}=45^{}`$. These results are the same as the result obtained by Super-Kamiokande collaboration . (5) It is interesting that the minimum of $`\chi ^2`$ is obtained at not $`\theta _{13}=0`$ but $`\theta _{13}=10^{}`$, though the difference between $`\chi ^2`$ for $`\theta _{13}=0`$ case and that for $`\theta _{13}=10^{}`$ case is not so large. This result is consistent with the CHOOZ experiment . The CHOOZE experiment predicts the results: $`\theta _{13}<13^{}`$ for $`\mathrm{\Delta }m_{23}^2=3\times 10^3\mathrm{eV}^2`$. We will now discuss about the interesting feature that the mixing angle $`\theta _{13}`$ is preferred to be about $`10^{}`$. If this feature is real, the detected $`\nu _e`$ events in the long baseline K2K experiment will be about 10 times as large as the events expected in $`\theta _{13}=0`$ case. This is caused from the fact that the $`\nu _\mu `$ flux produced at KEK is about 100 times as large as $`\nu _e`$ flux and the transition probability is $`P(\nu _\mu \nu _e)\mathrm{sin}^22\theta _{13}\mathrm{sin}^21.27\mathrm{\Delta }m_{23}^2L/E`$, $`E1.4\mathrm{GeV}`$ and $`L=250\mathrm{k}\mathrm{m}`$. ## 4 Conclusion We analyzed the atmospheric neutrino experimental data of Super-Kamiokande in the three-flavor neutrino framework with the mass hierarchy $`m_1m_2m_3`$ and obtained the allowed regions of parameters $`\mathrm{\Delta }m_{23}^2,\theta _{13}`$ and $`\theta _{23}`$, including the Earth matter effects thoroughly. We studied the event ratios of the sub- and multi-GeV, and upward thru- and stop-muons zenith angle distributions. From these atmospheric experiments, we can get the allowed region of mass parameter $`\mathrm{\Delta }m_{23}^2`$ restricted as $`0.002\mathrm{eV}^2\text{}0.01\mathrm{eV}^2`$, and mixing parameter $`\theta _{13}`$ as less than $`13^{}`$ and $`\theta _{23}`$ as $`35^{}\text{}55^{}`$. The value of $`\mathrm{\Delta }m_{23}^2`$ at the minimum $`\chi ^2=55`$ for 54 DOF is obtained as $`4\times 10^3\mathrm{eV}^2`$ at $`\theta _{13}=10^{}`$ and $`\theta _{23}=45^{}`$. The minimum $`\chi ^2=61`$ for 54 DOF is obtained with the restriction $`\theta _{13}=0`$ at the $`\mathrm{\Delta }m_{23}^2=3\times 10^3\mathrm{eV}^2`$ and $`\theta _{23}=45^{}`$. This fact seems very interesting for us because the mixing parameter $`\theta _{13}`$ may not be 0. For mixing parameter $`\theta _{12}`$, the difference between the large angle solution and the small one is less than 1%. If $`\theta _{13}`$ is about $`10^{}`$, the detected $`\nu _e`$ events in K2K experiment is about 10 times as large as events expected in $`\theta _{13}=0`$ case.
warning/0003/math0003100.html	ar5iv	text	# Quantum Strata of Coadjoint Orbits ## Introduction For locally compact groups, their C\-algebras contain exhausted informations about the groups them-selves and their representations, see \[D1\], \[D2\]. In some sense \[R1\]-\[R2\], the group algebras can be considered as C\-algebraic deformation quantization $`C_q^{}(G)`$ at the special value $`q=1`$. In \[D1\] and \[D2\], the group C\-algebras were described as repeated extensions of C\-algebras of strata of coadjoint orbits. Quantum groups are group Hopf algebras, i.e. replace C\-algebras by special Hopf algebras “of functions”. It is therefore interesting to ask whether we could describe quantum groups as some repeated extensions of some kind quantum strata of coadjoint orbits? We are attempting to give a positive answer to this question. It is not yet completely described but we obtained a reasonable answer. Let us describe the main ingredients of our approach. For the good strata, namely families of with some good enough parameter space, of coadjoint orbits, there exist always continuous fields of polarizations (in the sense of the representation theory), satisfying the L. Pukanszky irreducibility condition: for each orbit $`𝒪`$ and any point $`F_𝒪`$ in it, the affine subspace, orthogonal to some polarizations with respect to the symplectic form is included in orbits themselves, i.e. $$F_𝒪+𝔥_𝒪^{}𝒪$$ and $$dim𝔥_𝒪=\frac{1}{2}dim𝒪.$$ We choose the the canonical Darboux coordinates with impulse $`p`$’s-coordinates, following a vector structure basis of $`𝔥^{}`$. From this we can deduce that in this kind of Darboux coordinates, the Kirillov form $`\omega _𝒪`$ locally are canonical and every element $`X𝔤=LieG`$ can be considered as a function $`\stackrel{~}{X}`$ on $`𝒪`$, linear on $`p`$’s-coordinates, i.e. $$\stackrel{~}{X}=\underset{i=1}{\overset{n}{}}a_i(q)p_i+a_0(q).$$ This essential fact gives us a possibility to effectively write out the corresponding $``$-product of functions, define quantum strata $`C_q(V,)`$. On the strata acts our Lie group of symmetry. It induces therefore an action on equivariant differential operators. Using the indicate fields of polarizations we prove some kind of Poincaré-Birkhoff-Witt theorem and then provide quantization with separation of variable in sense of Karabegov \[Ka1\]. We can then express the corresponding representations of the quantum strata $`C_q(V,)`$, where $`V_{dim𝒪=const}𝒪`$, through the Feynman path integrals etc…,see \[D3\]. Our main result is Theorem 2.6 in which we construct the $``$-product exactly, Theorem 3.4, where the representations are obtained from Poincaré-Birkhoff-Witt separation of variable and $``$-products, and Theorem 5.5 where we express the product through Fourier integral operators. Let us now describe in brief the structure of sections. In section 1, we construct a canonical local coordinate system, where the generating functions $`\stackrel{~}{X}`$, for $`X𝔤`$ are linear in the impulse coordinates $`p`$’s, Theorem 1.4. In this kind of coordinates, we can in §2 construct construct the local $``$-product, Theorem 2.1, then globalize it into a $`\mathrm{\Gamma }`$-invariant $``$-product on the universal coverings of coadjoint orbits and then push down to the coadjoint orbits them-selves, Theorem 2.4, where $`\mathrm{\Gamma }=\pi _1(𝒪)`$ \- the fundamental group of the coadjoint orbit $`𝒪`$. IN section 3 we let this left $``$-multiplication of functions, acting on the quantum bundle sections. It manipulates an action of the generating functions as right $`G`$-invariant pseudo-differential operators of the first order. We then use the universal property of the universal enveloping algebra $`U(𝔤)`$ to construct a quantizing homomorphism $`U(𝔤)PDO_G(𝒪)`$, Theorem 3.1. We prove then a version of the Poincaré-Birkhoff-Witt Theorem, associated with a complex polarization. Theorem 3.3. In section 4, we show that this kind representations are the same as those obtained from the procedure of multidimensional quantization, \[D2\]. In section 5, we restrict to special case of strata of coadjoint orbits appeared from some versions of solvable cases of the Gelfand-Kirillov conjecture. In those cases we can express the associated unitary representations of the groups $`G`$ as some oscilatting Fourier integrals, Theorem 5.5. In Section 6, all the main ideas are demonstrated in 3 series of examples. ## 1. Canonical coordinates on a stratum Let us consider a connected and simply connected Lie group $`G`$ with Lie algebra $`𝔤`$. Denote the dual to $`𝔤`$ vector space by $`𝔤^{}`$. It is well-known that the action of $`G`$ on itself by conjugation $$A(g):GG,$$ defined by $`A(g)(h):=ghg^1`$ keeps the identity element $`h=e`$ unmoved. This induces the tangent map $`Ad(g):=A(g)_{}:𝔤=T_eG𝔤,`$ defined by $$Ad(g)X:=\frac{d}{dt}\|_{t=0}A(g)\mathrm{exp}(tX)$$ and the co-adjoint action $`K(g):=Ad(g^1)^{}`$ maps the dual space $`𝔤^{}`$ into itself. The orbit space $`𝒪(G):=𝔤^{}/G`$ is in general a bad topological space, namely non-Hausdorff, in general. Consider an arbitrary orbit $`\mathrm{\Omega }𝒪(G)`$ and an element $`F𝔤^{}`$ in it. The stabilizer is denote by $`G_F`$, its connected component by $`(G_F)_0`$ and its Lie algebra by $`𝔤_F:=Lie(G_F)`$. It is well-known that $$\begin{array}{ccc}G_F& & G\\ & & \\ & & \mathrm{\Omega }_F\end{array}$$ is a principal bundle with the structural group $`G_F`$. Let us fix some connection in this principal bundle, i.e. some trivialization of this bundle. We want to construct representations in some cohomology spaces with coefficients in the sheaves of sections of some vector bundle associated with this principal bundles. It is well know that every vector bundle is an induced one with respect to some representation of the structural group in the typical fiber. It is natural to fix some unitary representation $`\stackrel{~}{\sigma }`$ of $`G_F`$ such that its kernel contains $`(G_F)_0`$, the character $`\chi _F`$ of the connected component of stabilizer $$\chi _F(\mathrm{exp}X):=\mathrm{exp}(2\pi \sqrt{1}F,X)$$ and therefore the differential $`D(\stackrel{~}{\sigma }\chi _F)=\stackrel{~}{\rho }`$ is some representation of the Lie algebra $`𝔤_F`$. We suppose that the representation $`D(\stackrel{~}{\rho }\chi _F)`$ was extended to the complexification $`(𝔤_F)_{}`$. The whole space of all sections seems to be so large for the construction of irreducible unitary representations. One considers the invariant subspaces with the help of some so called polarizations, see \[D1\], \[D2\]. ###### Definition 1.1. We say that a triple $`(𝔭_𝒪,\rho _𝒪,\sigma _{0,𝒪})`$ is a $`(\stackrel{~}{\sigma },F)`$-polarization of $`𝒪`$ iff : 1. $`𝔭_𝒪`$ is some complex sub-algebra of the complexified $`𝔤_𝐂`$, containing $`𝔤_{F_𝒪}`$. 2. The sub-algebra $`𝔭_𝒪`$ is invariant under the action of all the operators of type $`Ad_{𝔤_𝐂}x,xG_{F_𝒪}.`$ 3. The vector space $`𝔭_𝒪+\overline{𝔭_𝒪}`$ is the complexification of some real Lie sub-algebra $`𝔪_𝒪:=(𝔭_𝒪+\overline{𝔭}_𝒪)𝔤.`$ 4. All the subgroups $`M_{0,𝒪}`$, $`H_{0,𝒪}`$, $`M_𝒪`$, $`H_𝒪`$ are closed, where by definition, $`M_{0,𝒪}`$ (resp., $`H_{0,𝒪}`$) is the connected subgroup of $`G`$ with the Lie algebra $`𝔪_𝒪`$ (resp., $`𝔥_𝒪:=𝔭_𝒪𝔤`$) and the semi-direct products $`M:=M_{0,𝒪}G_{F_𝒪}`$, $`H_𝒪:=H_{0,𝒪}G_{F_𝒪}`$. 5. $`\sigma _{0,𝒪}`$ is an irreducible representation of $`H_{0,𝒪}`$ in some Hilbert space $`V_𝒪`$ such that: (1.) the restriction $`\sigma _{0,𝒪}\|_{G_{F_𝒪}H_{0,𝒪}}`$ is some multiple of the restriction $`\chi _{F_𝒪}.\stackrel{~}{\sigma }_𝒪\|_{G_{F_𝒪}H_{0,𝒪}}`$, where the character $`\chi _𝒪`$ is by definition, $`\chi _𝒪(\mathrm{exp}X)=\mathrm{exp}(2\pi \sqrt{1}F_𝒪,X)`$; (2.) under the action of $`G_{F_𝒪}`$ on the dual $`\widehat{H}_{0,𝒪}`$, the point $`\sigma _{0,𝒪}`$ is fixed. 6. $`\rho _𝒪`$ is a representation of the complex Lie algebra $`𝔭_𝒪`$ in the same $`V_𝒪`$, which satisfies the Nelson conditions for $`H_{0,𝒪}`$ and $`\rho _𝒪\|_{𝔥_𝒪}D\sigma _{0,𝒪}`$. There is a natural order in the set of all $`(\stackrel{~}{\sigma }_𝒪,F_𝒪)`$-polarizations by inclusion and from now on speaking about polarizations we means always the maximal ones. It is not hard to prove that the (maximal) polarizations are also the Lagrangian distributions and in particular the co-dimension of $`𝔥_𝒪`$ in $`𝔤`$ is a half of the dimension of the coadjoint orbit $`𝒪`$, $$codim_𝔤𝔥_𝒪=\frac{1}{2}dim𝒪,$$ see e.g. \[D2\] or \[K\]. Let us now recall the Pukanszky condition. ###### Definition 1.2. We say that the $`(\stackrel{~}{\sigma }_𝒪,F_𝒪)`$-polarization $`(𝔭_𝒪,\rho _𝒪,\sigma _{0,𝒪})`$ satisfies the Pukanszky condition, iff $$F_𝒪+𝔥_𝒪^{}𝒪.$$ ###### Remark 1.3. Pukanszky conditions involve an inclusion of the Lagrangian affine subspace of $`p`$’s coordinates into the local Darboux coordinates. * The partial complex structure on orbits let us to use smaller subspaces of section in the induction construction, as subspaces of partially invariant, partially holomorphic sections of the induced bundles. ###### Theorem 1.4. There exists on each coadjoint orbit a local canonical system of Darboux coordinates, in which the Hamiltonian function $`\stackrel{~}{X}`$, $`X𝔤`$, are linear on $`p`$’s impulsion coordinates and in theses coordinates, $$\stackrel{~}{X}=\underset{i=1}{\overset{n}{}}a_i(q)p_i+a_0(q).$$ ###### Proof. We supposed that our coadjoint orbit admit at least a polarization $`(𝔭,\rho ,\sigma )`$, satisfying L. Pukanszky’s condition of irreducibility $$F_𝒪+𝔥_𝒪^{}𝒪.$$ The codimension of $`𝔥_𝒪`$ and therefore the dimension of $`𝔥_𝒪^{}`$ is $`\frac{1}{2}dim𝒪`$. Consider a canonical system of geodesics. The geodesics corresponding to the affine subspace $`𝔥_𝒪^{}`$ provide linear coordinates $`p_1,p_2,\mathrm{},p_n`$. The others are the corresponding $`q^1,\mathrm{},q^n`$. Therefore we can arrange a local system of coordinates, such that exponential map gives linear geodesics on $`p`$’s directions. ∎ ## 2. Poisson structure on strata Let us first recall the construction of strata of coadjoint orbit from \[D1\]-\[D2\]. The orbit space $`𝒪(G)`$ is a disjoint union of $`\mathrm{\Omega }_{2n}`$, each of which is a union of the coadjoint orbits of dimension $`2n`$, $`02ndimG`$. Denote $$V_{2n}=_{dim𝒪=2n}𝒪.$$ Then $`V_{2n}`$ is the set of points of fixed rank $`2n`$ of the Poisson structure bilinear function $$\{X,Y\}(F)=F,[X,Y].$$ Suppose that it is a foliation, at least it is true for $`V_{2n}`$, where $`2n`$ is the maximal dimension possible in $`𝒪(G)`$. It can be shown that the foliation $`V_{2n}`$ can be obtained from the group action of $`𝐑^{2n}`$ on $`V_{2n}`$. For this aim, let us consider a basis $`X_1,\mathrm{},X_{2n}`$ of the tangent space $`T_{F_𝒪}𝒪𝔤/𝔤_{F_𝒪}`$ at the point $`F_𝒪𝒪V_{2n}`$. We can define an action of $`^{2n}`$ on $`V_{2n}`$ as $$(t_1,\mathrm{},t_{2n})\mathrm{exp}(t_1X_1)\mathrm{exp}(t_2X_2)\mathrm{}\mathrm{exp}(t_{2n}X_{2n})F_𝒪.$$ We have therefore the Hamiltonian vector fields $$\xi _k:=\frac{d}{dt}\|_{t=0}\mathrm{exp}(t_kX_k)F,k=1,\mathrm{},2n$$ and their span $`_{2n}=\{\xi _1,\mathrm{},\xi _{2n}\}`$ provides a tangent distribution. It is easy to show that we have therefore a measurable (in sense of A. Connes) foliation. One can therefore define also the Connes C\-algebra $`C^{}(V_{2n},_{2n})`$, $`02ndimG`$. By introducing some technical condition in \[D1\], it is easy to reduce these C\-algebras to extension of other ones, those are in form of tensor product $`C(X)𝒦(H)`$ of algebras of continuous functions on compacts and the elementary algebra $`𝒦(H)`$ of compact operators in a separable Hilbert space $`H`$. The strata of these kinds we means good strata. Another kind of good strata of coadjoint orbits are obtained from relation with the cases where the Gelfand-Kirillov conjecture was solved, for examples for connected and simply connected solvable Lie groups, see §5. It is deduced from a result of Kontsevich that this Poisson structure can be quantized. This quantization however is formal. The question of convergence of the quantizing series is not clear. We show in this section that in the case of charts with the linear p’s impulse coordinates, the corresponding $``$-product is convergent. In this kind of special local chart systems of Darboux coordinates it is easy to deduce existence local convergent $``$-products. ###### Theorem 2.1. Locally on each coadjoint orbit, there exist a convergent $`star`$-product. ###### Proof. Let us denote by $`_p^1`$ the Fourier inverse transformation on variables $`p`$’s and by $`_p`$ the Fourier transformation. Let us denote by $`PDO_G(𝒪)`$ the algebra of $`G`$-invariant pseudodifferential operators on $`𝒪`$. Locally, the Fourier transformation maps symbols (as function on local coordinates of $`𝒪`$) to $`G`$-invariant pseudodifferential operators and conversely, the inverse Fourier transformation maps the pseudodifferential operators to some specific classes of symbols. For two symbols $`f,g_p^1(PDO_G(𝒪))=𝐂_q(𝒪)`$, their Fourier images $`_p(f)`$ and $`_p(g)`$ are operators and we can define their operator product and then take the Fourier inverse transforms, as $``$-product $$fg:=_p^1(_p(f)._p(g)).$$ So this product is again a symbol n the same class $`𝐂_q(𝒪)`$. Because of existence of special coordinate systems, linear on $`p`$’s coordinates we can treat for the good strata in the same way as in the cases of exponential groups. And the formal series of $``$-product is convergent. ∎ ###### Remark 2.2. The proof could be also done in the same scheme as in exponential or compact cases, see \[AC1\]-\[AC2\]. Let us denote by $`\mathrm{\Gamma }=\pi _1(𝒪)`$ the fundamental group of the orbit. Our next step is to globally extend this kind of local $``$-products. Our idea is as follows. We lift this $``$-product to the universal covering of coadjoint orbits as some $`\mathrm{\Gamma }`$-invariant $``$-products, globally extend them in virtue of the monodromy theorem and then pushdown to our orbits. We start with the following lemmas ###### Lemma 2.3. There is one-to-one correspondence between $``$-products on Poisson manifolds and $`\mathrm{\Gamma }`$-invariant $``$-products on their universal coverings. ###### Proof. By the lifting properties of the universal covering, we can easily lift each $``$-product on a Poisson manifold onto its universal covering. This correspondence is one-to-one. ∎ We use this lemma to describe existence of a $``$-product on coadjoint orbits. ###### Lemma 2.4. On a universal covering, each local $``$-product can be uniquely extended to some global $``$-product on this covering. ###### Proof. For local charts, there exist deformation quantization, as said above, $`fOp(f)`$ by using the formulas of Fedosov quantization. Also in the intersection of two local charts of coordinates $`(q,p)`$ and $`(\stackrel{~}{q},\stackrel{~}{p})`$, there is a symplectomorphism, namely $`\phi `$ such that $$(\stackrel{~}{q},\stackrel{~}{p})=\phi (q,p).$$ Using the local oscilatting integrals and by compensating the local Maslov’s index obstacles, one can exactly construct the unitary operator $`U`$ such that $$Op(f\phi )=U.Op(f).U^1,$$ see for example Fedosov’s book \[F\]. Because the universal coverings are simply connected, the extensions can be therefore produced because of Monodromy Theorem. ∎ ###### Remark 2.5. The operators “U” of this kind, depending on two local charts as parameters, provide some cohomological 2-class and therefore are classified by some $`2^{nd}`$ cohomology class with values in unitary operators. This reduces to some classification of all the possibilities of quantizations, upto the first cohomology classes, i.e. upto conjugations. ###### Theorem 2.6. There exists a convergent $``$-product on each orbit, which is a symplectic leaf of the Poisson structure on each stratum of coadjoint orbits. ###### Proof. From the description of the canonical coordinates in the previous section, we see that there exists at least a convergent $`\mathrm{\Gamma }`$-invariant local $``$-product. This local $``$-product then extended to a $`\mathrm{\Gamma }`$-invariant global one on the universal covering, which produces a convergent $``$-product on the coadjoint orbits, following Lemmas 2.3,2.4. ∎ ## 3. Star-product, Quantization and PBW Theorem We use the constructed in the previous section $``$-product to provide an action of functions on the spaces of partially invariant partially holomorphic sections of the corresponding partially invariant holomorphically induced bundles associated with polarizations. It is possibles because the quantum induced bundles are locally trivial and the spaces of partially invariant partially holomorphic sections with section in Hilbert spaces locally are finitely generated modules over the algebras of quantizing functions, see \[D2\]. We use then the construction of Karabegov and Fedosov to obtain a Hopf \-algebra of longitudinal pseudo-differential operators elliptic along the leaves of this measurable foliations. The main ingredient is that we use here the Poincaré-Birkhopf-Witt theorem to provide this quantization. ###### Theorem 3.1. There is a natural deformation quantization with separation of variables, corresponding to the Poincaré-Birkhoff-Witt Theorem for polarizations. We have from the multidimensional quantization, see e.g. \[D1\]-\[D2\], $`𝔤PDO^1(𝒪)`$, $`X\widehat{X}`$. More precisely, Let us denote by $`PDO(𝒪)`$ the algebra of right $`G`$-invariant pseudo-differential operators on $`𝒪`$, i.e. the continuous maps from $`C^{\mathrm{}}(𝒪)`$ into itself not extending support. We also denote $`PDO^1(𝒪)`$ the Lie algebras of right $`G`$-invariant first order pseudo-differential operators. By the procedure of multidimensional quantization, there is a homomorphism of Lie algebras $$𝔤PDO_G^1(𝒪)PDO_G(𝒪).$$ Following the universal property of $`U(𝔤)`$, there is a unique homomorphism of associative algebras $`U(𝔤)PDO_G(𝒪)`$making the following diagram commutative $$\begin{array}{ccc}𝔤& & PDO_G(𝒪)\\ & & & & \\ U(𝔤)& =& U(𝔤)\end{array}$$ Let us first describe the machinery applied in the method of Karabegov’s separation of variable. We use the idea about polarizations in multidimensional quantization, \[D2\]. Let us recall the root decomposition $$𝔤=𝔫_{}𝔞𝔫_+.$$ If $`𝔭`$ is a complex polarization, then from definition we have $`𝔪_𝐂=𝔭\overline{𝔭},`$ $`𝔥_𝐂=𝔭\overline{𝔭},`$ where $`𝔪=𝔤(𝔭\overline{𝔭})`$ and $`𝔥=𝔤(𝔭\overline{𝔭}).`$ The quotients subspaces $`𝔭/𝔥_𝐂`$ and $`\overline{𝔭}/𝔥_𝐂`$ are included in $`𝔪`$ as linear subspaces (not necessarily to be sub-algebras). Let us fix some (non-canonical) inclusions. Let us denote by $`U(𝔭/𝔥_𝐂)`$ (resp. $`U(\overline{𝔭}/𝔥_𝐂)`$) the sub-algebra, generated by elements from $`𝔭/𝔥_𝐂𝔪_𝐂`$ (resp. $`𝔭/𝔥_𝐂𝔪_𝐂`$) in the universal algebra $`U(𝔪)`$. We have therefore an analog of the well-known Poincaré-Birkhoff-Witt theorem. ###### Theorem 3.2 (Poincaré-Birkhoff-Witt Theorem). If $`𝔭`$ is as polarization, then $$U(𝔪_𝐂)U(𝔭/𝔥_𝐂)U(𝔥_𝐂)U(\overline{𝔭}/𝔥_𝐂)$$ ###### Proof. Let us fix in a basis each of $`𝔭/𝔥_𝐂`$, $`𝔥_𝐂`$ and $`\overline{𝔭}/𝔥_𝐂`$. We have therefore a basis of $`𝔪_𝐂=𝔭/𝔥_𝐂𝔥_𝐂\overline{𝔭}/𝔥_𝐂`$. Our theorem is therefore deduced from the original Poincaré-Birkhoff-Witt Theorem \[P\]. ∎ It is easy to see that by this reason, on the variety $`M/H`$ there is a natural complex structure and therefore on coadjoint orbits exists some partial complex structure. We use this complex structure and this Poincaré-Birkhoff-Witt theorem to do a separation of variables on $`M`$ and apply the machinery of Karabegov. ###### Remark 3.3. Because of BKW, $`U(𝔤)U(𝔭_𝒪/𝔥_𝒪)U(𝔥_𝒪)U(\overline{𝔭}_𝒪/𝔥_𝒪)`$ and because of this our quantizing map is coincided with that one used by Karabegov in the quantization with separation of variables in case of coadjoint orbits with totally complex polarizations $`𝔤_𝐂=𝔪_𝐂=𝔭\overline{𝔭}`$. ###### Theorem 3.4. In the case of totally complex polarizable coadjoint orbits, the quantization map from the theorem 3.1 is coincided with the rule of Karabegov’s quantization with separation of variables. ## 4. Representations Let us consider now a continuous fields of $`(\stackrel{~}{\sigma }_𝒪,F_𝒪)`$-polarizations $`(𝔭_𝒪,\rho _𝒪,\sigma _{0,𝒪})`$ satisfying the Pukanszky condition. On one hand side, we can use the multidimensional quantization procedure to obtain irreducible unitary representations $`\mathrm{\Pi }_𝒪=Ind(G;𝔭_𝒪,H_𝒪,\rho _𝒪,\sigma _{0,𝒪})`$ of $`G`$, \[D2\]. On other hand side, we can use $``$-product construction to provide the representations $`\mathrm{\Pi }_𝒪:\stackrel{~}{X\|_\mathrm{\Omega }}\widehat{\mathrm{}}_X`$ of the quantum strata $`𝐂_q(\mathrm{\Omega })`$. We’d like to show we have the same one. ###### Definition 4.1. Quantum coadjoint orbit $`𝐂_q(𝒪)`$ is defined as the Hopf algebra of symbols of differential operators $`U_{q,𝒪}(𝔤)=U(𝔤)\|_𝒪`$. The homomorphism $`Q:U(𝔤)PDO_G(𝒪)`$ is defined to be the second quantization homomorphism. ###### Theorem 4.2. The representation of the Lie algebra obtained from $``$-product is equal to the representations obtained from the multidimensional quantization procedure. ###### Proof. Let us recall \[D2\], that $$Lie_XInd(G;𝔭_𝒪,H_𝒪,\rho _𝒪,\sigma _{0,𝒪})\widehat{X}.$$ From the construction of quantization map $$U(𝔤)PDO_G(𝒪)$$ as the map arising from the universal property the map $`𝔤PDO_G^1(𝒪),`$ $$\widehat{\mathrm{}}_X=Op(\stackrel{~}{X})=\widehat{X},X𝔤=Lie(G).$$ The associated representation of $`𝐂_q(\mathrm{\Omega })`$ is the solution of the Cauchy problem for the differential equation $$\frac{}{t}U(t,q,p)=\mathrm{}_XU(t,q,p),$$ $$U(0,q,p)=Id.$$ The solution of this problem is uniquely defined. ∎ ## 5. Oscillating Fourier integrals Let us in this section consider the good family of coadjoint orbits arising from the solved cases of Gelfand-Kirillov Conjecture. The results in this section is revised from the \[D3\]. Consider a connected and simply connected Lie group $`G`$ with Lie algebra $`𝔤`$. ###### Theorem 5.1. There exists a $`G`$-invariant Zariski open set $`\mathrm{\Omega }`$ and a covering $`\stackrel{~}{\mathrm{\Omega }}`$ of $`\mathrm{\Omega }`$, with natural action of $`G`$ such that for each continuous field of polarizations $`(𝔭_𝒪,H_𝒪,\rho _𝒪,\sigma _{0,𝒪})`$, $`𝒪\mathrm{\Omega }/G`$, the Lie derivative of the direct integral of representations arized from the multidimensional quantization procedure is equivalent to the tensor product of the Schrödinger representation $`\mathrm{\Pi }=Sch`$ of the Gelfand-Kirillov basis of the enveloping field and a continuous field of trivial representations $`\{V_𝒪\}`$ on $`\stackrel{~}{\mathrm{\Omega }}`$. ###### Proof. Our proof is rather long and consists of several steps: 1. We apply the construction of Nghiem \[Ng5\] to the solvable radical $`{}_{}{}^{r}𝔤`$ of $`𝔤`$ to obtain the Zariski open set $`{}_{}{}^{r}\mathrm{\Omega }`$ in $`{}_{}{}^{r}𝔤_{}^{}`$ and its covering $`{}_{}{}^{r}\stackrel{~}{\mathrm{\Omega }}`$. 2. The general case is reduced to the semi-simple case $`{}_{}{}^{s}𝔤`$, see (\[nghien1\], Thm. C). Denote by $`{}_{}{}^{s}G`$ the corresponding analytic subgroup of $`G`$. 3. Take the Zariski open set $`𝒜_s𝒫(^sG)`$ of admissible and well-polarizable strongly regular functionals from $`{}_{}{}^{s}𝔤_{}^{}`$ and its covering $`_s(^sG)`$ via Duflo’s construction \[Du1\]. 4. The desired $`G`$-invariant Zariski open set and its covering are the corresponding Cartesian products $$\mathrm{\Omega }=𝒜_s𝒫(^sG)\times {}_{}{}^{r}\mathrm{\Omega }_{}^{0}\times {}_{}{}^{r}\mathrm{\Omega }_{}^{1}\times \mathrm{}\times {}_{}{}^{r}\mathrm{\Omega }_{}^{k}$$ $$\stackrel{~}{\mathrm{\Omega }}=_s(^sG)\times {}_{}{}^{r}\stackrel{~}{\mathrm{\Omega }^0}\times {}_{}{}^{r}\stackrel{~}{\mathrm{\Omega }^1}\times \mathrm{}\times {}_{}{}^{r}\stackrel{~}{\mathrm{\Omega }}_{}^{k}$$ 5. ###### Lemma 5.2. There exists a continuous field of polarizations of type $`(𝔭_𝒪,H_𝒪,\rho _𝒪,\sigma _{0,𝒪})`$, for each $`𝒪\stackrel{~}{\mathrm{\Omega }}/G`$. 6. ###### Lemma 5.3. The Lie derivative commutes with direct integrals, i.e. $$_{\stackrel{~}{\mathrm{\Omega }}/G}^{}Ind(G;𝔭_𝒪,H_𝒪,\rho _𝒪,\sigma _{0,𝒪})𝑑𝒪=_{\stackrel{~}{\mathrm{\Omega }}/G}^{}Ind(G;𝔭_𝒪,H_𝒪,\rho _𝒪,\sigma _{0,𝒪})𝑑𝒪.$$ 7. ###### Lemma 5.4. The restriction of the Schrödinger representation to coadjoint orbits provides a continuous field of polarizations. In particular, $$D\sigma _{0,𝒪}=\rho _𝒪\|_{𝔭_𝒪𝔤}multSch_𝒪.$$ ###### Theorem 5.5. 1. There exists an operator-valued phase $`\mathrm{\Phi }(t,z,x)`$ and an operator-valued amplitude $`a(t,z,x,y,\xi )`$ extended from the expression $$\mathrm{exp}(\mathrm{\Phi }(t,z,y)+\sqrt{1}\xi (g(t)xx))$$ in such a fashion that the action of the representation $$\mathrm{\Pi }_𝒪=Ind(G;𝔭_𝒪,H_𝒪,\rho _𝒪,\sigma _{0,𝒪})$$ can be expressed as an oscilatting Fourier integral $$\mathrm{\Pi }_𝒪(g(t))f(z,x)=c._{𝐑^M}_{𝐑^M}a(t,z,x,y,\xi )\mathrm{exp}(\sqrt{1}\xi (xy))f(z,y)𝑑y𝑑\xi ,$$ where $`c`$ is some constant. 2. For each function $`\phi `$ of Schwartz class $`𝒮(G)`$ satisfying the compactness criteria \[D2\] in every induction step \[Li\], \[D2\], the operator $`\mathrm{\Pi }_𝒪(\phi )`$ is of trace class and its action can be expressed as the oscilatting Fourier integral $$\mathrm{\Pi }_𝒪(\phi )f(z,x)=const._{𝐑^M}_{𝐑^M}(_𝐑^{\mathrm{}}a(t,z,x,y,\xi )\phi (g(t))dt)\times $$ $$\times \mathrm{exp}(\sqrt{1}\xi (xy))f(z,y)dyd\xi .$$ Hence, its trace is $$tr\mathrm{\Pi }_𝒪(\phi )=const._{𝐑^M}_{𝐑^M}(_𝐑^{\mathrm{}}a(t,z,x,x,\xi )\phi (g(t))𝑑t)𝑑x𝑑\xi $$ ###### Proof. The proof also consists of several steps: 1. From \[D1\] \- \[D2\] and \[Du1\] it is easy to see obtain a slight unipotent modification( i.e. a reduction to the unipotent radical) of the multidimensional quantization procedure. We refer the reader to \[TrV\]-\[TrV2\], and \[TDV\] for a detailed exposition. 2. From the unitary representations $`\mathrm{\Pi }_𝒪=Ind(G;𝔭_𝒪,H_𝒪,\rho _𝒪,\sigma _{0,𝒪})`$ in the unipotent context and its differential $`\pi _𝒪=D\mathrm{\Pi }_𝒪(G;𝔭_𝒪,H_𝒪,\rho _𝒪,\sigma _{0,𝒪})`$, it is easy to select the constant term and the vector fields term, (see \[Ng4\] for the simplest case and notations), $$\pi (L)=\pi _c(L)+\pi _d(L),\pi _s(L)=a_0(L,x),$$ $$\pi _d(L)=(L)=\underset{k}{}a_k(L)\frac{}{x_k}.$$ The operator-valued partial (i.e. depending on $`L`$) phase $$\mathrm{\Phi }_L(t_L,z,x)=_0^{t_L}a_0(z,g(s)x)𝑑s.$$ The phase $`\mathrm{\Phi }(t,z,x)`$ is the sum of the partial phases, on which the one-parameter group $`g_L(s)`$ operates by translations for all the factors on the left of $`g_L(t_L)`$ in the ordered product $$g(t)=\underset{L}{}g_L(t_L).$$ It is easy then to see that our induced representations $`\mathrm{\Pi }_𝒪`$ act as follows $$\mathrm{\Pi }_𝒪(g(t))f(z,x)=\sigma _{0,𝒪}(t,z,x)f(z,g(t)x).$$ 3. It suffices now to apply the Fourier transforms $$f(z,x)=(2\pi )^M_{𝐑^M}_{𝐑^M}\mathrm{exp}\{\sqrt{1}\xi (xy)\}f(z,y)𝑑y𝑑\xi $$ to the function $`\mathrm{\Pi }_𝒪(g(t))f(z,x)`$ to obtain $$\mathrm{\Pi }_𝒪(g(t))f(z,x)=$$ $$(2\pi )^M_{𝐑^M}_{𝐑^M}a(t,z,x,y,\xi )\mathrm{exp}\{\sqrt{1}\xi (xy)\}f(z,y)𝑑y𝑑\xi ,$$ where the amplitude $`a(t,z,x,y,\xi )`$ is the natural extension of the expression $$\mathrm{exp}\{\mathrm{\Phi }(t,z,y)+\sqrt{1}\xi (g(t)xx)\}$$ in the correspondence with the fields of polarizations. Hence for each $`\phi 𝒮(G)`$, $$\mathrm{\Pi }_𝒪(\phi )f(z,x)=(2\pi )^M_{𝐑^M}_{𝐑^M}(_𝐑^{\mathrm{}}a(t,z,x,y\xi )\phi (g(t))dt)\times $$ $$\times \mathrm{exp}\{\sqrt{1}\xi (xy)\}f(z,y)dyd\xi .$$ Remark that the integral $`_𝐑^{\mathrm{}}a(t,z,x,y,\xi )\phi (g(t))𝑑t`$ is just a type of Feynman path integrals. 4. ###### Lemma 5.6. If in every repeated induction step, see \[Li\], \[D2\], $`\mathrm{\Pi }_𝒪(\phi )`$ satisfies the compactness criteria, then the operator $`\mathrm{\Pi }_𝒪(\phi )`$ is trace class and hence $$tr\mathrm{\Pi }_𝒪(\phi )=(2\pi )^M_{𝐑^M}_{𝐑^M}tr(_𝐑^{\mathrm{}}a(t,z,x,x,\xi )\phi (g(t))dtdxd\xi .$$ The proof of the theorem is therefore achieved. ∎ ## 6. Examples In this section, we illustrate the main ideas in examples. ### 6.1. $`\overline{MD}`$-groups We expose in this subsection our joint works with Nguyen Viet Hai \[DH1\]-\[DH2\]. #### 6.1.1. The group of affine transformations of the real straight line We refer the reader to the work \[DH1\] for a detailed exposition with complete proof. Canonical coordinates on the upper half-planes. Recall that the Lie algebra $`𝔤=aff(𝐑)`$ of affine transformations of the real straight line is described as follows, see for example \[D2\]: The Lie group $`Aff(𝐑)`$ of affine transformations of type $$x𝐑ax+b,\text{ for some parameters }a,b𝐑,a0.$$ It is well-known that this group $`Aff(𝐑)`$ is a two dimensional Lie group which is isomorphic to the group of matrices $$Aff(𝐑)\{\left(\begin{array}{cc}a& b\\ 0& 1\end{array}\right)\|a,b𝐑,a0\}.$$ We consider its connected component $$G=Aff_0(𝐑)=\{\left(\begin{array}{cc}a& b\\ 0& 1\end{array}\right)\|a,b𝐑,a>0\}$$ of identity element. Its Lie algebra is $$𝔤=aff(𝐑)\{\left(\begin{array}{cc}\alpha & \beta \\ 0& 0\end{array}\right)\|\alpha ,\beta 𝐑\}$$ admits a basis of two generators $`X,Y`$ with the only nonzero Lie bracket $`[X,Y]=Y`$, i.e. $$𝔤=aff(𝐑)\{\alpha X+\beta Y\|[X,Y]=Y,\alpha ,\beta 𝐑\}.$$ The co-adjoint action of $`G`$ on $`𝔤^{}`$ is given (see e.g. \[AC2\], \[kirillov1\]) by $$K(g)F,Z=F,Ad(g^1)Z,F𝔤^{},gG\text{ and }Z𝔤.$$ Denote the co-adjoint orbit of $`G`$ in $`𝔤`$, passing through $`F`$ by $$\mathrm{\Omega }_F=K(G)F:=\{K(g)F\|FG\}.$$ Because the group $`G=Aff_0(𝐑)`$ is exponential (see \[D2\]), for $`F𝔤^{}=aff(𝐑)^{}`$, we have $$\mathrm{\Omega }_F=\{K(\mathrm{exp}(U)F\|Uaff(𝐑)\}.$$ It is easy to see that $$K(\mathrm{exp}U)F,Z=F,\mathrm{exp}(ad_U)Z.$$ It is easy therefore to see that $$K(\mathrm{exp}U)F=F,\mathrm{exp}(ad_U)XX^{}+F,\mathrm{exp}(ad_U)YY^{}.$$ For a general element $`U=\alpha X+\beta Y𝔤`$, we have $$\mathrm{exp}(ad_U)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\left(\begin{array}{cc}0& 0\\ \beta & \alpha \end{array}\right)^n=\left(\begin{array}{cc}1& 0\\ L& e^\alpha \end{array}\right),$$ where $`L=\alpha +\beta +\frac{\alpha }{\beta }(1e^\beta )`$. This means that $$K(\mathrm{exp}U)F=(\lambda +\mu L)X^{}+(\mu e^\alpha )Y^{}.$$ From this formula one deduces \[D2\] the following description of all co-adjoint orbits of $`G`$ in $`𝔤^{}`$: If $`\mu =0`$, each point $`(x=\lambda ,y=0)`$ on the abscissa ordinate corresponds to a 0-dimensional co-adjoint orbit $$\mathrm{\Omega }_\lambda =\{\lambda X^{}\},\lambda 𝐑.$$ * For $`\mu 0`$, there are two 2-dimensional co-adjoint orbits: the upper half-plane $`\{(\lambda ,\mu )\|\lambda ,\mu 𝐑,\mu >0\}`$ corresponds to the co-adjoint orbit (1) $$\mathrm{\Omega }_+:=\{F=(\lambda +\mu L)X^{}+(\mu e^\alpha )Y^{}\|\mu >0\},$$ and the lower half-plane $`\{(\lambda ,\mu )\|\lambda ,\mu 𝐑,\mu <0\}`$ corresponds to the co-adjoint orbit (2) $$\mathrm{\Omega }_{}:=\{F=(\lambda +\mu L)X^{}+(\mu e^\alpha )Y^{}\|\mu <0\}.$$ We shall work from now on for the fixed co-adjoint orbit $`\mathrm{\Omega }_+`$. The case of the co-adjoint orbit $`\mathrm{\Omega }_{}`$ is similarly treated. First we study the geometry of this orbit and introduce some canonical coordinates in it. It is well-known from the orbit method \[K\] that the Lie algebra $`𝔤=aff(𝐑)`$, realized by the complete right-invariant Hamiltonian vector fields on co-adjoint orbits $`\mathrm{\Omega }_FG_FG`$ with flat (co-adjoint) action of the Lie group $`G=Aff_0(𝐑)`$. On the orbit $`\mathrm{\Omega }_+`$ we choose a fix point $`F=Y^{}`$. It is well-known from the orbit method that we can choose an arbitrary point $`F`$ on $`\mathrm{\Omega }_F`$. It is easy to see that the stabilizer of this (and therefore of any) point is trivial $`G_F=\{e\}`$. We identify therefore $`G`$ with $`G_Y^{}G`$. There is a natural diffeomorphism $`Id_𝐑\times \mathrm{exp}(.)`$ from the standard symplectic space $`𝐑^2`$ with symplectic 2-form $`dpdq`$ in canonical Darboux $`(p,q)`$-coordinates, onto the upper half-plane $`𝐇_+𝐑𝐑_+`$ with coordinates $`(p,e^q)`$, which is, from the above coordinate description, also diffeomorphic to the co-adjoint orbit $`\mathrm{\Omega }_+`$. We can use therefore $`(p,q)`$ as the standard canonical Darboux coordinates in $`\mathrm{\Omega }_Y^{}`$. There are also non-canonical Darboux coordinates $`(x,y)=(p,e^q)`$ on $`\mathrm{\Omega }_Y^{}`$. We show now that in these coordinates $`(x,y)`$, the Kirillov form looks like $`\omega _Y^{}(x,y)=\frac{1}{y}dxdy`$, but in the canonical Darboux coordinates $`(p,q)`$, the Kirillov form is just the standard symplectic form $`dpdq`$. This means that there are symplectomorphisms between the standard symplectic space $`𝐑^2,dpdq)`$, the upper half-plane $`(𝐇_+,\frac{1}{y}dxdy)`$ and the co-adjoint orbit $`(\mathrm{\Omega }_Y^{},\omega _Y^{})`$. Each element $`Z𝔤`$ can be considered as a linear functional $`\stackrel{~}{Z}`$ on co-adjoint orbits, as subsets of $`𝔤^{}`$, $`\stackrel{~}{Z}(F):=F,Z`$. It is well-known that this linear function is just the Hamiltonian function associated with the Hamiltonian vector field $`\xi _Z`$, which represents $`Z𝔤`$ following the formula $$(\xi _Zf)(x):=\frac{d}{dt}f(x\mathrm{exp}(tZ))\|_{t=0},fC^{\mathrm{}}(\mathrm{\Omega }_+).$$ The Kirillov form $`\omega _F`$ is defined by the formula (3) $$\omega _F(\xi _Z,\xi _T)=F,[Z,T],Z,T𝔤=aff(𝐑).$$ This form defines the symplectic structure and the Poisson brackets on the co-adjoint orbit $`\mathrm{\Omega }_+`$. For the derivative along the direction $`\xi _Z`$ and the Poisson bracket we have relation $`\xi _Z(f)=\{\stackrel{~}{Z},f\},fC^{\mathrm{}}(\mathrm{\Omega }_+)`$. It is well-known in differential geometry that the correspondence $`Z\xi _Z,Z𝔤`$ defines a representation of our Lie algebra by vector fields on co-adjoint orbits. If the action of $`G`$ on $`\mathrm{\Omega }_+`$ is flat \[D2\], we have the second Lie algebra homomorphism from strictly Hamiltonian right-invariant vector fields into the Lie algebra of smooth functions on the orbit with respect to the associated Poisson brackets. Denote by $`\psi `$ the indicated symplectomorphism from $`𝐑^2`$ onto $`\mathrm{\Omega }_+`$ $$(p,q)𝐑^2\psi (p,q):=(p,e^q)\mathrm{\Omega }_+$$ ###### Proposition 6.1. 1. Hamiltonian function $`f_Z=\stackrel{~}{Z}`$ in canonical coordinates $`(p,q)`$ of the orbit $`\mathrm{\Omega }_+`$ is of the form $$\stackrel{~}{Z}\psi (p,q)=\alpha p+\beta e^q,\text{ if }Z=\left(\begin{array}{cc}\alpha & \beta \\ 0& 0\end{array}\right).$$ 2. In the canonical coordinates $`(p,q)`$ of the orbit $`\mathrm{\Omega }_+`$, the Kirillov form $`\omega _Y^{}`$ is just the standard form $`\omega =dpdq`$. Computation of generators $`\widehat{\mathrm{}}_Z`$ Let us denote by $`\mathrm{\Lambda }`$ the 2-tensor associated with the Kirillov standard form $`\omega =dpdq`$ in canonical Darboux coordinates. We use also the multi-index notation. Let us consider the well-known Moyal $``$-product of two smooth functions $`u,vC^{\mathrm{}}(𝐑^2)`$, defined by $$uv=u.v+\underset{r1}{}\frac{1}{r!}(\frac{1}{2i})^rP^r(u,v),$$ where $$P^r(u,v):=\mathrm{\Lambda }^{i_1j_1}\mathrm{\Lambda }^{i_2j_2}\mathrm{}\mathrm{\Lambda }^{i_rj_r}_{i_1i_2\mathrm{}i_r}u_{j_1j_2\mathrm{}j_r}v,$$ with $$_{i_1i_2\mathrm{}i_r}:=\frac{^r}{x^{i_1}\mathrm{}x^{i_r}},x:=(p,q)=(p_1,\mathrm{},p_n,q^1,\mathrm{},q^n)$$ as multi-index notation. It is well-known that this series converges in the Schwartz distribution spaces $`𝒮(𝐑^n)`$. We apply this to the special case $`n=1`$. In our case we have only $`x=(x^1,x^2)=(p,q)`$. ###### Proposition 6.2. In the above mentioned canonical Darboux coordinates $`(p,q)`$ on the orbit $`\mathrm{\Omega }_+`$, the Moyl $``$-product satisfies the relation $$i\stackrel{~}{Z}i\stackrel{~}{T}i\stackrel{~}{T}i\stackrel{~}{Z}=i\stackrel{~}{[Z,T]},Z,Taff(𝐑).$$ Consequently, to each adapted chart $`\psi `$ in the sense of \[AC2\], we associate a $`G`$-covariant $``$-product. ###### Proposition 6.3 (see \[G\]). Let $``$ be a formal differentiable $``$-product on $`C^{\mathrm{}}(M,𝐑)`$, which is covariant under $`G`$. Then there exists a representation $`\tau `$ of $`G`$ in $`AutN[[\nu ]]`$ such that $$\tau (g)(uv)=\tau (g)u\tau (g)v.$$ Let us denote by $`_pu`$ the partial Fourier transform \[meisevogt\] of the function $`u`$ from the variable $`p`$ to the variable $`x`$, i.e. $$_p(u)(x,q):=\frac{1}{\sqrt{2\pi }}_𝐑e^{ipx}u(p,q)𝑑p.$$ Let us denote by $`_p^1(u)(x,q)`$ the inverse Fourier transform. ###### Lemma 6.4. 1. $`_p_p^1(p.u)=i_p^1(x.u)`$ , 2. $`_p(v)=i_x_p(v)`$ , 3. $`P^k(\stackrel{~}{Z},_p^1(u))=(1)^k\beta e^q\frac{^k_p^1(u)}{^kp},\text{ with }k2.`$ For each $`Zaff(𝐑)`$, the corresponding Hamiltonian function is $`\stackrel{~}{Z}=\alpha p+\beta e^q`$ and we can consider the operator $`\mathrm{}_Z`$ acting on dense subspace $`L^2(𝐑^2,\frac{dpdq}{2\pi })^{\mathrm{}}`$ of smooth functions by left $``$-multiplication by $`i\stackrel{~}{Z}`$, i.e. $`\mathrm{}_Z(u)=i\stackrel{~}{Z}u`$. It is then continuated to the whole space $`L^2(𝐑^2,\frac{dpdq}{2\pi })`$. It is easy to see that, because of the relation in Proposition (6.2), the correspondence $`Zaff(𝐑)\mathrm{}_Z=i\stackrel{~}{Z}.`$ is a representation of the Lie algebra $`aff(𝐑)`$ on the space $`N[[\frac{i}{2}]]`$ of formal power series in the parameter $`\nu =\frac{i}{2}`$ with coefficients in $`N=C^{\mathrm{}}(M,𝐑)`$, see e.g. \[G\] for more detail. We study now the convergence of the formal power series. In order to do this, we look at the $``$-product of $`i\stackrel{~}{Z}`$ as the $``$-product of symbols and define the differential operators corresponding to $`i\stackrel{~}{Z}`$. It is easy to see that the resulting correspondence is a representation of $`𝔤`$ by pseudo-differential operators. ###### Proposition 6.5. For each $`Zaff(𝐑)`$ and for each compactly supported $`C^{\mathrm{}}`$ function $`uC_0^{\mathrm{}}(𝐑^2)`$, we have $$\widehat{\mathrm{}}_Z(u):=_p\mathrm{}_Z_p^1(u)=\alpha (\frac{1}{2}_q_x)u+i\beta e^{q\frac{x}{2}}u.$$ ###### Remark 6.6. Setting new variables $`s=q\frac{x}{2}`$, $`t=q+\frac{x}{2}`$, we have (4) $$\widehat{\mathrm{}}_Z(u)=\alpha \frac{u}{s}+i\beta e^su,$$ e.i. $$\widehat{\mathrm{}}_Z=\alpha \frac{}{s}+i\beta e^s,$$ which provides a representation of the Lie algebra $`aff(𝐑)`$. The associate irreducible unitary representations Our aim in this section is to exponentiate the obtained representation $`\widehat{\mathrm{}}_Z`$ of the Lie algebra $`aff(𝐑)`$ to the corresponding representation of the Lie group $`Aff_0(𝐑)`$. We shall prove that the result is exactly the irreducible unitary representation $`T_{\mathrm{\Omega }_+}`$ obtained from the orbit method or Mackey small subgroup method applied to this group $`Aff(𝐑)`$. Let us recall first the well-known list of all the irreducible unitary representations of the group of affine transformation of the real straight line. ###### Theorem 6.7 (\[GN\]). Every irreducible unitary representation of the group $`Aff(𝐑)`$ of all the affine transformations of the real straight line, up to unitary equivalence, is equivalent to one of the pairwise nonequivalent representations: * the infinite dimensional representation $`S`$, realized in the space $`L^2(𝐑^{},\frac{dy}{\|y\|})`$, where $`𝐑^{}=𝐑\{0\}`$ and is defined by the formula $$(S(g)f)(y):=e^{iby}f(ay),\text{ where }g=\left(\begin{array}{cc}a& b\\ 0& 1\end{array}\right),$$ * the representation $`U_\lambda ^\epsilon `$, where $`\epsilon =0,1`$, $`\lambda 𝐑`$, realized in the 1-dimensional Hilbert space $`𝐂^1`$ and is given by the formula $$U_\lambda ^\epsilon (g)=\|a\|^{i\lambda }(sgna)^\epsilon .$$ Let us consider now the connected component $`G=Aff_0(𝐑)`$. The irreducible unitary representations can be obtained easily from the orbit method machinery. ###### Theorem 6.8. The representation $`\mathrm{exp}(\widehat{\mathrm{}}_Z)`$ of the group $`G=Aff_0(𝐑)`$ is exactly the irreducible unitary representation $`T_{\mathrm{\Omega }_+}`$ of $`G=Aff_0(𝐑)`$ associated following the orbit method construction, to the orbit $`\mathrm{\Omega }_+`$, which is the upper half-plane $`𝐇𝐑𝐑^{}`$, i. e. + $$(\mathrm{exp}(\widehat{\mathrm{}}_Z)f)(y)=(T_{\mathrm{\Omega }_+}(g)f)(y)=e^{iby}f(ay),fL^2(𝐑^{},\frac{dy}{\|y\|}),$$ where $`g=\mathrm{exp}Z=\left(\begin{array}{cc}a& b\\ 0& 1\end{array}\right).`$ By analogy, we have also ###### Theorem 6.9. The representation $`\mathrm{exp}(\widehat{\mathrm{}}_Z)`$ of the group $`G=Aff_0(𝐑)`$ is exactly the irreducible unitary representation $`T_\mathrm{\Omega }_{}`$ of $`G=Aff_0(𝐑)`$ associated following the orbit method construction, to the orbit $`\mathrm{\Omega }_{}`$, which is the lower half-plane $`𝐇𝐑𝐑^{}`$, i. e. $$(\mathrm{exp}(\widehat{\mathrm{}}_Z)f)(y)=(T_\mathrm{\Omega }_{}(g)f)(y)=e^{iby}f(ay),fL^2(𝐑^{},\frac{dy}{\|y\|}),$$ where $`g=\mathrm{exp}Z=\left(\begin{array}{cc}a& b\\ 0& 1\end{array}\right).`$ ###### Remark 6.10. 1. We have demonstrated how all the irreducible unitary representation of the connected group of affine transformations could be obtained from deformation quantization. It is reasonable to refer to the algebras of functions on co-adjoint orbits with this $``$-product as quantum ones. 2. In a forthcoming work, we shall do the same calculation for the group of affine transformations of the complex straight line $`𝐂`$. This achieves the description of quantum $`\overline{MD}`$ co-adjoint orbits, see \[dndiep\] for definition of $`\overline{MD}`$ Lie algebras. #### 6.1.2. The group of affine transformations of the complex straight line Recall that the Lie algebra $`𝔤=aff(𝐂)`$ of affine transformations of the complex straight line is described as follows, see \[D\]. It is well-known that the group $`Aff(𝐂)`$ is a four (real) dimensional Lie group which is isomorphism to the group of matrices: $$Aff(𝐂)\left\{\left(\begin{array}{cc}a& b\\ 0& 1\end{array}\right)\right\|a,b𝐂,a0\}$$ The most easy method is to consider $`X`$,$`Y`$ as complex generators, $`X=X_1+iX_2`$ and $`Y=Y_1+iY_2`$. Then from the relation $`[X,Y]=Y`$, we get$`[X_1,Y_1][X_2,Y_2]+i([X_1Y_2]+[X_2,Y_1])=Y_1+iY_2`$. This mean that the Lie algebra $`aff(𝐂)`$ is a real 4-dimensional Lie algebra, having 4 generators with the only nonzero Lie brackets: $`[X_1,Y_1][X_2,Y_2]=Y_1`$; $`[X_2,Y_1]+[X_1,Y_2]=Y_2`$ and we can choose another basic noted again by the same letters to have more clear Lie brackets of this Lie algebra: $$[X_1,Y_1]=Y_1;[X_1,Y_2]=Y_2;[X_2,Y_1]=Y_2;[X_2,Y_2]=Y_1$$ ###### Remark 6.11. The exponential map $$\mathrm{exp}:𝐂𝐂^{}:=𝐂\backslash \{0\}$$ giving by $`ze^z`$ is just the covering map and therefore the universal covering of $`C^{}`$ is $`\stackrel{~}{𝐂}^{}𝐂`$. As a consequence one deduces that $$\stackrel{~}{Aff}(𝐂)𝐂𝐂\{(z,w)\|z,w𝐂\}$$ with the following multiplication law: $$(z,w)(z^{^{}},w^{^{}}):=(z+z^{},w+e^zw^{})$$ ###### Remark 6.12. The co-adjoint orbit of $`\stackrel{~}{Aff}(𝐂)`$ in $`𝔤^{}`$ passing through $`F𝔤^{}`$ is denoted by $$\mathrm{\Omega }_F:=K(\stackrel{~}{Aff}(𝐂))F=\{K(g)F\|g\stackrel{~}{Aff}(𝐂)\}$$ Then, (see \[D\]): 1. Each point $`(\alpha ,0,0,\delta )`$ is 0-dimensional co-adjoint orbit $`\mathrm{\Omega }_{(\alpha ,0,0,\delta )}`$ 2. The open set $`\beta ^2+\gamma ^2`$ 0 is the single 4-dimensional co-adjoint orbit $`\mathrm{\Omega }_F=\mathrm{\Omega }_{\beta ^2+\gamma ^20}`$. We shall also use $`\mathrm{\Omega }_F`$ in form $`\mathrm{\Omega }_F𝐂\times 𝐂^{}`$. ###### Remark 6.13. Let us denote: $$𝐇_k=\{w=q_1+iq_2𝐂\|\mathrm{}<q_1<+\mathrm{};2k\pi <q_2<2k\pi +2\pi \};k=0,\pm 1,\mathrm{}$$ $$L=\{\rho e^{i\phi }𝐂\|0<\rho <+\mathrm{};\phi =0\}\text{ and }𝐂_k=𝐂\backslash L$$ is a univalent sheet of the Riemann surface of the complex variable multi-valued analytic function $`Ln(w)`$, ($`k=0,\pm 1,\mathrm{}`$) Then there is a natural diffeomorphism $`w𝐇_ke^w𝐂_k`$ with each $`k=0,\pm 1,\mathrm{}.`$ Now consider the map: $$𝐂\times 𝐂\mathrm{\Omega }_F=𝐂\times 𝐂^{}$$ $$(z,w)(z,e^w),$$ with a fixed $`k𝐙`$. We have a local diffeomorphism $$\phi _k:𝐂\times 𝐇_k𝐂\times 𝐂_k$$ $$(z,w)(z,e^w)$$ This diffeomorphism $`\phi _k`$ will be needed in the all sequel. On $`𝐂\times 𝐇_k`$ we have the natural symplectic form (5) $$\omega =\frac{1}{2}[dzdw+d\overline{z}d\overline{w}],$$ induced from $`𝐂^2`$. Put $`z=p_1+ip_2,w=q_1+iq_2`$ and $`(x^1,x^2,x^3,x^4)=(p_1,q_1,p_2,q_2)𝐑^4`$, then $$\omega =dp_1dq_1dp_2dq_2.$$ The corresponding symplectic matrix of $`\omega `$ is $$=\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right)\text{ and }^1=\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right)$$ We have therefore the Poisson brackets of functions as follows. With each $`f,g𝐂^{\mathrm{}}(\mathrm{\Omega })`$ $$\{f,g\}=^{ij}\frac{f}{x^i}\frac{g}{x^j}=^{12}\frac{f}{p_1}\frac{g}{q_1}+^{21}\frac{f}{q_1}\frac{g}{p_1}+^{34}\frac{f}{p_2}\frac{g}{q_2}+^{43}\frac{f}{q_2}\frac{g}{p_2}=$$ $$=\frac{f}{p_1}\frac{g}{q_1}\frac{f}{q_1}\frac{g}{p_1}\frac{f}{p_2}\frac{g}{q_2}+\frac{f}{q_2}\frac{g}{p_2}=$$ $$=2\left[\frac{f}{z}\frac{g}{w}\frac{f}{w}\frac{g}{z}+\frac{f}{\overline{z}}\frac{g}{\overline{w}}\frac{f}{\overline{w}}\frac{g}{\overline{z}}\right]$$ ###### Proposition 6.14. Fixing the local diffeomorphism $`\phi _k(k𝐙)`$, we have: 1. For any element $`Aaff(𝐂)`$, the corresponding Hamiltonian function $`\stackrel{~}{A}`$ in local coordinates $`(z,w)`$ of the orbit $`\mathrm{\Omega }_F`$ is of the form $$\stackrel{~}{A}\phi _k(z,w)=\frac{1}{2}[\alpha z+\beta e^w+\overline{\alpha }\overline{z}+\overline{\beta }e^{\overline{w}}]$$ 2. In local coordinates $`(z,w)`$ of the orbit $`\mathrm{\Omega }_F`$, the symplectic Kirillov form $`\omega _F`$ is just the standard form (1). Computation of Operators $`\widehat{\mathrm{}}_A^{(k)}`$. ###### Proposition 6.15. With $`A,Baff(𝐂)`$, the Moyal $``$-product satisfies the relation: (6) $$i\stackrel{~}{A}i\stackrel{~}{B}i\stackrel{~}{B}i\stackrel{~}{A}=i[\stackrel{~}{A,B}]$$ For each $`A\text{aff}(𝐂`$), the corresponding Hamiltonian function is $$\stackrel{~}{A}=\frac{1}{2}[\alpha z+\beta e^w+\overline{\alpha }\overline{z}+\overline{\beta }e^{\overline{w}}]$$ and we can consider the operator $`\mathrm{}_A^{(k)}`$ acting on dense subspace $`L^2(𝐑^2\times (𝐑^2)^{},\frac{dp_1dq_1dp_2dq_2}{(2\pi )^2})^{\mathrm{}}`$ of smooth functions by left $``$-multiplication by $`i\stackrel{~}{A}`$, i.e: $`\mathrm{}_A^{(k)}(f)=i\stackrel{~}{A}f`$. Because of the relation in Proposition 3.1, we have ###### Corollary 6.16. (7) $$\mathrm{}_{[A,B]}^{(k)}=\mathrm{}_A^{(k)}\mathrm{}_B^{(k)}\mathrm{}_B^{(k)}\mathrm{}_A^{(k)}:=[\mathrm{}_A^{(k)},\mathrm{}_B^{(k)}]^{}$$ From this it is easy to see that, the correspondence $`Aaff(𝐂)\mathrm{}_A^{(k)}=`$i$`\stackrel{~}{A}`$. is a representation of the Lie algebra $`aff(𝐂`$) on the space N$`\left[[\frac{i}{2}]\right]`$ of formal power series, see \[G\] for more detail. Now, let us denote $`_z`$(f) the partial Fourier transform of the function f from the variable $`z=p_1+ip_2`$ to the variable $`\xi =\xi _1+i\xi _2`$, i.e: $$_z(f)(\xi ,w)=\frac{1}{2\pi }_{R^2}e^{iRe(\xi \overline{z})}f(z,w)𝑑p_1𝑑p_2$$ Let us denote by $$_z^1(f)(\xi ,w)=\frac{1}{2\pi }_{R^2}e^{iRe(\xi \overline{z})}f(\xi ,w)𝑑\xi _1𝑑\xi _2$$ the inverse Fourier transform. ###### Lemma 6.17. Putting $`g=g(z,w)=_z^1(f)(z,w)`$ we have: 1. $$_zg=\frac{i}{2}\overline{\xi }g;_z^rg=(\frac{i}{2}\overline{\xi })^rg,r=2,3,$$ 2. $$_{\overline{z}}g=\frac{i}{2}\xi g;_{\overline{z}}^rg=(\frac{i}{2}\xi )^rg,r=2,3,$$ 3. $$_z(zg)=2i_{\overline{\xi }}_z(g)=2i_{\overline{\xi }}f;_z(\overline{z}g)=2i_\xi _z(g)=2i_\xi f$$ 4. $$_wg=_w(_z^1(f))=_{z}^{}{}_{}{}^{1}(_wf);_{\overline{w}}g=_{\overline{w}}(_z^1(f)=_{z}^{}{}_{}{}^{1}(_{\overline{w}}f)$$ We also need another Lemma which will be used in the sequel. ###### Lemma 6.18. With $`g=_z^1`$$`(f)(`$$`z,w)`$, we have: 1. $$_z(P^0(\stackrel{~}{A},g))=i(\alpha _{\overline{\xi }}+\overline{\alpha }_\xi )f+\frac{1}{2}\beta e^wf+\frac{1}{2}\overline{\beta }e^{\overline{w}}f.$$ 2. $$_z(P^1(\stackrel{~}{A},g))=\overline{\alpha }_{\overline{w}}f+\alpha _wf\overline{\beta }e^{\overline{w}}(\frac{i}{2}\xi )f\beta e^w(\frac{i}{2}\overline{\xi })f.$$ 3. $$_z(P^r(\stackrel{~}{A},g))=(1)^r.2^{r1}[\overline{\beta }e^{\overline{w}}(\frac{i}{2}\xi )^r+\beta e^w(\frac{i}{2}\overline{\xi })^r]fr2.$$ ###### Proposition 6.19. For each $`A=\left(\begin{array}{cc}\alpha & \beta \\ 0& 0\end{array}\right)aff(𝐂)`$ and for each compactly supported $`C^{\mathrm{}}`$-function $`fC_0^{\mathrm{}}(𝐂\times 𝐇_k)`$, we have: (8) $$\mathrm{}_A^{(k)}f:=_z\mathrm{}_A^{(k)}_z^1(f)=[\alpha (\frac{1}{2}_w_{\overline{\xi }})f+\overline{\alpha }(\frac{1}{2}_{\overline{w}}_\xi )f+$$ $$+\frac{i}{2}(\beta e^{w\frac{1}{2}\overline{\xi }}+\overline{\beta }e^{\overline{w}\frac{1}{2}\xi })f]$$ ###### Remark 6.20. Setting new variables u = $`w\frac{1}{2}\overline{\xi }`$;$`v=w+\frac{1}{2}\overline{\xi }`$ we have (9) $$\widehat{\mathrm{}}_A^{(k)}(f)=\alpha \frac{f}{u}+\overline{\alpha }\frac{f}{\overline{u}}+\frac{i}{2}(\beta e^u+\overline{\beta }e^{\overline{u}})f\|_{(u,v)}$$ i.e $`\widehat{\mathrm{}}_A^{(k)}=\alpha \frac{}{u}+\overline{\alpha }\frac{}{\overline{u}}+\frac{i}{2}(\beta e^u+\overline{\beta }e^{\overline{u}})`$,which provides a ( local) representation of the Lie algebra aff(C). The Irreducible Representations of $`\stackrel{~}{Aff}(𝐂)`$. Since $`\widehat{\mathrm{}}_A^{(k)}`$ is a representation of the Lie algebra $`\stackrel{~}{\text{Aff}}(𝐂)`$, we have: $$\mathrm{exp}(\widehat{\mathrm{}}_A^{(k)})=\mathrm{exp}\left(\alpha \frac{}{u}+\overline{\alpha }\frac{}{\overline{u}}+\frac{i}{2}(\beta e^u+\overline{\beta }e^{\overline{u}})\right)$$ is just the corresponding representation of the corresponding connected and simply connected Lie group $`\stackrel{~}{Aff}(𝐂)`$. Let us first recall the well-known list of all the irreducible unitary representations of the group of affine transformation of the complex straight line, see \[D\] for more details. ###### Theorem 6.21. Up to unitary equivalence, every irreducible unitary representation of $`\stackrel{~}{\text{Aff}}(𝐂)`$ is unitarily equivalent to one the following one-to-another nonequivalent irreducible unitary representations: 1. The unitary characters of the group, i.e the one dimensional unitary representation $`U_\lambda ,\lambda 𝐂`$, acting in $`𝐂`$ following the formula $`U_\lambda (z,w)=e^{i\mathrm{}(z\overline{\lambda })},(z,w)\stackrel{~}{Aff}(𝐂),\lambda 𝐂.`$ 2. The infinite dimensional irreducible representations $`T_\theta ,\theta 𝐒^1`$, acting on the Hilbert space $`L^2(𝐑\times 𝐒^1)`$ following the formula: (10) $$[T_\theta (z,w)f](x)=\mathrm{exp}(i(\mathrm{}(wx)+2\pi \theta [\frac{\mathrm{}(x+z)}{2\pi }]))f(xz),$$ Where $`(z,w)\stackrel{~}{Aff}(𝐂)`$ ; $`x𝐑\times 𝐒^1=𝐂\backslash \{0\};fL^2(𝐑\times 𝐒^1);`$ $$xz=Re(x+z)+2\pi i\{\frac{\mathrm{}(x+z)}{2\pi }\}$$ In this section we will prove the following important Theorem which is very interesting for us both in theory and practice. ###### Theorem 6.22. The representation $`\mathrm{exp}(\widehat{\mathrm{}}_A^{(k)})`$ of the group $`\stackrel{~}{Aff}(𝐂)`$ is the irreducible unitary representation $`T_\theta `$ of $`\stackrel{~}{Aff}(𝐂)`$ associated, following the orbit method construction, to the orbit $`\mathrm{\Omega }`$, i.e: $$\mathrm{exp}(\widehat{\mathrm{}}_A^{(k)})f(x)=[T_\theta (\mathrm{exp}A)f](x),$$ where $`fL^2(𝐑\times 𝐒^1);A=\left(\begin{array}{cc}\alpha & \beta \\ 0& 0\end{array}\right)aff(𝐂);\theta 𝐒^1;k=0,\pm 1,`$ ###### Remark 6.23. We say that a real Lie algebra $`𝔤`$ is in the class $`\overline{MD}`$ if every K-orbit is of dimension, equal 0 or dim $`𝔤`$. Further more, one proved that (\[D, Theorem 4.4\]) Up to isomorphism, every Lie algebra of class $`\overline{MD}`$ is one of the following: 1. Commutative Lie algebras. 2. Lie algebra $`aff(𝐑)`$ of affine transformations of the real straight line 3. Lie algebra $`aff(𝐂)`$ of affine transformations of the complex straight line. Thus, by calculation for the group of affine transformations of the real straight line $`Aff(𝐑)`$ in \[DH\] and here for the group affine transformations of the complex straight line $`Aff(𝐂)`$ we obtained a description of the quantum $`\overline{MD}`$ co-adjoint orbits. ### 6.2. $`MD_4`$-groups We refer the reader to the results of Nguyen Viet Hai \[H3\]-\[H4\] for the class of $`MD_4`$-groups (i.e. 4-dimensional solvable Lie groups, all the coadjoint of which are of dimension 0 or maximal). It is interesting that here he obtained the same exact computation for $``$-products and all representations. ### 6.3. $`SO(3)`$ As an typical example of compact Lie group, the author proosed Job A. Nable to consider the case of $`SO(3)`$. We refer the reader to the results of Job Nable \[Na1\]-\[Na3\]. In these examples, it is interesting that the $``$-products, in some how as explained in these papers, involved the Maslov indices and Monodromy Theorem. ## Acknowledgments This work was completed during the stay of the first author as a visiting mathematician at the Department of mathematics, The University of Iowa. The author would like to express the deep and sincere thanks to Professor Tuong Ton-That and his spouse, Dr. Thai-Binh Ton-That for their effective helps and kind attention they provided during the stay in Iowa, and also for a discussion about the PBW Theorem. The deep thanks are also addressed to the organizers of the Seminar on Mathematical Physics, Seminar on Operator Theory in Iowa and the Iowa-Nebraska Functional Analysis Seminar (INFAS), in particular the professors Raul Curto, Palle Jorgensen, Paul Muhly and Tuong Ton-That for the stimulating scientific atmosphere. The author would like to thank the University of Iowa for the hospitality and the scientific support, the Alexander von Humboldt Foundation, Germany, for an effective support.
warning/0003/hep-th0003215.html	ar5iv	text	# Untitled Document hep-th/0003215 CALT-68-2267 CITUSC/00-016 SU-ITP-00-11 Noncommutative Gauge Dynamics From The String Worldsheet Jaume Gomis<sup>1</sup>, Matthew Kleban<sup>2,3</sup>, Thomas Mehen<sup>1</sup>, Mukund Rangamani<sup>1,4</sup>, Stephen Shenker<sup>2</sup> <sup>1</sup>California Institute of Technology, Pasadena, CA 91125 CIT-USC Center for Theoretical Physics <sup>2</sup>Department of Physics, Stanford University, CA 94305-4060 <sup>3</sup>Department of Physics, University of California, Berkeley, CA 94720 <sup>4</sup>Department of Physics, Princeton University, NJ 08544 We show how string theory can be used to reproduce the one-loop two-point photon amplitude in noncommutative $`U(1)`$ gauge theory. Using a simple realization of the gauge theory in bosonic string theory, we extract from a string cylinder computation in the decoupling limit the exact one loop field theory result. The result is obtained entirely from the region of moduli space where massless open strings dominate. Our computation indicates that the unusual IR/UV singularities of noncommutative field theory do not come from closed string modes in any simple way. March 2000 1. Introduction It is striking that field theories on noncommutative spaces are naturally embedded in string theory. The first complete example of this phenomenon was found in toroidal compactifications of Matrix Theory with a nonzero $`B`$ field . Of course the remarkable early construction of membranes by large matrices is very much in this spirit. After additional work was done on extracting noncommutative Yang-Mills theory more directly from open string theory -. In a sense this work culminated in where, among other things, a decoupling limit was carefully formulated in which perturbative open string theory reduced to noncommutative Yang-Mills theory. A large amount of work, partially inspired by these developments, has been done on the perturbative dynamics of noncommutative field theories -. These field theories have a very interesting and unusual perturbative behavior -. The noncommutativity of the underlying space gives rise to a strong mixing between the ultraviolet and the infrared . There are analogs of this IR/UV mixing in string theory which provides one motivation to study these systems. When loop diagrams are evaluated in these theories, large momentum regions of the loop integration lead to terms in the effective action that are infrared divergent and nonanalytic in the noncommutativity parameter $`\theta `$. In conventional field theory, singularities in the low energy effective action usually reflect omission of relevant low energy degrees of freedom, and the low energy description is cured once the the missing degrees of freedom are added. In , it was proposed that at least some of the novel IR/UV divergences of one loop diagrams in the noncommutative theory could be understood as arising from tree level exchange of new degrees of freedom. This phenomenon is analogous to open-closed channel duality in one loop string graphs, where ultraviolet divergences in the open string channel can be interpreted as infrared divergences arising from tree level exchange of massless closed strings. This work is motivated in part by trying to understand this interpretation of the IR/UV singularities in noncommutative field theories. Noncommutative gauge theories can be realized in string theory by taking a low energy decoupling limit of theories on D-branes, in the presence of a constant magnetic field . In this stringy setup one can try to reproduce the noncommutative perturbative expansion from string theory. This embedding confronts the issue of whether extra degrees of freedom – apart from the obvious massless open string modes – are needed to make sense of the low energy effective action. As we will show, there seems to be no need to add any further degrees of freedom. We can account for the entire field theory result by looking at the region of cylinder moduli space which is dominated by massless open strings. The opposite region of moduli space, where massless closed strings dominate, gives a vanishing contribution in the zero-slope limit and therefore seem to decouple in the field theory limit. This seems to indicate that the degrees of freedom proposed in , even though they reproduce the low energy effective action after integrating them out, do not have a natural interpretation as massless closed string states. In section $`2`$ we compute the two-point function of the pure noncommutative $`U(1)`$ gauge theory in 3+1 dimensions at one loop in the background field gauge. The background field gauge is very useful for comparing field theory amplitudes with the zero slope limit of string amplitudes because the effective action obtained in the background field gauge is manifestly gauge invariant, as is the answer obtained in string theory. At one loop, the two point function of the noncommutative gauge theory has terms which contain logarithmic and quadratic infrared divergences which do not appear in conventional gauge theories. The appearance of quadratic infrared divergences is surprising, but nevertheless compatible with gauge invariance. In section $`3`$ we embed this gauge theory in bosonic string theory by considering the low energy limit of the theory on a single D3-brane stuck at an $`R^{22}/Z_2`$ orbifold singularity and in the presence of a magnetic field along the worldvolume directions <sup>1</sup> In this paper we do not discuss the physics of noncommutative field theories with $`\theta ^{0i}0`$. Such a gauge theory can be realized by having a constant electric field along the brane. However, in the decoupling limit the upper bound of the allowed electric field vanishes and the string theory realization is ill-defined.. We compute at one loop the planar and non-planar contributions to the two-point function of photons in the string theory. The field theory answer is reproduced by isolating from the string theory amplitude the contribution coming from the boundary of cylinder moduli space where massless open strings are important. Our result indicates that closed string modes decouple from the low energy noncommutative field theory, just as they decouple from conventional gauge theories realized on branes. We conclude with a discussion of our results in section $`4`$. 2. Gauge Theory Calculation In this section we describe the calculation of the two point function in noncommutative $`U(1)`$ gauge theory in the background field gauge. The action for the U(1) noncommutative Yang-Mills theory in a background metric $`G^{\mu \nu }`$ is given by $$S=\frac{1}{4}d^4x\sqrt{G}G^{\mu \rho }G^{\nu \sigma }F_{\mu \nu }F_{\rho \sigma },$$ where the field strength is $$\begin{array}{cc}\hfill F_{\mu \nu }=_\mu A_\nu _\nu A_\mu ig[A_\mu ,A_\nu ]& \\ \hfill [A_\mu ,A_\nu ]=A_\mu A_\nu A_\nu A_\mu .& \end{array}$$ The action in (2.1) is invariant under the gauge transformation $$\delta _\lambda A_\mu =_\mu \lambda ig[A_\mu ,\lambda ]D_\mu \lambda .$$ The noncommutative star product appearing in (2.1)(2.1) is defined by $$f(x)g(x)=e^{\frac{i}{2}\theta ^{ij}\frac{}{\alpha _i}\frac{}{\beta _j}}f(x+\alpha )g(x+\beta )\|_{\alpha =\beta =0},$$ where the parameter $`\theta ^{ij}`$ is related to the commutator of the coordinates in the noncommutative space: $$[x^\mu ,x^\nu ]=i\theta ^{\mu \nu }.$$ We will quantize the theory in background field gauge . The gauge field $`A_\mu `$ is split into classical and a quantum pieces denoted $`A_\mu `$ and $`Q_\mu `$ respectively. The path integral is performed over the quantum fields while the classical fields are kept fixed. The generating functional for Green’s functions is given by $$Z[J,A]=[dQ]\mathrm{det}[\frac{\delta \mathrm{\Delta }}{\delta \lambda }]\mathrm{Exp}\left[id^4x\sqrt{}G\left(\frac{1}{4}G^{\mu \rho }G^{\nu \sigma }F_{\mu \nu }F_{\rho \sigma }\frac{1}{2\alpha }\mathrm{\Delta }^2+G^{\mu \nu }J_\mu Q_\nu \right)\right],$$ where $`\mathrm{\Delta }`$ is the gauge fixing condition, $`\mathrm{det}[\delta \mathrm{\Delta }/\delta \lambda ]`$ is the Faddeev-Popov determinant and $`J_\mu `$ is an external current coupled only to the quantum fields. In (2.1), $`F_{\mu \nu }`$ is understood to be a function of both $`A_\mu `$ and $`Q_\mu `$. The background field gauge effective action is defined by $$\mathrm{\Gamma }[\overline{Q},A]=W[J,A]d^4x\sqrt{}GG^{\mu \nu }J_\mu \overline{Q}_\nu ,$$ where $$W[J,A]=i\mathrm{ln}Z[J,A]\mathrm{and}\overline{Q}_\mu =\frac{\delta W}{\delta J^\mu }.$$ The background field gauge effective action is invariant under the transformations $$\begin{array}{cc}\hfill \delta A_\mu & =_\mu \lambda ig[A_\mu ,\lambda ]\hfill \\ \hfill \delta \overline{Q}_\mu & =[\lambda ,\overline{Q}_\mu ],\hfill \end{array}$$ so $`\mathrm{\Gamma }[0,A]`$ is a manifestly gauge invariant functional of the classical field $`A_\mu `$. Fig. 1: Feynman rules for Noncommutative $`U(1)`$ in background field gauge. In the calculation of section 2, we use Feynman gauge, $`\alpha =1`$. We can compute $`\mathrm{\Gamma }[0,A]`$ by summing one-particle irreducible Feynman diagrams with classical fields $`A_\mu `$ on external legs and quantum fields $`Q_\mu `$ appearing only in internal lines. The Feynman rules for conventional non-Abelian gauge theory in the background field gauge are given in . The Feynman rules for the noncommutative $`U(1)`$ theory can be obtained from by simply replacing the structure constants $`f_{abc}`$ by $`\mathrm{sin}(\frac{1}{2}p_\mu \theta ^{\mu \nu }k_\nu )`$, where $`p,k`$ are the momenta of two of the gluons entering the vertex. The Feynman rules relevant for the calculations of this paper are shown in fig. 1. The explicit gauge invariance of $`\mathrm{\Gamma }[0,A]`$ simplifies computations in the gauge theory. The bare field strength, when expressed in terms of the renormalized coupling and gauge field, is $$F_{\mu \nu }^{bare}=Z_A^{1/2}(_\mu A_\nu _\nu A_\mu igZ_gZ_A^{1/2}[A_\mu ,A_\nu ]),$$ where the $`Z_g,Z_A`$ are the coupling constant and field renormalizations $$g^{bare}=Z_gg,A_\mu ^{bare}=Z_A^{1/2}A_\mu ^{bare}.$$ Gauge invariance of (2.1) implies that $`Z_A^{1/2}=Z_g`$. This means that the $`\beta `$-function can be computed from $`Z_A`$, which only requires knowledge of the two-point function of the theory. The background field gauge is also very useful for comparing field theory amplitudes with the zero slope limit of string theory amplitudes. Ref. derived the $`\beta `$-function of Yang-Mills theory by calculating the effective action of strings in a background magnetic field. Ref. pointed out a correspondence between loop amplitudes in the background field gauge and loops calculated using string motivated rules. More recently, calculated two, three and four point gauge boson amplitudes in open bosonic string theory. Using a suitable prescription for continuing the amplitudes off-shell, the renormalization constants obtained from the zero slope limit of the string theory amplitudes were observed to respect the background field gauge Ward identities. We will see below that the low energy limit of the string theory amplitudes in our D-brane construction reproduce the field theory loop amplitudes calculated in the background field gauge. Fig. 2: One loop contributions to the two-point function. In fig. 2 we show the one loop diagrams for the two-point function in U(1) noncommutative gauge theory. In ordinary gauge theory the tadpole diagrams give vanishing contributions because $$\frac{d^dp}{p^2}=0$$ in dimensional regularization. However, in the noncommutative theory these graphs are nonvanishing and must be included to obtain a gauge invariant answer. The sum of the diagrams of Fig. 1 is $$\begin{array}{cc}\hfill \mathrm{\Pi }_{\mu \nu }=2g^2\frac{d^dq}{(2\pi )^d}\mathrm{sin}^2(\frac{\stackrel{~}{p}q}{2})\left[\frac{8(p^2G_{\mu \nu }p_\mu p_\nu )}{q^2(p+q)^2}+(d2)\left(\frac{(p+2q)_\mu (p+2q)_\nu }{q^2(p+q)^2}\frac{2G_{\mu \nu }}{q^2}\right)\right],& \end{array}$$ where $`\stackrel{~}{p}^\mu =\theta ^{\mu \nu }p_\nu `$. Using the identity $`\mathrm{sin}^2(x)=\frac{1}{2}(1\mathrm{cos}(2x))`$ the field theory expression separates into two parts, one independent of $`\theta `$ and one with a $`\mathrm{cos}(\stackrel{~}{p}q)`$ in the integrand. The term independent of $`\theta `$ corresponds to the planar diagrams , and gives an expression identical to ordinary Yang-Mills theory with the usual group theory factor $`f_{acd}f_{bcd}=N\delta _{ab}`$ replaced by 2. This piece is divergent in four dimensions and gives a $`1/ϵ`$ pole. From this term one can extract the $`\beta `$-function of the noncommutative theory . The term with the $`\mathrm{cos}(\stackrel{~}{p}q)`$ corresponds to the nonplanar graphs and is ultraviolet finite. To compare with the string theory calculation in section $`3`$ it is best to combine the propagators using Feynman parameters and then do the momentum integral via the method of Schwinger parameters. The contribution of the planar graphs is $$\mathrm{\Pi }_{\mu \nu }^P=\frac{ig^2\mu ^{4d}}{(4\pi )^{d/2}}(p^2G_{\mu \nu }p_\mu p_\nu )_0^{\mathrm{}}𝑑t_0^1𝑑xt^{1d/2}e^{p^2tx(1x)}(8(d2)(12x)^2).$$ The contribution from the nonplanar graphs is $$\mathrm{\Pi }_{\mu \nu }^{NP}=\frac{ig^2}{(4\pi )^2}_0^{\mathrm{}}𝑑t_0^1𝑑xt^1e^{p^2tx(1x)\stackrel{~}{p}^2/4t}\left[(p^2G_{\mu \nu }p_\mu p_\nu )(82(12x)^2)\frac{2}{t^2}\stackrel{~}{p}_\mu \stackrel{~}{p}_\nu \right]$$ where since these diagrams are finite we have set $`d=4`$. The planar diagrams are ultraviolet divergent and to regulate this divergence we take $`d=42ϵ`$. The result of evaluating the integral in (2.1) is $$\begin{array}{cc}\hfill \mathrm{\Pi }_{\mu \nu }^P=\frac{i2g^2}{(4\pi )^2}(p^2G_{\mu \nu }p_\mu p_\nu )\left(\frac{117ϵ}{32ϵ}\right)\left(\frac{p^2}{4\pi \mu ^2}\right)^ϵ\frac{\mathrm{\Gamma }[ϵ]\mathrm{\Gamma }[1ϵ]^2}{\mathrm{\Gamma }[22ϵ]}& \\ \hfill =\frac{ig^2}{(4\pi )^2}\frac{22}{3}(p^2G_{\mu \nu }p_\mu p_\nu )\left(\frac{1}{ϵ}\mathrm{ln}\left(\frac{p^2}{\mu ^2}\right)+\mathrm{}\right),& \end{array}$$ where $`\mathrm{}`$ is a constant. The nonplanar diagrams evaluate to $$\begin{array}{cc}\hfill \mathrm{\Pi }_{\mu \nu }^{NP}& =\frac{ig^2}{(4\pi )^2}_0^1dx[2(p^2G_{\mu \nu }p_\mu p_\nu )(82(12x)^2)K_0(p\stackrel{~}{p}\sqrt{x(1x)})\hfill \\ & 16\frac{\stackrel{~}{p}_\mu \stackrel{~}{p}_\nu }{\stackrel{~}{p}^4}p^2\stackrel{~}{p}^2x(1x)K_2(p\stackrel{~}{p}\sqrt{x(1x)})]\hfill \\ & =\frac{ig^2}{(4\pi )^2}\left[(p^2G_{\mu \nu }p_\mu p_\nu )\left(\frac{22}{3}\mathrm{ln}\left(p^2\stackrel{~}{p}^2\right)+\mathrm{}\right)32\frac{\stackrel{~}{p}_\mu \stackrel{~}{p}_\nu }{\stackrel{~}{p}^4}+\mathrm{}\right]\hfill \end{array},$$ where we have expanded in $`p^2\stackrel{~}{p}^2`$ to lowest order and kept the most infrared singular terms. The nonplanar graphs give rise to $`\mathrm{ln}(p^2\stackrel{~}{p}^2)`$, as first observed in , as well as the correction to the photon polarization tensor of the form $`\stackrel{~}{p}_\mu \stackrel{~}{p}_\nu /\stackrel{~}{p}^4`$, observed in . This last term is interesting since it modifies the photon dispersion relation. We will see in the next section that the zero slope limit of the one loop string theory amplitude exactly reproduces the field theory answer (2.1)(2.1). The Schwinger parameter $`t`$ in the field theory calculation is proportional to the modulus of the string world sheet while the Feynman parameter $`x`$ is related to the separation of the vertex operators on the worldsheet. 3. The String Theory Calculation In this section we reproduce the field theory results of section $`2`$ using string theory. We will find a simple realization of the noncommutative U(1) gauge theory using D3-branes in bosonic string theory and perform a one loop string calculation which will yield, in the massless open string boundary of moduli space, the results of the previous section. A four dimensional pure U(1) gauge theory can be realized by taking the low energy limit of the theory on a single D3-brane of bosonic string theory stuck at an $`R^{22}/Z_2`$ orbifold singularity. This can be accomplished by choosing the action of the orbifold group on the Chan-Paton factors to be represented by either of the two one-dimensional representations of $`Z_2`$. The projection equation projects out the transverse scalars and we are left with only a gauge field. The noncommutative version of the pure U(1) gauge theory can be obtained by applying a constant magnetic field along the D-brane worldvolume. The quantum effective action of this gauge theory is encoded in the string theory effective action. An appropriate truncation of a string loop diagram should provide, in the low energy limit, the field theory answer. We will explicitly verify at the one loop level that the planar and non-planar two-point amplitudes for the photon on the cylinder in a background magnetic field reproduce the corresponding computations in the noncommutative gauge theory in the background field gauge. The correlation function of the photon field vertex operators on the disk shows that the low energy classical action on the brane is a noncommutative gauge theory with the usual replacement of conventional products by $``$-products . Comparison between the string theory calculation and the field theory fixes the normalization of the photon vertex operator to be $$V(z,k)=ig_\mathrm{\Sigma }𝑑se_sXe^{ikX},$$ where $`g`$ is the tree level Yang-Mills coupling constant, $`e_\mu `$ is the polarization of the photon, $`s`$ is the coordinate on the worldsheet boundary $`\mathrm{\Sigma }`$ and indices are contracted with the $`G^{\mu \nu }`$ metric. The worldsheet topology of the one-loop diagram is a cylinder, which we represent in the complex z-plane as a rectangle of width $`\pi `$ and height $`2\pi it`$ – where $`0t\mathrm{}`$ is the modulus of the cylinder – and with the edges $`y=0`$ and $`y=2\pi it`$ identified: $$\begin{array}{cc}\hfill z& =x+iy0x\pi \hfill \\ \hfill y& y+2\pi t0y2\pi t.\hfill \end{array}$$ Open string vertex operators must be inserted on the boundaries of the cylinder at either $`x=0`$ or $`x=\pi `$, with the positions given by $`w=iy`$ and $`w=\pi +iy`$ respectively. The full two-point function is obtained by summing over the planar (two vertex operators on the same boundary) and non-planar (each vertex operator on a different boundary) diagrams. These diagrams are given by $$A=_0^{\mathrm{}}\frac{dt}{2t}_0^{2\pi t}𝑑y_1_0^{2\pi t}𝑑y_2Z(t)<V(y_1,k_1)V(y_2,k_2)>e_1^\mu e_2^\nu \mathrm{\Pi }_{\mu \nu },$$ where one should keep in mind when computing $`<\mathrm{}>`$ if the diagram is planar or non-planar. The different terms in (3.1) are easily understood. The $`1/2t`$ factor arises from explicitly gauge fixing the path integral and dividing by the conformal Killing volume of the cylinder (which allows all vertex operators to be unfixed). $`Z(t)`$ is the open string partition function of the vacuum under consideration, which in our case is that of an open string ending on D3-brane of bosonic string theory stuck at an orbifold singularity and in a background magnetic field. The correlator $`<\mathrm{}>`$ is computed by contracting the fields using the Green’s function on the cylinder with boundary conditions modified by the background magnetic field. The open string partition function is a key ingredient in the measure of the correlation function (3.1). Worldsheet consistency conditions require projecting the open string spectrum onto states invariant under the action of the orbifold group. For our $`Z_2`$ orbifold this is reflected in the partition function $$Z(t)=\text{Tr}(\frac{1+g}{2}e^{2tH_o})=Z_1(t)+Z_g(t),$$ where $`H_o`$ is the open string Hamiltonian and $`g`$ is the $`Z_2`$ generator. The $`Z_2`$ action on the endpoints of the string corresponding to a stuck D3-brane just multiplies (3.1) by unity . It is straightforward to show that $$\begin{array}{cc}\hfill Z_1(t)& =i\text{det}(g+2\pi \alpha ^{}F)\frac{V_{p+1}}{2}(8\pi ^2\alpha ^{}t)^{\frac{p+1}{2}}\eta (it)^{24}\hfill \\ \hfill Z_g(t)& =i\text{det}(g+2\pi \alpha ^{}F)2^{\frac{25p}{2}}\frac{V_{p+1}}{2}(8\pi ^2\alpha ^{}t)^{\frac{p+1}{2}}\vartheta _2(0,it)^{\frac{p25}{2}}\eta (it)^{\frac{273p}{2}},\hfill \end{array}$$ where we have left explicit the dimensionality of the brane (which will turn out to be useful when comparing to the field theory results from section $`2`$). $`F_{\mu \nu }`$ are the components of the background magnetic field. In order to compute the correlation function (3.1) we must solve for the Green’s function of the worldsheet scalars on the cylinder. The background magnetic field along the brane does not modify the equations of motion of the open strings ending on it, but it does change the boundary conditions on the worldsheet fields. Worldsheet coordinates along the brane satisfy the following boundary conditions<sup>2</sup> The coordinates transverse to the brane are projected out by the orbifold quotient.: $$g_{\mu \nu }_nX^\nu +2\pi i\alpha ^{}F_{\mu \nu }_sX^\nu \|_\mathrm{\Sigma }=0.$$ Here $`g_{\mu \nu }`$ is the closed string metric (the metric that appears in the string sigma model action). The operators $`_n`$ and $`_s`$ are derivatives normal and tangential to the worldsheet boundaries $`\mathrm{\Sigma }`$. The correlation function of vertex operators $`<\mathrm{}>`$ on a given worldsheet is computed from the propagators of the worldsheet fields, which can be found by solving the worldsheet wave equation while taking into account the boundary conditions (3.1). On the cylinder, the wave equation to be solved is $$\frac{2}{\alpha ^{}}_w_{\overline{w}}G^{\rho \sigma }(w,w^{})=2\pi \delta ^2(ww^{})g^{\rho \sigma }+\frac{1}{2\pi t}g^{\rho \sigma }.$$ The last term, which is proportional to the inverse area of the cylinder, is included in order to satisfy Gauss’ law and is compatible with the boundary conditions (3.1). The propagator we are interested in should solve (3.1), satisfy (3.1) at both boundaries of the cylinder, and respect the identification $`yy+2\pi t`$ of the cylinder. The solution is given by $$\begin{array}{cc}\hfill 𝒢& {}_{}{}^{\mu \nu }(w,w^{})<X^\mu (w)X^\nu (w^{})>=\alpha ^{}[g^{\mu \nu }\left(\mathrm{log}\right\|\vartheta _1(\frac{ww^{}}{2\pi it}\|\frac{i}{t})\|\mathrm{log}\|\vartheta _1(\frac{w+\overline{w}^{}}{2\pi it}\|\frac{i}{t})\|\hfill \\ \hfill & \frac{\text{Re}^2(w+\overline{w}^{})+\text{Re}^2(ww^{})}{4\pi t})+G^{\mu \nu }\mathrm{log}\|\vartheta _1(\frac{w+\overline{w}^{}}{2\pi it}\|\frac{i}{t})\|^2+\frac{\theta ^{\mu \nu }}{2\pi \alpha ^{}}\mathrm{log}(\frac{\vartheta _1(\frac{w+\overline{w}^{}}{2\pi it}\|\frac{i}{t})}{\vartheta _1^{}(\frac{w+\overline{w}^{}}{2\pi it}\|\frac{i}{t})})],\hfill \end{array}$$ where $`w`$ and $`w^{}`$ are points on the cylinder. The propagator on the cylinder has a similar structure to the propagator on the disk. Here $`G^{\mu \nu }`$ is the open string metric (the metric that defines the dispersion relation for open string fields), and $`\theta ^{\mu \nu }`$ is the noncommutativity parameter which appears in the definition of the $``$-product. They are defined by $$\begin{array}{cc}\hfill G^{\mu \nu }& =\left(\frac{1}{g+2\pi \alpha ^{}F}g\frac{1}{g2\pi \alpha ^{}F}\right)^{\mu \nu }\hfill \\ \hfill \theta ^{\mu \nu }& =(2\pi \alpha ^{})^2\left(\frac{1}{g+2\pi \alpha ^{}F}F\frac{1}{g2\pi \alpha ^{}F}\right)^{\mu \nu }.\hfill \end{array}$$ The noncommutative field theory is obtained by taking the limit $`\alpha ^{}ϵ0`$ and $`gϵ^{1/2}0`$ , with the magnetic field kept constant . For the two-point function of photon vertex operators we only need the propagator for points on the boundaries. The correlator is given by $$\begin{array}{cc}\hfill <& _sX^\mu e^{ikX}(w)_sX^\nu e^{ikX}(w^{})>\hfill \\ \hfill =& \left(_w_w^{}𝒢^{\mu \nu }+k_\rho k_\sigma _w𝒢^{\mu \rho }_w^{}𝒢^{\sigma \nu }\right)e^{k_\rho k_\sigma (𝒢^{\rho \sigma }(w,w^{})\frac{1}{2}𝒢_r^{\rho \sigma }(w,w)\frac{1}{2}𝒢_r^{\rho \sigma }(w^{},w^{}))}\hfill \\ \hfill =& \left(k_\rho k_\sigma _w𝒢^{\rho \sigma }_w^{}𝒢^{\mu \nu }+k_\rho k_\sigma _w𝒢^{\mu \rho }_w^{}𝒢^{\nu \sigma }\right)e^{k_\rho k_\sigma (𝒢^{\rho \sigma }(w,w^{})\frac{1}{2}𝒢_r^{\rho \sigma }(w,w)\frac{1}{2}𝒢_r^{\rho \sigma }(w^{},w^{}))}\hfill \end{array}$$ where we have integrated the first term by parts. $`𝒢_r^{mn}(w,w)`$ is the renormalized propagator, which regulates the divergences in the self-contractions by subtracting the short distance behaviour of the propagator. The proper renormalized propagator for open string vertex operators is given by $$𝒢_r^{\rho \sigma }(w,w^{})=𝒢^{\rho \sigma }(w,w^{})+\alpha ^{}G^{\rho \sigma }(\mathrm{log}\|ww^{}\|^2),$$ where $`w`$ and $`w^{}`$ are points on the same boundary. We will denote the combination in the exponent in (3.1) as $$\stackrel{~}{𝒢}^{\rho \sigma }(w,w^{})𝒢^{\rho \sigma }(w,w^{})\frac{1}{2}𝒢_r^{\rho \sigma }(w,w)\frac{1}{2}𝒢_r^{\rho \sigma }(w^{},w^{})$$ We will now consider in turn the results for the planar and non-planar diagrams. 3.1. Planar Two-point function $`\stackrel{~}{𝒢}^{\rho \sigma }(w,w^{})`$ differs from $`𝒢^{\rho \sigma }(w,w^{})`$ by a term that is independent of the position of the vertex operators. On the $`x=0`$ boundary it is given by<sup>3</sup> Unlike Eq.(3.8) above, this expression is not manifestly periodic. However, it is straightforward to rewrite this expression in a form in which periodicity is manifest. $$\begin{array}{cc}\hfill \stackrel{~}{𝒢}^{\rho \sigma }(w,w^{})& =\alpha ^{}\left[G^{\rho \sigma }\mathrm{log}\left\|2\pi \frac{\vartheta _1(\frac{i(yy^{})}{2\pi },it)}{\vartheta _1^{}(0,it)}\right\|^2\frac{(yy^{})^2}{2\pi t}\right]+\frac{i}{2}\theta ^{\rho \sigma }ϵ(yy^{})\hfill \\ & \alpha ^{}G^{\rho \sigma }\mathrm{\Gamma }(yy^{})+\frac{i}{2}\theta ^{\rho \sigma }ϵ(yy^{}),\hfill \end{array}$$ where $`w=y,w^{}=y^{}`$ (or $`w=\pi +iy,w^{}=\pi +iy^{}`$), $`G^{\rho \sigma }`$ is the open string metric and $`ϵ(x)`$ is $`1`$ for $`x>0`$ and -1 for $`x<0`$. On the $`x=\pi `$ boundary the sign of the term proportional to $`\theta ^{\rho \sigma }`$ changes sign<sup>4</sup> We thank H. Dorn for correspondence on this point.. Note that we have used a $`\vartheta `$-function identity to rewrite the propagator in a form conducive to taking the $`t\mathrm{}`$ limit. Plugging this expression into the correlator (3.1) we see that it has the familiar form of the vacuum polarization diagram of the photon $$<_sX^\mu e^{ikX}(w_1)_sX^\nu e^{ikX}(w_2)>=\alpha ^2\left(k^2G^{\mu \nu }k^\mu k^\nu \right)(_y\mathrm{\Gamma })^2e^{\alpha ^{}k^2\mathrm{\Gamma }},$$ where $`k^2=G^{\rho \sigma }k_\rho k_\sigma `$ and $`y=y_1y_2`$. Combining all the terms in (3.1) one is led to the following expression for the planar two-point function $$\mathrm{\Pi }_{\mu \nu }^P=g^2_0^{\mathrm{}}\frac{dt}{2t}_0^{2\pi t}𝑑y_1_0^{2\pi t}𝑑y_2Z(t)\alpha ^2\left(k^2G_{\mu \nu }k_\mu k_\nu \right)(_y\mathrm{\Gamma })^2e^{\alpha ^{}k^2\mathrm{\Gamma }},$$ where $`Z(t)`$ is given by (3.1) and $`\mathrm{\Gamma }`$ by (3.1). The task at hand is to identify in this string computation the noncommutative field theory result. We have to examine (3.1) in the decoupling limit specified in , which in particular requires taking the $`\alpha ^{}0`$ limit. In this limit we only get contributions from corners of string moduli space. We will now show that we obtain the exact planar field theory answer from the boundary of moduli space of the cylinder which is dominated by massless open strings<sup>5</sup> Note that since we are using bosonic string theory the field theory result is obtained only after removing by hand the divergence caused by the open string tachyon., which comes from the $`t\mathrm{}`$ limit. We therefore need the large $`t`$ expression of the integrand in (3.1). The large $`t`$ expansions of $`Z(t)`$ and of $`_y\mathrm{\Gamma }`$ are given by $$\begin{array}{cc}\hfill Z(t)& V_{p+1}(8\pi ^2\alpha ^{}t)^{\frac{p+1}{2}}\left(e^{2\pi t}+p1+𝒪(e^{2\pi t})\right)\hfill \\ \hfill _y\mathrm{\Gamma }& 12x+2\left(e^{2\pi xt}e^{2\pi xt}e^{2\pi t}\right),\hfill \end{array}$$ where $`x=y/2\pi t`$. Plugging these expressions into (3.1) and tossing out the contribution due to the tachyon, one gets (with $`d=p+1`$ and $`gg\mu ^{4d}`$) $$\mathrm{\Pi }_{\mu \nu }^P=i\frac{g^2\mu ^{4d}}{(4\pi )^{d/2}}\left(k^2G_{\mu \nu }k_\mu k_\nu \right)_0^{\mathrm{}}𝑑t_0^1𝑑xt^{1d/2}\left(8(d2)(12x)^2\right)e^{k^2tx(1x)},$$ which is precisely the field theory answer (2.1). To obtain this result we rescaled $`tt/\alpha ^{}`$ and $`yy/\alpha ^{}`$. In these new variables (with $`\alpha ^{}0`$) any finite value of $`t`$ is in the extreme open string limit of the moduli space. In particular the excited open string corrections in (3.1) become $`𝒪(e^{2\pi t/\alpha ^{}})`$ which vanish in the decoupling limit. This will be discussed further in Section $`4`$. 3.2. Non-Planar Two-point function The nonplanar propagator with $`w=\pi +iy`$ and $`w^{}=y^{}`$ is given by $$\begin{array}{cc}\hfill \stackrel{~}{𝒢}^{\rho \sigma }=& \alpha ^{}[G^{\rho \sigma }\left(\mathrm{log}\right\|2\pi \frac{\vartheta _2(\frac{i(yy^{})}{2\pi },it)}{\vartheta _1^{}(0,it)}\|^2+\frac{\pi }{2t}\frac{(yy^{})^2}{2\pi t})+i\frac{\theta ^{\rho \sigma }}{2\pi \alpha ^{}}\frac{(yy^{})}{t}\hfill \\ & g^{\rho \sigma }\frac{\pi }{2t}].\hfill \end{array}$$ As in the planar diagram computation we expand (3.1) in the large $`t`$ region. The asymptotics of the term proportional to the open string metric $`G^{\rho \sigma }`$ is the same as for the planar diagram. The important differences are in the terms proportional to the closed string metric $`g^{\mu \nu }`$ and the noncommutativity parameter $`\theta ^{\mu \nu }`$. The term proportional to $`\theta ^{\mu \nu }`$ does not contribute to the exponential in (3.1), but it plays a very important role in the derivative terms. Since this is a non-planar diagram for an oriented string, there is an overall factor of $`1`$ since the ends of the string carry opposite charges. The final answer is given by $$\begin{array}{cc}\hfill \mathrm{\Pi }_{\mu \nu }^{NP}& =i\frac{g^2\mu ^{4d}}{(4\pi )^{d/2}}_0^{\mathrm{}}𝑑t_0^1𝑑xt^{1d/2}e^{k^2tx(1x)\stackrel{~}{k}^2/4t}\hfill \\ & \times \left(\left(k^2G_{\mu \nu }k_\mu k_\nu \right)\left(8(d2)(12x)^2\right)\frac{2}{t^2}\stackrel{~}{k}^\mu \stackrel{~}{k}^\nu \right).\hfill \end{array}$$ This expression is identical to that obtained from the nonplanar field theory graphs (2.1). 4. Discussion In the previous sections we have seen how the annulus (cylinder) amplitude of string theory in a background $`F`$ field reproduces, in the decoupling limit of , the planar and nonplanar results of noncommutative gauge theory. Let us discuss this in more detail. Schematically the two-point function on the annulus is given in the open string channel by $$A_0^{\mathrm{}}𝑑tt\underset{I}{}\mathrm{exp}(\mathrm{\Delta }_It).$$ The index $`I`$ labels all open string states and $`\mathrm{\Delta }_I`$, basically the $`L_0`$ eigenvalue, is the mass squared of state $`I`$ plus momentum dependence. In the decoupling limit of $`\alpha ^{}0`$ and hence the string mass scale is sent to infinity. The oscillator contribution to $`\mathrm{\Delta }_I`$ is unaffected by $`F`$ and hence is just $`N_I/\alpha ^{}`$ where $`N_I`$ is the total oscillator occupation number. So in the $`\alpha ^{}0`$ limit all the excited open string states become much heavier than the massless one and hence should decouple from the dynamics. The vanishing of the exponential corrections in (3.1) and (3.1) illustrates this. This is similar to other decoupling limits such as those which show that field theories arise from branes separated by short distances. As pointed out in and further discussed in there are peculiar singularities indicating IR/UV mixing in the nonplanar noncommutative gauge theory results. These results have been interpreted to mean that some closed string residue remains in the field theory, even in the decoupling limit. So it is important to examine if and how decoupling is breaking down here. Generally, the only way decoupling can fail is for the interactions of the decoupled theory to have bad high energy behavior. If loops of massless open string states are UV divergent, then massive open string states will generally be excited, violating decoupling. If the decoupled field theory is divergent but renormalizable then there will be a mild violation of decoupling, but all the effects of the massive string states can be absorbed in a few “renormalized” couplings. For instance, $`g^2\mathrm{ln}(p^2/\mu ^2)`$ in (2.1) becomes, in string theory, $`g^2(\mathrm{ln}(p^2\alpha ^{})+𝒪(1))`$. Now let us examine decoupling in the nonplanar diagrams that display mysterious UV/IR singularities. We can write a caricature of the string amplitude (3.1) by suppressing the $`x`$ integral, all numerical factors, and adding back in the effect of the first excited open string state. $$A_{NP}𝑑tt^{1d/2}e^{k^2t\stackrel{~}{k}^2/4t}(1+e^{2\pi t/\alpha ^{}}+\mathrm{}).$$ Note that the $`\stackrel{~}{k}^2/4t`$ term in the exponential renders the UV region of the modular integration ($`t0`$) completely finite for any nonzero $`\stackrel{~}{k}^2`$. This term is present in the field theoretic amplitude (2.1) and represents the effect of the rapidly oscillating phases in the noncommutative gauge theory interaction vertices. These phases are enough to render the nonplanar amplitude finite for any nonzero $`\stackrel{~}{k}^2`$. The smallest important value of $`t`$ in (4.1) is roughly $`t\stackrel{~}{k}^2`$. The contribution of the first excited state is then $`e^{2\pi \stackrel{~}{k}^2/\alpha ^{}}.`$ In the decoupling limit $`\alpha ^{}0`$ and this contribution vanishes for any nonzero $`\stackrel{~}{k}^2`$. The decoupled field theory amplitude is UV finite so decoupling cannot fail. To further investigate this question let us keep $`\alpha ^{}`$ finite. There are then two regimes to consider. If $`\stackrel{~}{k}^2\alpha ^{}`$ then the excited open string state contribution is negligible and the decoupled field theory result is accurate. If $`\stackrel{~}{k}^2\alpha ^{}`$, however, the small $`t`$ region of the integral may be important. This depends on whether the field theory graph without phase factors is UV divergent. It will be, for instance, if the space-time dimension $`d`$ is large enough. If there is a small $`t`$ UV divergence then all the excited open string states will become important. In this case the correct way to analyze the situation is to use channel duality and rewrite the amplitude in terms of closed string states. At small $`t`$ only the lightest closed string states will contribute, giving a massless propagator $`1/k^2`$ (assuming we drop the closed string tachyon). The region where the closed string description is valid becomes smaller and smaller as we approach the decoupling limit $`\alpha ^{}0`$. In this limit the region of validity shrinks to a set of measure zero. For all finite $`\stackrel{~}{k}^2`$ the decoupled field theory describing only the lightest open string mode is exact. So the complete structure of the mysterious IR/UV singularities is contained in the open string description. Of course one can formally represent the behavior of the lightest open string in the dual closed string channel. But this requires a sum over closed string states of arbitrarily high mass and does not seem very transparent. This is the usual situation in dualities. A regime that has a simple description in one set of variables typically has a complicated description in the dual variables. To illustrate this point consider (4.1) for general $`d`$. The log divergence in $`d=4`$ becomes a $`1/\stackrel{~}{k}^{d4}`$ divergence. The massless closed strings will produce a $`1/k^2`$ behavior for any $`d`$. To produce the open string behavior will require a sum over all the closed string states. There is at least one situation where decoupled field theory results can be reproduced from the lightest closed string state<sup>6</sup> This observation is due to Lenny Susskind.. This is the case where a nonrenormalization theorem exists . The contribution of excited open string states vanishes, typically because they are in long multiplets of an extended supersymmetry algebra. The exact amplitude is given by the lightest open string state, and so this must also agree with the closed string answer. 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warning/0003/math0003027.html	ar5iv	text	# Proper forcing and 𝐿⁢(ℝ) ## 1 Introduction It is well known that under the existence of large cardinals the theory of $`L()`$ —possibly with real parameters— is absolute, and in particular cannot be changed by small forcings. Things may be different if one considers the theory of $`L()`$ with ordinal parameters. By results of Woodin and Shelah this theory can be changed by semiproper forcing, even granted large cardinals. In fact the truth value of the formula $`\varphi (\stackrel{~}{\alpha })=\text{}\stackrel{~}{\alpha }\text{ is equal to }\omega _2\text{}`$ can easily be changed as follows: Start with $`V`$ having, e.g., a supercompact. Then $`L(^V)\varphi [\alpha ]`$ for some $`\alpha \omega _2^V`$, simply because any cardinal of $`V`$ must be a cardinal of $`L(^V)`$. Using the supercompact force to make the SemiProper Forcing Axiom hold in the generic extension. The supercompact of $`V`$ becomes $`\omega _2`$ of $`V[G]`$, and in $`V[G]`$ every set still has a sharp. All this can be done with a semiproper forcing. From SPFA it follows in $`V[G]`$ that the non-stationary ideal on $`\omega _1`$ is saturated (see \[Jec87\]; this result is due to Shelah), and by results of Woodin this (with sharps) implies that $`\omega _2`$ as computed in $`L(^{V[G]})`$ is equal to $`\omega _2^{V[G]}`$ (see \[Woo99\]). Now this is greater than $`\omega _2^V`$, and so certainly greater than $`\alpha `$. Thus $`L(^V)\varphi [\alpha ]`$ while $`L(^{V[G]})\vDash ̸\varphi [\alpha ]`$. This example demonstrates that semiproper forcings can change the theory of $`L()`$ with ordinal parameters greater than or equal to $`\omega _2`$.<sup>1</sup><sup>1</sup>1Ordinal parameters below $`\omega _2`$ can, modulo $`\omega _1`$, be coded by reals. Thus, assuming large cardinals, forcing notions which preserve $`\omega _1`$ cannot change the theory of $`L()`$ with ordinal parameters below $`\omega _2`$. Any attempt to prove the preservation of this theory must therefore be restricted to a class of forcing strictly smaller than semiproper. ###### Theorem 1 (Embedding Theorem) (Under large cardinal assumption $`A_\kappa `$, see below.) Let $`P`$ be a proper forcing notion of size $`\kappa `$, and let $`G`$ be $`P`$-generic/$`V`$. Then there exists an elementary embedding $`j:L(^V)L(^{V[G]})`$ which is the identity on all ordinals. The large cardinal assumption, $`A_\kappa `$, of Theorem 1 is the following: $`(A_\kappa )`$ There exists a class inner model $`M`$ and a countable ordinal $`\delta `$ so that * $`M=L(V_\delta ^M)`$; * $`M\text{}\delta \text{ is the supremum of }\omega \text{ Woodin cardinals;”}`$ and * $`M`$ is uniquely iterable for iteration trees of length $`\kappa ^+`$. Uniquely iterable means basically that in the course of our proof we are free to create iteration trees on $`M`$, without having to worry about the existence of cofinal well-founded branches. More precisely $`M`$ must be iterable, meaning that the good player must win the full iteration game of \[MS94\] on $`M`$ of length $`\kappa ^++1`$. Furthermore the choice of cofinal branches must be unique, in the sense that for any iteration tree on $`M`$ of size $`<\kappa ^+`$ there must be a unique cofinal branch $`b`$ such that the direct limit model $`M_b`$ is itself iterable. The technical assumption $`A_\kappa `$ is weaker than the existence (above $`\kappa `$) of $`\omega `$ Woodin cardinals and a measurable cardinal above them. It is closely connected to the large cardinal strength of $`AD^{L()}`$. Theorem 1 implies in particular that the example given above cannot be carried out with proper (as opposed to semiproper) forcing; the full theory of $`L()`$, with ordinal parameters, cannot be changed by proper forcing. It is a further immediate corollary of the Embedding Theorem that $`\mathrm{HOD}^{L()}`$ cannot be changed by proper forcing. Proper forcings also cannot “code” into $`L()`$ a set of ordinals $`AVL()`$: ###### Theorem 2 (Anti-coding Theorem) (Under large cardinal assumption $`A_\kappa `$). Let $`P`$ be a proper forcing notion of size $`\kappa `$, and let $`G`$ be $`P`$-generic/$`V`$. Suppose that $`A\mathrm{ON}`$ is in $`V`$; then $`AL(^V)AL(^{V[G]})`$. As with the Embedding Theorem, the Anti-Coding Theorem cannot be extended much further. By a result of Woodin it fails for semiproper forcings (provably from large cardinals). Both the Anti-Coding Theorem and the Embedding Theorem do however extend to the class of reasonable forcings — a class slightly bigger than proper. The proofs in this paper apply to reasonable forcings. Our Theorems are similar in flavor to results of Foreman and Magidor \[FM95\], who investigated the possibility of forcing to change the definable continuum — the supremum of all ordinals $`\gamma `$ such that $`\gamma `$ is the order type of some prewellorderings of reals in $`L()`$. (This ordinal is commonly denoted as $`\theta ^{L()}`$). In \[FM95\] it is shown that (granted large cardinals) reasonable forcings cannot change the definable continuum. This result can be obtained also from our Embedding Theorem, since clearly $`\theta ^{L(^V)}=j(\theta ^{L(^V)})=\theta ^{L(^{V[G]})}`$ where $`G`$ is $`P`$-generic/$`V`$ for a reasonable forcing $`P`$, and $`j`$ is the elementary embedding given by Theorem 1. \[FM95\] prove a more general result concerning prewellorderings which are homogeneously Suslin. Using additional work of Woodin’s it is possible to derive the Embedding Theorem from their result. The proof of the Embedding Theorem which we include here is different, and its methods are needed later to obtain the Anti-Coding Theorem. As with many results involving $`L()`$ and large cardinals there are (at least) two alternative routes to proving our Theorems; one which uses stationary tower forcing, and another which uses iteration trees. The latter is presented in this paper while the former can be found in \[NZ98\]. It is interesting that even though iteration trees and stationary tower forcing are technically entirely different there are several similarities between the two approaches. Historically the Embedding and Anti-Coding Theorems were conjectured by the second author, who from a weakly compact Woodin cardinal proved the first for c.c.c forcings and the second for c.c.c. forcing as well as proper forcing notions contained in $`\omega _1`$. Both proofs used the techniques of stationary tower forcing. Those results were presented during the 1996 Set Theory meeting in Luminy, France. The first author subsequently used iteration trees to prove the full Theorems as they appear in the present paper, while the second author strengthened the stationary tower proofs to prove roughly the same results as they appear in \[NZ98\]. The structure of this paper is such that most of the use of large cardinal assumptions is exiled into two “black boxes” (Woodin’s genericity iterations) which are quoted and then used. The proofs relating to these black boxes are due entirely to Hugh Woodin. Readers who are not experts on large cardinals may still be able to follow the proofs of the Embedding and Anti-Coding Theorems if they are willing to accept these black boxes. In Section 2 we present the proof of the Embedding Theorem, and in Section 3 the proof of the Anti-Coding Theorem. The proof in Section 3 uses the techniques of Section 2 as its backbone. The proofs of the black boxes are included in an Appendix to the e-print of this paper at http://arXiv.org. ## 2 The Embedding Theorem We begin now the proof of the Embedding Theorem. Fix a proper forcing notion $`P`$ and a generic $`G`$. Fix $`M`$ which witnesses $`A_\kappa `$. To prove the Theorem we must construct the elementary embedding $`j:L(^V)L(^{V[G]})`$. The requirement that $`j\mathrm{ON}`$ be the identity essentially tells us what $`j`$ is. We must have $`j(x)=x`$ for any $`x^V`$, and since all elements of $`L()`$ are definable from a real and some ordinals this fixes the map $`j`$ completely. Any element of $`L(^V)`$ definable in $`L(^V)`$ from the real $`z`$ and the ordinals $`\alpha _0,\mathrm{}\alpha _k`$ using the formula $`\varphi `$, must be mapped to the element of $`L(^{V[G]})`$ definable from $`z,\alpha _0,\mathrm{},\alpha _k`$ using the same formula $`\varphi `$ in $`L(^{V[G]})`$. All we must prove is that this gives a well-defined elementary embedding $`j`$, and this amounts to showing that For any $`z^V`$, any $`\alpha _0,\mathrm{},\alpha _k\mathrm{ON}`$, and any formula $`\varphi `$, $$\begin{array}{c}L(^V)\varphi [z,\alpha _0,\mathrm{},\alpha _k]\hfill \\ L(^{V[G]})\varphi [z,\alpha _0,\mathrm{},\alpha _k].\hfill \end{array}$$ Fix $`z`$, $`\alpha _0,\mathrm{},\alpha _k`$, and a formula $`\varphi `$. We shall prove the above equivalence using a symmetric collapse. Given a model $`N`$ and some ordinal $`\lambda `$, we consider the Lévy Collapse $`col(\omega ,<\lambda )`$ — the finite support product of the forcings $`col(\omega ,\xi )`$ for $`\xi <\lambda `$. Define the name $`\dot{}_{symm}=_{\xi <\lambda }^{N[\dot{H}col(\omega ,<\xi )]}`$, where $`\dot{H}`$ is a name for the generic object. $`\dot{}_{symm}`$ are the reals in the symmetric collapse up to $`\lambda `$. Those were first investigated by Solovay who used a symmetric collapse to construct a model where all sets of reals are Lebesgue measurable. The important property of the collapse is its homogeneity — any statement about $`\dot{}_{symm}`$ which involves only parameters from $`N`$ is true in the generic extension iff it is forced by the empty condition (see \[Jec78\]). Our strategy is to construct a model $`N`$ and two different generics $`H_1`$ and $`H_2`$ such that 1. $`zN`$; 2. $`H_1`$ and $`H_2`$ are both $`col(\omega ,<\lambda )`$-generic/$`N`$; 3. $`\dot{}_{symm}[H_1]=^V`$; and 4. $`\dot{}_{symm}[H_2]=^{V[G]}`$. This will immediately complete the proof, as $`L(^V)\varphi [z,\alpha _0,\mathrm{},\alpha _k]`$ $`_1`$ $`N[H_1]\text{}L(\dot{}_{symm}[H_1])\varphi [z,\alpha _0,\mathrm{},\alpha _k]\text{}`$ $`_2`$ $`N\text{}^{col(\omega ,<\lambda )}L(\dot{}_{symm})\varphi [z,\alpha _0,\mathrm{},\alpha _k]\text{}`$ $`_3`$ $`N[H_2]\text{}L(\dot{}_{symm}[H_2])\varphi [z,\alpha _0,\mathrm{},\alpha _k]\text{}`$ $`_4`$ $`L(^{V[G]})\varphi [z,\alpha _0,\mathrm{},\alpha _k].`$ The implications $`1`$ and $`4`$ follow from items (3) and (4) above. The implications $`2`$ and $`3`$ follow from the homogeneity of the forcing. We construct $`N`$ as an iterate of the model $`M`$, in $`\omega `$ stages. Each stage will be carried out in $`V`$ while the full construction will exist in $`V^{col(\omega ,)}`$. Our main tool is the following Theorem of Woodin’s (see \[HMW\] or http://www.???.???). ###### Theorem (Woodin’s first genericity iteration) Let $`Q`$ be an $`\omega _1+1`$-iterable inner model and let $`\tau <\eta `$ be countable (in $`V`$) ordinals such that $`Q\text{}\eta \text{ is a Woodin cardinal”}`$. Then there exists a forcing notion $`𝕎_{\tau ,\eta }^QQ`$ of size $`\eta `$, such that for any real $`x`$ it is possible to construct an iteration embedding $`j:Q\stackrel{~}{Q}`$ with the property that * $`x`$ is $`j(𝕎_{\tau ,\eta }^Q)`$-generic/$`\stackrel{~}{Q}`$; * $`j(\eta )`$ is countable in $`V`$, indeed $`j^{\prime \prime }(\omega _1^V)\omega _1^V`$; and * $`crit(j)>\tau `$. Furthermore for any small forcing $`𝕆V_\tau ^Q`$ there exists an $`𝕆`$ name for a forcing notion $`\dot{𝕎}_{\tau ,\eta }^{Q,𝕆}`$ so that for any $`o`$ which is $`𝕆`$-generic/$`Q`$, there exists an iteration embedding $`j:Q\stackrel{~}{Q}`$ satisfying the above except that now $`x`$ is made $`j(\dot{𝕎}_{\tau ,\eta }^{Q,𝕆})[o]`$-generic/$`\stackrel{~}{Q}[o]`$.<sup>2</sup><sup>2</sup>2 This does not follow from the previous part of the Theorem, since it gives a $`j`$ which is an iteration of $`Q`$, and this is more restrictive than being an iteration of $`Q[o]`$. Woodin’s genericity Theorem immediately tells us how to iterate $`M`$ so as to satisfy condition (3) above. Fix some $`g:\omega ^V`$ which is $`col(\omega ,)`$-generic/$`V`$, and so enumerates all the reals of $`V`$. Our plan is to apply Woodin’s Theorem using the $`2i`$-th Woodin cardinal of $`M`$ to make $`g(i)`$ generic over an iterate of $`M`$. (The reason we use only the even Woodin cardinals will become clear presently.) More precisely, working in $`M`$ we let $`\{\delta _i\}_{i\omega }`$ be an increasing sequence of Woodin cardinals with supremum $`\delta `$. Inductively define $`\dot{𝔹}_i`$ to be Woodin’s forcing $`\dot{𝕎}_{\delta _{2i1},\delta _{2i}}^{M,\dot{𝔹}_0\mathrm{}\dot{𝔹}_{i1}}`$ (defined in $`M`$). Let $`𝔹`$ be the finite support iteration of the forcings $`\dot{𝔹}_i`$. Inductively construct an iteration of $`M`$. Begin by letting $`M_0=M`$, and construct models $`M_i`$ and embeddings $`j_i:M_iM_{i+1}`$ so that * $`g(0),\mathrm{},g(i1)`$ is $`j_{0,2i}(\dot{𝔹}_0\mathrm{}\dot{𝔹}_{i1})`$-generic over $`M_{2i}`$, where $`j_{0,2i}`$ is obtained through composition of the $`j_i`$’s. * $`j_{2i}:M_{2i}M_{2i+1}`$ is an embedding to make the real $`g(i)`$ generic for the forcing notion $`j_{0,2i+1}(\dot{𝔹}_i)[g(0),\mathrm{},g(i1)]`$ over the model $`M_{2i+1}[g(0),\mathrm{},g(i1)]`$. (Such an embedding always exists by Woodin’s first genericity iteration.) We take $`j_{2i}`$ to be the identity whenever possible. * For the time being, let $`j_{2i+1}:M_{2i+1}M_{2i+2}`$ be the identity. When iterating $`M`$ we use the unique iteration strategy. Thus by an “iteration embedding of $`M`$” we mean only embeddings obtained through those iteration trees on $`M`$ which choose the unique iterable branch at every limit stage. By our iterability assumption on $`M`$ this guarantees that direct limits of iteration embeddings in $`V`$ are well founded. Let $`M_{\mathrm{}}`$ be the direct limit model of the $`M_i`$-s and let $`j_{i,\mathrm{}}:M_iM_{\mathrm{}}`$ be the direct limit maps. Observe that $`M_{\mathrm{}}`$ is well founded. This is not entirely trivial as the sequence $`j_i`$ does not belong to $`V`$. However if this sequence gave rise to an ill founded direct limit one could use Schoenfield absoluteness to pull the existence of such a “bad” sequence back to $`V`$. Notice further that $`j_{2i,\mathrm{}}`$ has critical point greater than $`j_{0,2i}(\delta _{2i1})`$ so that $`j_{0,\mathrm{}}(\dot{𝔹}_0\mathrm{}\dot{𝔹}_{i1})=j_{0,2i}(\dot{𝔹}_0\mathrm{}\dot{𝔹}_{i1})`$ and $`j_{0,\mathrm{}}(\delta _{2i1})=j_{0,2i}(\delta _{2i1})`$. In particular, $`j_{0,\mathrm{}}(\delta _{2i1})`$ is countable (in $`V`$) and so $`j_{0,\mathrm{}}(\delta )\omega _1^V`$. It is clear from the construction that $`g(0),\mathrm{},g(i1)`$ is $`j_{0,\mathrm{}}(\dot{𝔹}_0\mathrm{}\dot{𝔹}_{i1})`$-generic/$`M_{\mathrm{}}`$ for all $`i`$. Using the fact that $`g`$ is $`col(\omega ,^V)`$-generic/$`V`$ one can verify further that $`g(i)i<\omega `$ is $`j_{0,\mathrm{}}(𝔹)`$-generic over $`M_{\mathrm{}}`$. Specifically, fix any dense set $`D`$ in $`j_{0,\mathrm{}}(𝔹)`$ and assume $`()`$ that the filter given by $`g(i)i<\omega `$ does not intersect $`D`$. Fix some $`n`$ large enough so that $`g(0),\mathrm{},g(n1)`$ forces $`()`$ in $`col(\omega ,)`$, and such that $`D=j_{2n,\mathrm{}}(\overline{D})`$ for some $`\overline{D}M_{2n}`$. Now by condition (a), $`g(0),\mathrm{},g(n1)`$ is $`j_{0,2n}(\dot{𝔹}_0\mathrm{}\dot{𝔹}_{n1})`$-generic/$`M_{2n}`$. Working in $`M_{2n}[g(0),\mathrm{},g(n1)]`$ we can therefore find a condition $`b=\dot{b}_0,\mathrm{},\dot{b}_{k1}\overline{D}`$ (with $`kn`$) such that $`\dot{b}_0,\mathrm{},\dot{b}_{n1}`$ belongs to the $`j_{0,2n}(𝔹_0\mathrm{}𝔹_{n1})`$-generic given by $`g(0),\mathrm{},g(n1)`$. Next let us force over $`M_{2n}`$ with $`j_{0,2n}(\dot{𝔹}_n\mathrm{}\dot{𝔹}_{k1})[g(0),\mathrm{},g(n1)]`$ below the condition $`\dot{b}_n,\mathrm{},\dot{b}_{k1}[g(0),\mathrm{},g(n1)]`$, and obtain reals $`y_n,\mathrm{},y_{k1}`$ such that * $`g(0),\mathrm{},g(n1),y_n,\mathrm{},y_{k1}`$ is $`j_{0,2n}(\dot{𝔹}_0\mathrm{}\dot{𝔹}_{k1})`$-generic/$`M_{2n}`$, and * this generic contains the condition $`b\overline{D}`$. Such reals can be found in $`V`$ since the level of $`M_{2n}`$ involved is countable in $`V`$ (see conditions (i,ii) below). Consider finally the condition $`g(0),\mathrm{},g(n1),y_n,\mathrm{},y_{k1}`$ in the forcing $`col(\omega ,)`$. This condition forces our construction to produce a model $`M_{2k}`$ which is equal to $`M_{2n}`$, and an embedding $`j_{2n,2k}`$ which equals the identity (note our use here of the requirement in (b) that $`j_{2i}`$ be the identity whenever possible). From this together with (G1,G2) it follows easily that $`()`$ is forced to fail, but this is a contradiction. Observe next that the forcing $`𝔹`$ can be replaced by a symmetric collapse. In other words it is possible to find $`HM_{\mathrm{}}[g(i)i<\omega ]`$ which is $`col(\omega ,<j_{0,\mathrm{}}(\delta ))`$-generic/$`M_{\mathrm{}}`$ and so that $`\dot{}_{symm}^M_{\mathrm{}}[H]=\{g(i)i<\omega \}`$. In fact it is well known that in general (for $`\delta `$ a strong limit cardinal) whenever $`𝔸`$ is a direct limit of a regular chain of forcings $`𝔸_i`$, each of size $`<\delta `$, such that each cardinal below $`\delta `$ is collapsed to $`\omega `$ by some $`𝔸_i`$, then $`𝔸`$ is isomorphic to $`col(\omega ,<\delta )`$ in such a way that the symmetric reals are exactly those added by the forcings $`𝔸_i`$. In our case the reals added by the forcings $`𝔹_0\mathrm{}𝔹_i`$ are all in $`V`$, and eventually all reals of $`V`$ are added. Thus we finally have $`^V=\dot{}_{symm}^M_{\mathrm{}}[H]`$. The argument so far is not new. It was first presented by Steel who used it in \[Ste93\] to derive several absoluteness results for $`L()`$, among them the generic absoluteness of the theory of $`L()`$ with real —but not ordinal— parameters. For our purposes however this argument is not sufficient. We have made $`^V`$ the set of reals in the symmetric collapse of an iterate of $`M`$, but we must simultaneously make $`^{V[G]}`$ the set of reals in a different symmetric collapse of the same iterate. For this reason exactly we left ourselves some space during the construction, in the form of the embeddings $`j_{2i+1,2i+2}`$ and the Woodin cardinals $`\delta _{2i+1}`$. Let $`\dot{}_i`$ be Woodin’s forcing $`\dot{𝕎}_{\delta _{2i},\delta _{2i+1}}^{M,\dot{}_0\mathrm{}\dot{}_{i1}}`$ defined in $`M`$, and $``$ their finite support iteration. We will use those to make the reals of $`V[G]`$ generic, as we made the reals of $`V`$ generic. We must however take care not to spoil the part of the construction we have completed — we want to define $`j_{2i+1,2i+2}`$ in a way that still allows us to argue that $`g(i)i<\omega `$ is generic for $`j_{0,\mathrm{}}(𝔹)`$. For that argument to work we needed to know that the reals $`y_n,\mathrm{},y_{k1}`$ could be chosen in $`V`$, and this followed from * $`V_{j_{0,\mathrm{}}(\delta _i)}^M_{\mathrm{}}`$ is an element of $`V`$ for all $`i`$; and * $`j_{0,\mathrm{}}(\delta _i)`$ is countable in $`V`$, for all $`i`$. Either one of (i),(ii) can easily be maintained using Woodin’s first and second (see below) genericity iterations. The difficulty is in maintaining both conditions simultaneously, and it is here that we must make use of our assumption that $`P`$ is proper. ###### Lemma 3 (Assuming $`G`$ is $`P`$–generic/$`V`$ for some proper $`P`$.) Let $`Q=L(V_\delta ^Q)`$ be uniquely iterable, in $`V`$, for trees of size $`\kappa ^++1`$. Assume $`\delta `$ is countable in $`V`$, let $`\tau <\eta <\delta `$ be ordinals such that $`Q\text{}\eta \text{ is a Woodin cardinal”}`$, and consider Woodin’s forcing $`𝕎=𝕎_{\tau ,\eta }^Q`$. Then for any real $`xV[G]`$ it is possible to construct an iteration embedding $`j:Q\stackrel{~}{Q}`$ in $`V`$ with the property that * $`x`$ is $`j(𝕎)`$-generic/$`\stackrel{~}{Q}`$; * $`j(\eta )`$ is countable in $`V`$, indeed $`j^{\prime \prime }(\omega _1^V)\omega _1^V`$; and * $`crit(j)>\tau `$. Furthermore for any small forcing $`𝕆V_\tau ^Q`$, if we let $`\dot{𝕎}=\dot{𝕎}_{\tau ,\eta }^{Q,𝕆}`$ then for any $`oV[G]`$ which is $`𝕆`$-generic/$`Q`$ it is possible to construct an iteration embedding $`j:Q\stackrel{~}{Q}`$ satisfying the above except that now $`x`$ is made $`j(\dot{𝕎})[o]`$-generic over $`\stackrel{~}{Q}[o]`$. It is worthwhile emphasizing the difference between Woodin’s Theorem and Lemma 3. In Lemma 3 we allow $`xV[G]`$ (and also $`oV[G]`$ for the second part), and still obtain an iteration embedding $`j`$ in $`V`$. Fix $`g:\omega ^V`$ and $`h:\omega ^{V[G]}`$ so that the pair $`g,h`$ is $`col(\omega ,^V)\times col(\omega ,^{V[G]})`$-generic/$`V[G]`$. Granted the Lemma we may repeat our construction replacing condition (c) with * $`j_{2i+1}:M_{2i+1}M_{2i+2}`$ is an embedding to make the real $`h(i)`$ generic for the forcing $`j_{0,2i+2}(\dot{}_i)[h(0),\mathrm{},h(i1)]`$ over the model $`M_{2i+2}[h(0),\mathrm{},h(i1)]`$. We take $`j_{2i+1}`$ to be the identity if possible. Otherwise we take the embedding given by Lemma 3. This modified construction produces $`M_{\mathrm{}}`$ and $`j_{i,\mathrm{}}`$ satisfying 1. For all $`n<\omega `$ $`g(0),\mathrm{},g(n1)`$ is $`j_{0,\mathrm{}}(\dot{𝔹}_0\mathrm{}\dot{𝔹}_{n1})`$-generic/$`M_{\mathrm{}}`$; 2. For all $`n<\omega `$ $`h(0),\mathrm{},h(n1)`$ is $`j_{0,\mathrm{}}(\dot{}_0\mathrm{}\dot{}_{n1})`$-generic/$`M_{\mathrm{}}`$; and 3. For $`\xi <j_{0,\mathrm{}}(\delta )`$, $`V_\xi ^M_{\mathrm{}}`$ belongs to $`V`$ and is countable in $`V`$. Condition (3) and the genericity of $`g,h`$ allow us as before to argue that in fact $`g(i)i<\omega `$ is $`j_{0,\mathrm{}}(𝔹)`$-generic/$`M_{\mathrm{}}`$; and $`h(i)i<\omega `$ is $`j_{0,\mathrm{}}()`$-generic/$`M_{\mathrm{}}`$. As before we can now convert the forcings $`𝔹`$ and $``$ into symmetric collapses — finding $`H_1`$ and $`H_2`$ which are $`col(\omega ,<j_{0,\mathrm{}}(\delta ))`$-generic/$`M_{\mathrm{}}`$ so that $`\dot{}_{symm}^M_{\mathrm{}}[H_1]=\{g(i)i<\omega \}`$ and $`\dot{}_{symm}^M_{\mathrm{}}[H_2]=\{h(i)i<\omega \}`$. Letting $`N=M_{\mathrm{}}`$ this completes the proof of Theorem 1, at least if $`z`$ belongs to $`M`$ — but if not, before the beginning of the construction simply iterate $`M`$ to make $`z`$ generic, and then continue to realize $`^V`$ and $`^{V[G]}`$ as the reals of a symmetric collapse over $`N[z]`$. It remains therefore only to prove Lemma 3. We use the following: ###### Theorem (Woodin’s second genericity iteration) Let $`Q`$ be a $`\kappa ^++1`$-iterable inner model, let $`\tau <\eta <\kappa ^+`$ be ordinals such that $`Q\text{}\eta \text{ is a Woodin cardinal”}`$, and let $`𝔸V`$ be any forcing notion of size $`\kappa `$. Let $`𝕎=𝕎_{\tau ,\eta }^Q`$ be Woodin’s forcing of the first genericity iteration, defined in $`Q`$ from $`\tau `$ and $`\eta `$. Then for any $`\dot{x}`$ which is a name for a real in $`V^𝔸`$, it is possible to construct an iteration embedding $`j:Q\stackrel{~}{Q}`$ (in $`V`$) with the property that * For any $`F`$ which is $`𝔸`$-generic/$`V`$, the real $`\dot{x}[F]`$ is $`j(𝕎)`$-generic/$`\stackrel{~}{Q}`$; * $`j(\eta )<(\kappa ^+)^V`$, indeed $`j^{\prime \prime }(\kappa ^+)\kappa ^+`$; and * $`crit(j)>\tau `$. Furthermore For any small forcing $`𝕆V_\tau ^Q`$, if we let $`𝕎=\dot{𝕎}_{\tau ,\eta }^{M,𝕆}`$ then for any $`\dot{o}`$ which is an $`𝔸`$ name for an $`𝕆`$-generic filter/$`Q`$, it is possible to construct an iteration embedding $`j:Q\stackrel{~}{Q}`$ (in $`V`$) satisfying the above except that now $`\dot{x}[F]`$ is made generic over $`\stackrel{~}{Q}[\dot{o}[F]]`$ (for all $`F`$ which are $`𝔸`$-generic/$`V`$). Using Woodin’s second genericity iteration let us prove Lemma 3. Let $`j^{}:QQ^{}`$ be the embedding given by Woodin’s second genericity iteration applied with a name $`\dot{x}`$ for the real $`xV[G]`$. Then $`j^{}V`$, but $`j^{}(\eta )`$ need not be countable. To overcome this: Fix an elementary submodel $`Y`$ of $`V_\lambda `$ for some sufficiently large $`\lambda `$ so that $`P,\dot{x},Q,j^{},Q^{}Y`$ <sup>3</sup><sup>3</sup>3$`Q`$ is a class model of course, but it is coded by a real, and we can throw this real into $`Y`$.; $`Y`$ belongs to $`V`$ and is countable in $`V`$; $`GY`$ is $`P`$-generic/Y; and $`Y[GY]V_\lambda [G]`$. The existence of $`Y`$ follows from the properness of $`P`$. In fact it is enough (by the very definition) to assume that $`P`$ is reasonable. Let $`\overline{Y}`$ be the transitive collapse of $`Y`$ and $`\pi :\overline{Y}Y`$ the inverse collapse embedding. Let $`\overline{G}=\pi _{}^{1}{}_{}{}^{\prime \prime }G`$ and $`\overline{\dot{x}},j,\stackrel{~}{Q}=\pi ^1(\dot{x},j^{}.Q^{})`$. Let $`\overline{x}=\overline{\dot{x}}[\overline{G}]`$. Notice that $`Q`$ is not moved by $`\pi ^1`$, so we have $`j:Q\stackrel{~}{Q}`$. $`\pi `$ induces an embedding from $`\overline{Y}[\overline{G}]`$ onto $`Y[GY]`$ which we also call $`\pi `$. Thus $`\pi :\overline{Y}[\overline{G}]V_\lambda [G]`$ is elementary. By the elementarity of $`\pi `$, $`\overline{x}`$ is $`j(𝕎_{\tau ,\eta }^Q)`$-generic/$`\stackrel{~}{Q}`$. Of course $`\overline{x}`$ is a real and is not moved by $`\pi `$, so $`\overline{x}=\pi (\overline{x})=x`$. Thus the embedding $`j`$ makes $`x`$ generic for Woodin’s forcing. As $`j\overline{Y}`$ it is clear that $`j(\eta )`$ is countable in $`V`$. The reader can now easily check the remaining requirements of Lemma 3. Let us here only verify that $`j`$ is an iteration embedding. This is not obvious — by elementarity $`\overline{Y}\text{}j\text{ is an iteration embedding”}`$, but this does not mean $`j`$ is an iteration embedding in $`V`$. Let $`𝒯`$ be the iteration tree giving rise to $`j`$. We must show that the branches $`𝒯`$ chooses are according to the iteration strategy for $`Q`$ which we have in $`V`$. But $`Q`$ is uniquely iterable, so this strategy chooses at every limit stage the unique branch with iterable direct limit. Thus it is sufficient to show that every model $`Q_\xi `$ on the tree $`𝒯`$ is iterable (in $`V`$). Remember that $`\pi `$ maps $`Q_\xi `$ into a model $`Q_\xi ^{}^{}`$ on the tree $`𝒯^{}`$ which gives rise to $`j^{}`$. $`Q_\xi ^{}^{}`$ is iterable and by \[MS94\] every model which embeds into an iterable model is iterable. Thus $`Q_\xi `$ is iterable and we are done. The second part of Lemma 3 is proved in a similar fashion. Note that since $`Q`$ is countable and $`\dot{o},𝕆Q`$, both $`\dot{o}`$ and $`𝕆`$ are automatically in $`Y`$ and $`\pi (\dot{o})=\dot{o}`$. Thus both $`\dot{x}[G]`$ and $`\dot{o}[G]`$ are not moved by $`\pi :\overline{Y}[\overline{G}]V_\lambda [G]`$. We take $`j^{}`$ to be the iteration from the second part of Woodin’s second genericity Theorem, and immediately by the elementarity of $`\pi `$ can conclude that $`\dot{x}[G]`$ is generic over $`\stackrel{~}{Q}[\dot{o}[G]]`$. $`\mathrm{}`$(Lemma 3, Theorem 1) ## 3 The Anti-Coding Theorem Next let us prove Theorem 2. Fix a set $`A\mathrm{ON}`$ in $`V`$. We must show that $`AL(^V)`$ iff $`AL(^{V[G]})`$. Now the implication from left to right follows immediately from the Embedding Theorem. Assume then that $`AL(^{V[G]})`$. We must show $`AL(^V)`$. As all sets in $`L()`$ are definable from a real and some ordinals, we may fix a name $`\dot{x}`$, ordinals $`\stackrel{}{\alpha }`$ and a formula $`\varphi `$, so that $$\stackrel{ˇ}{A}[G]=A=\{\gamma L(^{V[G]})\varphi [\stackrel{}{\alpha },\dot{x}[G],\gamma ]\}.$$ Without loss of generality we may assume that this is forced by the empty condition in $`P`$. It is convenient to replace $`A_\kappa `$ with the large cardinal assumption $`B_\kappa `$ stated below. It can be seen (using Woodin’s second genericity iteration and some fine structure) that $`B_\kappa `$ follows from $`A_\kappa `$. ($`B_\kappa `$) For any $`K\kappa `$ there exists a class model $`M`$ such that * $`M=L(V_\delta ^M)`$, for some $`\kappa <\delta <(\kappa ^+)^V`$; * $`V_\kappa M`$, and $`KM`$; * $`M\text{}\delta \text{ is the supremum of }\omega \text{ Woodin cardinals;”}`$ and * $`M`$ is uniquely iterable above $`\kappa `$ for trees of length $`(\kappa ^+)^V`$ (i.e., the good player wins the iteration game when the bad player is restricted to playing extenders with critical points above $`\kappa `$). As the forcing $`P`$ has size $`\kappa `$ we may take it to be a subset of $`\kappa `$, and so can fix a model $`M`$ satisfying the conditions of assumption $`B_\kappa `$ with $`P,\dot{x}M`$. Notice that from $`B_\kappa `$ it follows that every subset of $`\kappa `$ has a sharp, and so $`V_\delta ^M`$ has a sharp. We now pass to work in a countable elementary submodel $`YV_\lambda `$ (for $`\lambda `$ sufficiently large) which belongs to $`V`$, and contains all relevant objects (including $`V_\delta ^M`$ and its sharp). Let $`\overline{Y}`$ be the transitive collapse of $`Y`$, and $`\overline{M}`$ the image under the collapse map of $`M`$.<sup>4</sup><sup>4</sup>4Again, $`M`$ is a class model. What we mean is that $`\overline{M}=L(\overline{V_\delta ^M})`$ where $`\overline{V_\delta ^M}`$ is the collapse of $`V_\delta ^M`$. Let $`\pi :\overline{Y}Y`$ be the inverse collapse embedding. Let $`\overline{\dot{x}},\overline{P}`$, and $`\overline{G}`$ be the collapse of $`\dot{x}`$, $`P`$, and $`GY`$. Then $`\overline{\dot{x}}[\overline{G}]=x`$, and by properness (reasonability) of $`P`$ we may assume that $`\overline{G}`$ is $`\overline{P}`$-generic/$`\overline{Y}`$. We will attempt to replace the real $`\overline{\dot{x}}[\overline{G}]`$ in the definition of $`A`$ with a real $`\overline{\dot{x}}[h]`$ for some $`hV`$ which is $`\overline{P}`$-generic/$`\overline{Y}`$. The fact that $`hV`$ will then imply that $`AL(^V)`$. It is simple to find $`hV`$ which is $`\overline{P}`$-generic/$`\overline{Y}`$ (since $`\overline{Y}`$ is countable). The difficulty of course is to do this in such a way that $`\overline{\dot{x}}[\overline{G}]`$ and $`\overline{\dot{x}}[h]`$ still define the same set of ordinals. Let us find $`hV`$ which is $`\overline{P}`$-generic/$`\overline{Y}`$ with the property that for any $`\overline{p}h`$, there exists a condition $`q\pi (\overline{p})`$ which is $`(Y,P)`$-generic. If $`P`$ is proper this can be done trivially (perhaps at the price of modifying $`Y`$). If $`P`$ is only known to be reasonable this is a bit less trivial. Fix in this case some $`q_0G`$ which is $`(Y,P)`$-generic. In $`V[G]`$ there exists an $`h`$ which is $`\overline{P}`$-generic/$`\overline{Y}`$ such that all conditions in $`\pi ^{\prime \prime }(h)`$ are compatible with $`q_0`$ (e.g. take $`h=\overline{G}`$). By absoluteness then such $`h`$ exists in $`V`$, and it is easy to see that any such $`h`$ satisfies our requirement above. Through our choice of $`h`$ we may, for any condition $`\overline{p}h`$, fix in some external generic extension of $`V`$ a filter $`G^{\overline{p}}`$ which is $`P`$-generic/$`V`$; contains the condition $`\pi (\overline{p})`$; and such that $`\overline{G}^{\overline{p}}=\pi _{}^{1}{}_{}{}^{\prime \prime }(G^{\overline{p}}Y)`$ is $`\overline{P}`$-generic over $`\overline{Y}`$ (and hence also over $`\overline{M}`$ <sup>5</sup><sup>5</sup>5We are using here the existence of $`\overline{V_\delta ^M}^{\mathrm{}}`$ inside $`\overline{Y}`$ to see that all subsets of $`\overline{P}`$ in $`\overline{M}`$ belong to $`\overline{Y}`$.). By dovetailing together constructions of the sort used in Section 2 iterate $`\overline{M}`$ to a model $`N`$ so that * For each $`\overline{p}h`$ the reals of $`V[G^{\overline{p}}]`$ can be realized as the symmetric collapse over $`N[\overline{G}^{\overline{p}}]`$; and * The reals of $`V`$ can be realized as the symmetric collapse over $`N[h]`$. Let $`j:\overline{M}N`$ be the iteration embedding, which we construct to have critical point above $`\overline{\kappa }`$, so that $`j(\overline{P})=\overline{P}`$, $`j(\overline{\dot{x}})=\overline{\dot{x}}`$ etc. As in Section 2 $`j`$ is a composition of $`\omega `$ maps, each of which is in $`V`$, and $`j`$ itself exists only in some external model. Let us denote $`j(\overline{\delta })`$ by $`\stackrel{~}{\delta }`$.<sup>6</sup><sup>6</sup>6This is easily seen to be equal to $`\omega _1^V`$. “The symmetric collapse” in (a,b) above refers to the collapse up to $`\stackrel{~}{\delta }`$. ###### Claim 4 Working in $`N[h]`$ let $`y=\overline{\dot{x}}[h]`$ and consider the forcing $`col(\omega ,<\stackrel{~}{\delta })`$. We claim that for any ordinal $`\gamma `$ the following are equivalent: 1. $`\gamma A`$ 2. In the forcing $`col(\omega ,<\stackrel{~}{\delta })`$ over $`N[h]`$, it is forced that “$`L(\dot{}_{symm})\varphi [\stackrel{}{\alpha },y,\gamma ]`$”. Otherwise we may fix some ordinal $`\gamma `$, and a condition $`\overline{p}h`$, such that $`\gamma A`$ say, and nonetheless $`\overline{p}`$ forces in $`\overline{P}`$ that $`L(\dot{}_{symm})\varphi [\stackrel{}{\alpha },\overline{\dot{x}},\gamma ]`$ holds in the symmetric collapse. (Alternatively $`\gamma A`$ and $`\overline{p}`$ forces $`L(\dot{}_{symm})\vDash ̸\varphi [\stackrel{}{\alpha },y,\gamma ]`$, but the proof in this case is similar.) Since $`\overline{p}`$ is an element of $`\overline{G}^{\overline{p}}`$ it follows that over $`N[\overline{G}^{\overline{p}}]`$ it is forced in $`col(\omega ,<\stackrel{~}{\delta })`$ that $`L(\dot{}_{symm})\vDash ̸\varphi [\stackrel{}{\alpha },\overline{\dot{x}}[\overline{G}^{\overline{p}}],\gamma ]`$. But now by (a) we may fix $`H_{\overline{p}}`$ which is $`col(\omega ,<\stackrel{~}{\delta })`$-generic/$`N[\overline{G}^{\overline{p}}]`$ such that $`\dot{}_{symm}[H_{\overline{p}}]=^{V[G^{\overline{p}}]}`$. As furthermore $`\overline{\dot{x}}[\overline{G}^{\overline{p}}]=\dot{x}[G^{\overline{p}}]`$ it follows that $`L(^{V[G^{\overline{p}}]})\vDash ̸\varphi [\stackrel{}{\alpha },\dot{x}[G^{\overline{p}}],\gamma ]`$. But this implies $`\gamma A`$, a contradiction. $`\mathrm{}`$(Claim 4) By (b) we may fix $`H`$, a filter which is $`col(\omega ,<\stackrel{~}{\delta })`$-generic over $`N[h]`$ so that $`\dot{}_{symm}[H]=^V`$. By Claim 4 $`\gamma A`$ iff $`L(\dot{R}_{symm}[H])\varphi [\stackrel{}{\alpha },y,\gamma ]`$, i.e., $`L(^V)\varphi [\stackrel{}{\alpha },y,\gamma ]`$. Thus, $$A=\{\gamma L(^V)\varphi [\stackrel{}{\alpha },y,\gamma ]\}L(^V)$$ completing the proof of the Anti-Coding Theorem. $`\mathrm{}`$(Theorem 2) ## Appendix A Black Boxes We include here a proof of Woodin’s genericity Theorems. The results in this Appendix are due to Hugh Woodin (circa 1987, to be published in \[HMW\]). The reader may easily verify that Woodin’s first genericity iteration is an immediate corollary of the second (taking $`𝔸`$ to be the trivial forcing for adding nothing and $`\kappa =\omega `$), and so we prove here only the second. For the rest of this section $`\eta `$ is assumed to be a Woodin cardinal in $`Q`$. Consider the algebra $`_\eta `$ of all transfinite formulae formed by starting with “$`n\stackrel{~}{x}`$” (for $`n\omega `$) and closing under negation and wellordered disjunctions of length $`<\eta `$. The forcing $`𝕎_{\tau ,\eta }^Q`$ is similar to the Lindenbaum algebra on $`_\eta `$, but rather than simply setting $`\varphi \psi \text{}\varphi \psi \text{}`$ Woodin introduces a set of axioms $`𝒜_\eta `$ and then defines: $`\varphi \psi `$ iff $`𝒜\varphi \psi `$; and $`[\varphi ][\psi ]`$ iff $`𝒜\varphi \psi `$. $`𝕎_{\tau ,\eta }^Q`$ is defined to be the forcing notion consisting of equivalence classes $`[\varphi ]`$ for $`\varphi _\eta `$, ordered by $``$ as above. Before writing down the set of axioms $`𝒜`$ note that with this definition, if $`\varphi `$ is any formula such that $`[\varphi ]=0`$ and $`x`$ any real such that $`x\varphi `$, then there must exist an axiom $`a𝒜`$ such that $`x\neg a`$. Thus any real satisfying the axioms cannot satisfy the $`0`$ condition. The set $`𝒜`$ is defined as follows: For any $`\lambda `$ and $`\rho `$ satisfying $`\tau <\lambda <\rho <\eta `$, any $`\rho `$-strong extender $`EV_\eta ^Q`$ with $`crit(E)=\lambda `$, and any sequence $`\stackrel{}{\varphi }=\{\varphi _\xi \}_{\xi <\lambda }`$ of formulae in $`_\eta V_\lambda ^Q`$, let $`i_E:QUlt(Q,E)`$ be the ultrapower embedding of $`Q`$, and let $`\nu `$ be least such that $`i_E(\stackrel{}{\varphi })_\nu `$ is not in $`V_\rho ^Q`$. (Notice that $`\nu \lambda `$, and certainly a strict inequality is possible.) The following formula is taken to be an axiom: $$a_{\lambda ,\rho ,E,\stackrel{}{\varphi }}=\text{}[_{\xi <\nu }i_E(\stackrel{}{\varphi })_\xi ][_{\xi <\lambda }\varphi _\xi ]\text{.”}$$ (This is a formula in $`_\eta `$, and in fact one which is an element of $`V_{\rho +1}^Q`$.) We denote $`\nu `$ by $`\nu _{\lambda ,\rho ,E,\stackrel{}{\varphi }}`$. It is worthwhile observing that $`i_E(\stackrel{}{\varphi })_\xi =\varphi _\xi `$ for $`\xi <\lambda `$, so that the disjunction $`_{\xi <\nu }i_E(\stackrel{}{\varphi })_\xi `$ is always weaker than (or equal to) the disjunction $`_{\xi <\lambda }\varphi _\xi `$. Woodin then proves the following Claim ###### Claim In $`Q`$, the forcing $`𝕎_{\tau ,\eta }^Q`$ is $`\eta `$-c.c. Proof. Assume for contradiction that the Claim fails and fix an anti-chain $`\{[\psi _\xi ]\}_{\xi <\eta }`$ witnessing this. Let $`f:\eta \eta `$ be the function defined by setting $`f(\xi )`$ to be least $`\alpha `$ such that $`\psi _\xi V_\alpha ^Q`$. Since $`\eta `$ is a Woodin cardinal we can now find $`\lambda <\rho `$ between $`\tau `$ and $`\eta `$ and an extender $`EV_\eta ^Q`$ such that 1. $`crit(E)=\lambda `$; 2. $`E`$ is $`\rho `$ strong, and indeed even $`\rho `$ strong wrt $`\{\psi _\xi \xi <\eta \}`$; and 3. $`\rho >i_E(f)(\lambda )`$. Let $`\stackrel{}{\varphi }=\stackrel{}{\psi }\lambda `$, and consider the axiom $`a_{\lambda ,\rho ,E,\stackrel{}{\varphi }}`$. By condition (3) $`\nu _{\lambda ,\rho ,E,\stackrel{}{\varphi }}\lambda +1`$ so $`a_{\lambda ,\rho ,E,\stackrel{}{\varphi }}`$ clearly proves $`i_E(\stackrel{}{\varphi })_\lambda _{\xi <\lambda }\varphi _\xi `$. But by condition (2) $`i_E(\stackrel{}{\varphi })_\lambda =\psi _\lambda `$ and so $`𝒜\text{}\psi _\lambda _{\xi <\lambda }\psi _\xi \text{}`$. Thus $`[\psi _\lambda ]_{\xi <\lambda }[\psi _\xi ]`$ — a contradiction since $`\{[\psi _\xi ]\}_{\xi <\eta }`$ is an anti-chain. $`\mathrm{}`$ ###### Lemma (Woodin) Let $`x`$ be any real and assume that $`xa`$ for all $`a𝒜`$. Then $`x`$ generates a $`𝕎_{\tau ,\eta }^Q`$-generic filter $`W_x`$, such that $`xQ[W_x]`$. Proof. Define $`W_x=\{[\varphi ]𝕎_{\tau ,\eta }^Qx\varphi \}`$. This is well defined since $`x𝒜`$. To see that $`W_x`$ is a generic filter: Let $`\{[\psi _\xi ]\}_{\xi <\beta }`$ be a maximal anti-chain in $`𝕎_{\tau ,\eta }^Q`$ and assume for contradiction that $`[\psi _\xi ]W_x`$ for all $`\xi <\beta `$. Note that $`\beta <\eta `$ by the previous Claim, so $`\phi =_{\xi <\beta }\neg \psi _\xi `$ is a formula in $`_\eta `$. $`[\phi ]`$ is therefore a condition, and $`[\phi ]=0`$ since $`\{[\psi _\xi ]\}_{\xi <\beta }`$ is a maximal anti-chain. But $`x\phi `$ and this is a contradiction since $`x𝒜`$. Finally to see $`xQ[W_x]`$, note that $`x=\{n[\text{}n\stackrel{~}{x}\text{}]W_x\}`$. $`\mathrm{}`$ At last we are in a position to prove Woodin’s second genericity Theorem. Fix a forcing notion $`𝔸`$ of size $`\kappa `$ and let $`\dot{x}`$ be a name for a real in $`V^𝔸`$. By the previous Lemma, the real $`\dot{x}`$ is generic over $`Q`$ unless it contradicts some of the axioms in $`𝒜`$. The reader can easily verify that if a real $`z`$ contradicts some axiom $`a_{\lambda ,\rho ,E,\stackrel{}{\varphi }}`$, then $`z`$ does not contradict the image axiom $`i_E(a_{\lambda ,\rho ,E,\stackrel{}{\varphi }})`$, where $`i_E:QUlt(Q,E)`$ is the ultrapower map. Thus forming the ultrapower by $`E`$ “removes” the obstruction caused by the axiom $`a_{\lambda ,\rho ,E,\stackrel{}{\varphi }}`$. The second genericity iteration is proved by forming an iteration tree, hitting at every stage the first extender $`E`$ which defines an axiom that $`\dot{x}`$ does not satisfy. A comparison type argument is then used to show that this iteration terminates. The key to this comparison type argument is the fact that once an obstructing axiom has been removed its image will never again become an obstructing axiom. Thus with each step of the construction we come closer to having no obstructing axioms at all. This argument requires an iteration tree; if instead we attempt to use linear iterations then each step may undo previous steps, and the image of an axiom that was handled previously may become obstructing again. Let us begin the construction. We construct a normal iteration tree $`𝒯=E_\alpha ,\rho _\alpha \alpha <\beta `$ with models $`Q_\alpha `$ and tree structure $`T`$. The construction is inductive. At limit $`\alpha `$ we use our iteration strategy for $`Q`$ to pick a cofinal branch of the tree $`𝒯\alpha `$, and set $`Q_\alpha `$ to be the direct limit of the models along this branch. At successor stages $`\alpha +1`$ we must specify $`E_\alpha `$ and $`\rho _\alpha `$ (the tree structure is then determined by finding the least $`\alpha ^{}`$ such that $`\rho _\alpha ^{}>crit(E)`$ and setting $`\alpha ^{}𝑇\alpha +1`$). We shall use only extenders with critical point above $`\tau `$. At successor stages we distinguish between two cases. Case 1: If $`\dot{x}[F]`$ is $`j_{0,\alpha }(𝕎_{\tau ,\eta }^Q)`$-generic/$`Q_\alpha `$ (for all $`𝔸`$-generic/$`V`$ filters $`F`$) then we let $`\beta =\alpha `$, $`j=j_{0,\alpha }`$, and we are done proving the Theorem. Case 2: Otherwise, working in $`V^𝔸`$ we apply the previous Lemma to $`Q_\alpha `$ and $`𝕎_{\tau ,j_{0,\alpha }(\eta )}^{Q_\alpha }=j_{0,\alpha }(𝕎_{\tau ,\eta }^Q)`$, and conclude that there must be some axiom $`aj_{0,\alpha }(𝒜)`$ such that $`\dot{x}\vDash ̸a`$. This axiom must have the form $`a_{\lambda ,\rho ,E,\stackrel{}{\varphi }}^{Q_\alpha }`$ for some $`\lambda ,\rho ,E,\stackrel{}{\varphi }Q_\alpha `$. Let us pick a condition $`q_\alpha 𝔸`$ forcing this, and forcing value for the unsatisfied axiom $`a`$, say $`q`$ forces $`\dot{x}\vDash ̸a_{\lambda _\alpha ,\rho _\alpha ,E_\alpha ,\stackrel{}{\varphi }^\alpha }`$. Pick $`q_\alpha `$ so that $`\rho _\alpha `$ is minimal. We extend the tree by setting $`Q_{\alpha +1}=Ult(Q_\alpha ^{},E_\alpha )`$ for $`\alpha ^{}`$ least so that $`critE_\alpha =\lambda _\alpha <\rho _\alpha ^{}`$. The genericity iteration Theorem will be proved by showing that the second case in the construction cannot hold for all $`\alpha <(\kappa ^+)^V`$. This is very similar to the usual proof that comparisons of mice of size $`\kappa `$ must terminate before reaching $`\kappa ^+`$. Assume for contradiction that the construction continues to $`(\kappa ^+)^V`$, and let $`𝒯`$ be the tree of length $`(\kappa ^+)^V`$ constructed. Since $`Q`$ is assumed to be $`(\kappa ^+)^V+1`$-iterable there exists a cofinal branch through the tree. Let $`b`$ denote this branch. Note that $`b(\kappa ^+)^V`$ is closed-unbounded. For every $`\alpha b`$ let $`\alpha _b^+`$ be the least ordinal such that $`\alpha 𝑇\alpha _b^++1`$. Then $`E_{\alpha _b^+}`$ has critical point ($`\lambda _{\alpha _b^+}`$) below $`\rho _\alpha `$, and is applied to $`Q_\alpha `$ in the tree $`𝒯`$ to form the ultrapower $`Q_{\alpha _b^++1}`$. Note that $`\stackrel{}{\varphi }^{\alpha _b^+}`$ is in $`V_{\lambda _{\alpha _b^+}+1}^{Q_{\alpha _b^+}}`$, and since $`Q_\alpha `$ and $`Q_{\alpha _b^+}`$ agree on subsets of $`\lambda _{\alpha _b^+}`$ it follows that $`\stackrel{}{\varphi }^{\alpha _b^+}Q_\alpha `$. Let us denote $`\stackrel{}{\varphi }^{\alpha _b^+}`$ by $`\stackrel{}{\psi }^\alpha `$. Let $`S_1b`$ be the set of limit points of $`b`$. For $`\alpha S_1`$ the model $`Q_\alpha `$ is a direct limit and so $`Q_\alpha =_{\beta <\alpha ,\beta b}j_{\beta ,\alpha }^{\prime \prime }Q_\beta `$. As $`\stackrel{}{\psi }^\alpha Q_\alpha `$ there must exist some $`h(\alpha )<\alpha `$ such that $`\stackrel{}{\psi }^\alpha j_{h(\alpha ),\alpha }^{\prime \prime }Q_{h(\alpha )}`$. A standard pressing down argument now produces $`\beta <\kappa ^+`$ and stationary $`S_2S_1`$ so that $`h(\alpha )=\beta `$ for all $`\alpha S_2`$. Since $`Q_\beta `$ has cardinality $`\kappa `$, further thinning of $`S_2`$ produces stationary $`S_3S_2`$ and a fixed $`\phi Q_\beta `$ such that $`\stackrel{}{\psi }^\alpha =j_{\beta ,\alpha }(\stackrel{}{\phi })`$ for all $`\alpha S_3`$. Since $`𝔸`$ too has cardinality $`\kappa `$ we may assume further that for some fixed $`q𝔸`$ we have $`q_{\alpha _b^+}=q`$ for all $`\alpha S_3`$. Let $`\alpha `$ be any element of $`S_3`$, and let $`\gamma `$ be $`\alpha _b^+`$ (so $`\gamma +1b`$). Now $`q`$ forces the real $`\dot{x}`$ to contradict the axiom $`a_{\lambda _\gamma ,\rho _\gamma ,E_\gamma ,j_{\beta ,\alpha }(\stackrel{}{\phi )}}^{Q_\gamma }`$. This means that necessarily ($`q`$ forces) $`\dot{x}\vDash ̸_{\xi <\lambda _\gamma }j_{\beta ,\alpha }(\stackrel{}{\phi })_\xi `$ $`()`$, and $`\dot{x}_{\xi <\nu _\gamma }i_{E_\gamma }^{Q_\gamma }(j_{\beta ,\alpha }(\stackrel{}{\phi }))_\xi `$. But $`i_{E_\gamma }^{Q_\gamma }(j_{\beta ,\alpha }(\stackrel{}{\phi }))=i_{E_\gamma }^{Q_\alpha }(j_{\beta ,\alpha }(\stackrel{}{\phi }))`$ <sup>7</sup><sup>7</sup>7We replaced $`Q_\gamma `$ with $`Q_\alpha `$. since $`j_{\beta ,\alpha }(\stackrel{}{\phi })`$ is an element of $`V_{\lambda _\gamma +1}^{Q_\alpha }`$. Thus $`\dot{x}_{\xi <\nu _\gamma }i_{E_\gamma }^{Q_\alpha }(j_{\beta ,\alpha }(\stackrel{}{\phi }))`$. $`i_{E_\gamma }^{Q_\alpha }`$ is simply $`j_{\alpha ,\gamma +1}`$, so we can rewrite the above as ($`q`$ forces) $`\dot{x}_{\xi <\nu _\gamma }j_{\beta ,\gamma +1}(\stackrel{}{\phi })_\xi `$. Consider now any $`\alpha ^{}b`$ such that $`\alpha ^{}>\gamma +1`$. Then $`crit(j_{\gamma +1,\alpha ^{}})\rho _\gamma `$ (it is to secure this fact that we are forced to use iteration trees, and cannot manage with the simpler linear iterations), and so for $`\xi <\nu _\gamma `$, $`j_{\beta ,\gamma +1}(\stackrel{}{\phi })_\xi `$ is not moved by $`j_{\gamma +1,\alpha ^{}}`$. Thus ($`q`$ forces) $`\dot{x}_{\xi <\nu _\gamma }j_{\beta ,\alpha ^{}}(\stackrel{}{\phi })`$. But then clearly $`\dot{x}_{\xi <\lambda _{\alpha _{}^{}{}_{b}{}^{+}}}j_{\beta ,\alpha ^{}}(\stackrel{}{\phi })`$, and we now obtain a contradiction to $`()`$ by taking $`\alpha ^{}S_3`$. This concludes the proof of the first part of the second genericity Theorem. We leave the second half to the reader, and indicate here only how to define the forcing $`\dot{𝕎}_{\tau ,\eta }^{Q,𝕆}`$ when $`𝕆`$ is a forcing notion in $`V_\tau ^Q`$. Working in $`Q^𝕆`$, again consider the algebra of all formulae obtained from “$`n\stackrel{~}{x}`$” closing under negations and wellordered disjunctions (in $`Q^𝕆`$) of length $`<\eta `$. Let $`\dot{}`$ be the set of axioms (computed in $`Q^𝕆`$) $`a_{\stackrel{ˇ}{\lambda },\stackrel{ˇ}{\rho },\dot{E},\dot{\stackrel{}{\varphi }}}`$ as before, with the restriction that $`\dot{E}`$ must be an extender (of $`Q^𝕆`$) induced by an extender of $`Q`$. I.e., there must exist an extender $`FQ`$ such that the embedding $`i_{\dot{E}}^{Q^𝕆}:Q^𝕆Ult(Q^𝕆,\dot{E})`$ extends the embedding $`i_F^Q:QUlt(Q,F)`$. Set then $`\dot{\varphi }\dot{}\dot{\psi }`$ iff $`\dot{}\dot{\varphi }\dot{\psi }`$; and $`\dot{\varphi }\dot{}\dot{\psi }`$ iff $`\dot{}\dot{\varphi }\dot{\psi }`$. $`\dot{𝕎}_{\tau ,\eta }^{Q,𝕆}`$ is then defined to be the set of equivalence classes of $`\dot{}`$, ordered by $`\dot{}`$. The proof of the genericity Theorem proceeds as before. The reader can verify this, noting that there are many extenders $`\dot{E}`$ in $`Q^𝕆`$ which are induced by extenders in $`Q`$ — in fact there are enough such extenders to witness that $`\eta `$ is a Woodin cardinal (because $`𝕆`$ is a “small” forcing). This allows carrying out the argument of the Claim above, and subsequently the rest of the proof.
warning/0003/cond-mat0003234.html	ar5iv	text	# Role of free energy landscape in the dynamics of mean field glassy systems ## I Introduction Generally in the study of thermodynamics much attention is payed to the free energy landscape . This landscape, which can be interpreted as the effective potential whose minima represent different possible states, gives an intuitive and quantitative description of the equilibrium properties. Consider for example ferromagnetic systems. In this case the effective potential is a function of magnetisation. The ferromagnetic transition corresponds to the splitting of the paramagnetic minimum in the two ferromagnetic minima. At low temperature a vanishing external magnetic field breaks the up-down symmetry and fixes the system in one of the two possible ferromagnetic states. Generally, glassy systems are characterised by a complicated energy landscape, which can give rise eventually to the existence of many possible states. Frequently, qualitative explanations of glassy behaviours are based on some assumptions on the properties of the free energy landscape . Consider for example the Kirkpatrick-Thirumalai-Wolynes scenario for the glass transition in which the (exponential) number of states with a given free energy plays a crucial role. However, if the relevance of the free energy landscape for the equilibrium properties is clear, the relationship between the free energy landscape and the dynamical behaviour is not completely understood, especially for glassy systems which remain out of equilibrium also at long time. For instance what we can learn on the (out of equilibrium) dynamical behaviour starting from the knowledge of the free energy landscape is not clear. All the explanations based on the free energy landscape remain often at a qualitative level, because in general this landscape cannot be computed and studied exactly. Only for mean field frustrated systems this “Landscape Paradigm” has received a firm theoretical basis. In this case an analytic solution of the thermodynamics and of the asymptotic out of equilibrium dynamics is available. For these systems the free energy landscape can be computed and the partition function, and therefore the equilibrium properties, can be recovered as a sum over the free energy minima weighted with the Boltzmann factor . This approach to the thermodynamics of mean field spin glasses is called the TAP approach, because it was introduced by Thouless, Anderson and Palmer for the Sherrington-Kirkpatrick model . In this paper we focus on the class of spin glass models which reproduce phenomenologically some features of structural glasses . To understand the relationship between free energy landscape and dynamical behaviour we generalise the TAP approach to dynamics. ## II The TAP approach ### A Static TAP equations In the following we show how the free energy landscape, also called TAP free energy, can be derived for the p-spin spherical model . The aim of this section is to present in a simple case the strategy which we have followed to compute the dynamical TAP equations. The p-spin Hamiltonian reads: $$H(\{S_i\})=\underset{1i_1<\mathrm{}<i_pN}{}J_{i_1,\mathrm{},i_p}S_{i_1}\mathrm{}S_{i_p},$$ (1) where the couplings are Gaussian variables with zero mean and average $`\overline{J_{i_1,\mathrm{},i_p}^2}=\frac{p!}{2N^{p1}}`$. The TAP free energy $`\mathrm{\Gamma }(\beta ,m_i,l)`$, which depends on the magnetisation $`m_i`$ at each site $`i`$ and on the spherical parameter $`l`$, is the Legendre transform of the “true” free energy: $$\beta \mathrm{\Gamma }(\beta ,m_i,l)=\mathrm{ln}_{\mathrm{}}^+\mathrm{}\underset{i=1}{\overset{N}{}}dS_i\mathrm{exp}\left(\beta H(\{S_i\})\underset{i}{}h_i(S_im_i)\frac{\lambda }{2}\underset{i=1}{\overset{N}{}}(S_i^2l)\right).$$ (2) The Lagrange multipliers $`h_i(\beta )`$ fix the magnetisation at each site $`i`$: $`S_i=m_i`$ and $`\lambda (\beta )`$ enforces the condition $`_{i=1}^NS_i^2l=0`$. $``$ denotes the thermal average and $`N`$ is the number of spins. Once $`\mathrm{\Gamma }`$ is known, the equation $`\frac{2}{N}\frac{\beta \mathrm{\Gamma }}{l}\|_{l=1}=\lambda `$ fixes the spherical constraint ($`_iS_i^2=N`$) and gives the spherical multiplier as a function of $`m_i`$, whereas $`\frac{\beta \mathrm{\Gamma }}{m_i}\|_{l=1}=h_i`$ are the TAP equations, which fix the values of local magnetisations. The standard perturbation expansion for the generalised potential $`\mathrm{\Gamma }`$ is rather involved and cannot be directly applied to the Ising case. Thus, we prefer to follow the approach developed for the Sherrington-Kirkpatrick model by T. Plefka and A. Georges and S. Yedidia because it is simple and can be directly applied to all mean field spin glass models. They obtained the TAP free energy for the Sherrington-Kirkpatrick model expanding $`\beta \mathrm{\Gamma }`$ in powers of $`\beta `$ around $`\beta =0`$. For a general system this corresponds to a $`\frac{1}{d}`$ expansion ($`d`$ being the spatial dimension) around mean field theory ; so it is not surprising that for mean field spin glass models only a finite number of terms survives. The zeroth- and first-order terms give the “naïve” TAP free energy, whereas the second term is the Onsager reaction term. From the definition of $`\beta \mathrm{\Gamma }`$ given in equation (2), we find<sup></sup><sup></sup>We are neglecting a useless additive constant in $`\mathrm{\Gamma }`$. A term in $`\mathrm{\Gamma }`$, that does not depend on $`l`$ and $`m_i`$, has no influence on thermodynamics. that the zeroth-order term is the entropy of non interacting spherical spins constrained to have magnetisation $`m_i`$: $$\beta \mathrm{\Gamma }(\beta ,m_i,l)\|_{\beta =0}=\frac{N}{2}\mathrm{ln}\left(l\frac{1}{N}\underset{i=1}{\overset{N}{}}m_i^2\right).$$ (3) Using the Lagrange conditions and that the spins are decoupled at $`\beta =0`$ we find that the linear term in the power expansion of the TAP free energy equals: $$\beta \frac{(\beta \mathrm{\Gamma })}{\beta }\|_{\beta =0}=\beta \underset{1i_1<\mathrm{}<i_pN}{}J_{i_1,\mathrm{},i_p}m_{i_1}\mathrm{}m_{i_p}.$$ (4) This “mean field” energy together with the zeroth-order term gives the standard mean field theory, which becomes exact for infinite-ranged ferromagnetic system. The Onsager reaction term comes from the second derivative of $`\mathrm{\Gamma }`$: $$\frac{\beta ^2}{2}\frac{^2(\beta \mathrm{\Gamma })}{\beta ^2}\|_{\beta =0}=\frac{\beta ^2}{2}\left(\underset{1i_1<\mathrm{}<i_pN}{}Y_{i_1,\mathrm{},i_p}\right)^2_{\beta =0}^c,$$ (5) $`Y_{i_1,\mathrm{},i_p}=J_{i_1,\mathrm{},i_p}S_{i_1}\mathrm{}S_{i_p}(S_{i_1}m_{i_1})m_{i_2}\mathrm{}m_{i_p}`$ $`\mathrm{}m_{i_1}\mathrm{}m_{i_{p1}}(S_{i_p}m_{i_p}).`$ To compute (5) we have used the following Maxwell relations: $`{\displaystyle \frac{h_i}{\beta }}\|_{\beta =0}`$ $`=`$ $`{\displaystyle \frac{}{m_i}}{\displaystyle \frac{(\beta \mathrm{\Gamma })}{\beta }}\|_{\beta =0}`$ (6) $`{\displaystyle \frac{\lambda }{\beta }}\|_{\beta =0}`$ $`=`$ $`{\displaystyle \frac{2}{N}}{\displaystyle \frac{}{l}}{\displaystyle \frac{(\beta \mathrm{\Gamma })}{\beta }}\|_{\beta =0}`$ (7) Using the statistical properties of the couplings it is easy to check that the only terms giving a contribution of the order of $`N`$ correspond to the squares of $`J_{i_1,\mathrm{},i_p}`$: $$\frac{\beta ^2}{2}\frac{^2(\beta \mathrm{\Gamma })}{\beta ^2}\|_{\beta =0}=\frac{\beta ^2}{2}\underset{1i_1<\mathrm{}<i_pN}{}Y_{i_1,\mathrm{},i_p}^2_{\beta =0}^c.$$ (8) Using again the statistical properties of the couplings and neglecting terms giving a contribution of an order smaller than $`N`$ we find that the reaction term depends on $`m_i`$ through the overlap $`q=\frac{1}{N}_im_i^2`$ only: $$\frac{\beta ^2}{2}\frac{^2(\beta \mathrm{\Gamma })}{\beta ^2}\|_{\beta =0}=\frac{\beta ^2N}{4}\left(l^pq^pp(lq^{p1}q^p)\right).$$ (9) Higher derivatives lead to terms which can be neglected because they are not of the order of N ; so collecting (3),(4) and (9) we find the TAP free energy for spherical p-spin models. Differentiating the free energy with respect to magnetisations $`m_i`$ and the spherical parameter $`l`$ one finds the TAP equations. These equations admit for certain temperatures an infinite number of solutions. This is a fundamental characteristic and difficulty of mean field spin glasses. It has been shown that the weighted sum of the local minima of the TAP free energy gives back equilibrium results found by the replica or the cavity method : $`Z=_\alpha e^{N\beta f_\alpha }`$, where $`f_\alpha `$ is the TAP free energy of a stable solution $`\{m_i^\alpha \}`$ of TAP equations. Note that states which do not have the minimum free energy can dominate the previous sum if their number is very large. Let us conclude this section with few comments on the derivation of the static TAP equations. First of all we remark that we have improperly called $`\mathrm{\Gamma }(\beta ,m_i,l)`$ a Legendre transform. Indeed the function $`\mathrm{\Gamma }(\beta ,m_i,l)`$ is the generating functional of proper vertices . This function may have many minima and is not convex in general. Finally we want to point out a striking difference, which arises in the computation of $`\mathrm{\Gamma }`$, between completely connected and finite connectivity mean field models. For the former the expansion of $`\mathrm{\Gamma }`$ in powers of $`\beta `$ stops at the second order in $`\beta `$. Whereas for the latter the expansion contains all the powers of $`\beta `$. Roughly speaking for the Sherrington-Kirkpatrick model the only non trivial term in $`\mathrm{\Gamma }`$ is the reaction term, which represents the contribution to the effective field of the ith spin due to the influence of the ith spin on the others. Whereas for its counterpart on a Bethe lattice the interaction between two neighbouring spins has to be taken into account exactly, i.e. one has to take into account not only the reaction of the neighbours of $`S_i`$ due to the presence of $`S_i`$, but also the reaction of the reaction and so on. ### B Dynamical TAP equations In the following we focus on a Langevin relaxation dynamics for mean field glassy systems. Standard field theoretical manipulations lead to the Martin-Siggia-Rose generating functional for the expectation values of $`s_i(t)`$. Within the superspace notation the dynamics and the static theory are formally very similar . As a consequence dynamical TAP equations can be derived straightforwardly generalising the method described in the previous section. We refer to for a detailed derivation. Once the dynamical TAP free energy is known, the dynamical TAP equation are obtained from the Lagrange relation for the supermagnetisation. In the following we simply quote the result : $`{\displaystyle \frac{}{t}}\left(C(t,t^{})Q(t,t^{})\right)`$ $`=`$ $`2R(t^{},t)\lambda (t)\left(C(t,t^{})Q(t,t^{})\right)+\mu {\displaystyle _0^t^{}}𝑑t^{\prime \prime }\left(C(t,t^{\prime \prime })^{p1}Q(t,t^{\prime \prime })^{p1}\right)R(t^{},t^{\prime \prime })`$ (10) $`+`$ $`\mu (p1){\displaystyle _0^t}𝑑t^{\prime \prime }\left(C(t^{\prime \prime },t^{})Q(t^{\prime \prime },t^{})\right)R(t,t^{\prime \prime })C(t,t^{\prime \prime })^{p2},`$ (11) $`{\displaystyle \frac{}{t}}R(t,t^{})`$ $`=`$ $`\lambda (t)R(t,t^{})+\delta (tt^{})+\mu (p1){\displaystyle _t^{}^t}𝑑t^{\prime \prime }R(t,t^{\prime \prime })R(t^{\prime \prime },t^{})C(t,t^{\prime \prime })^{p2},`$ (13) $`\left({\displaystyle \frac{}{t}}+\lambda (t)\right)m_i(t)`$ $`=`$ $`\beta h_i(t)+\beta {\displaystyle \underset{1i_2<\mathrm{}<i_pN}{\overset{}{}}}J_{i,i_2,\mathrm{},i_p}m_{i_2}(t)\mathrm{}m_{i_p}(t)`$ (15) $`+`$ $`\mu (p1){\displaystyle _0^t}dt^{\prime \prime }(C(t,t^{\prime \prime })^{p2}Q(t,t^{\prime \prime })^{p2}))R(t,t^{\prime \prime })m_i(t^{\prime \prime }),`$ (16) where $`C(t,t^{})=\frac{1}{N}_{i=1}^Ns_i(t)s_i(t^{})`$ is the correlation function, $`R(t,t^{})=\frac{1}{N}_{i=1}^N\frac{s_i(t)}{h_i(t^{})}`$ is the response function to magnetic fields $`h_i(t)`$ coupled to the spins $`S_i`$, $`Q(t,t^{})=\frac{1}{N}_{i=1}^Nm_i(t)m_i(t^{})`$ is the overlap function, $`m_i(t)`$ are the local magnetisation, $`\mu =p\beta ^2/2`$ and $`\lambda (t)`$ is the spherical constraint which fixes $`C(t,t)=1`$. The correlation function satisfies the boundary condition $`C(t,0)=Q(t,0)`$ and magnetisations fulfil the initial conditions $`m_i(0)=s_i^0`$ . Note that now $``$ means the average over the thermal noise. Moreover the spherical condition $`C(t,t)=1`$ fixes $`\lambda `$ as a function of time through the equation: $`\lambda (t)\left(1q(t)\right)`$ $`=`$ $`1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{dq}{dt}}+\mu {\displaystyle _0^t}𝑑t^{\prime \prime }\left(C(t,t^{\prime \prime })^{p1}Q(t,t^{\prime \prime })^{p1}\right)R(t,t^{\prime \prime })`$ (17) $`+`$ $`\mu (p1){\displaystyle _0^t}𝑑t^{\prime \prime }\left(C(t^{\prime \prime },t)Q(t^{\prime \prime },t)\right)R(t,t^{\prime \prime })C(t,t^{\prime \prime })^{p2},`$ (18) where $`q(t)=Q(t,t)`$. Three important remarks are in order on these equations. First of all if one takes for initial condition a uniform average over all possible configurations as in , then the magnetisations are equal to zero at $`t=0`$ and there is no boundary condition on the correlation function . In this case we find that the equation (15) is trivially satisfied and equations (10), (13) and (17) reduce to the ones considered in . Moreover we notice that in the zero temperature limit the equation (15) coincides with a simple gradient descent, as should be when the thermal noise is absent. Finally, it is interesting to remark that the equations on local magnetisations do not have at finite times the form of a gradient descent in the free energy landscape since the Onsager reaction term is non-Markovian. This is natural because it represents the contribution to the effective field of the $`i`$th spin due to the influence at previous times of the $`i`$th spin on the others. ### C Asymptotic analysis In the following we perform an asymptotic analysis of the equations (10), (13), (15) and (17). For the sake of simplicity we will take $`h_i(t)=0`$ in (15). Two asymptotic behaviour have been found for the p-spin spherical model depending on the choice of the initial conditions : True ergodicity breaking: the system equilibrates in a separate ergodic component. Asymptotically time homogeneity and fluctuation-dissipation theorem (FDT) hold . In this case, following , we take for the asymptotic form of the two time quantities the Ansatz: $`C(t,t^{})=C_{FDT}(tt^{})`$ $`,`$ $`R(t,t^{})=R_{FDT}(tt^{})`$ (19) $`R_{FDT}(\tau )=\theta (\tau ){\displaystyle \frac{dC_{FDT}(\tau )}{d\tau }}`$ $`,`$ $`Q(t,t^{})=q`$ (20) $`\underset{\tau \mathrm{}}{lim}C_{FDT}(\tau )=q.`$ (21) * Slow dynamics: the system does not equilibrate. Asymptotically two time sectors can be identified. In the first one (FDT regime), which corresponds to finite time differences $`\|tt^{}\|O(1)`$, ($`t>>1`$, $`t^{}>>1`$), the system has a pseudo-equilibrium dynamics since FDT and time translation invariance hold asymptotically. In the second one (ageing regime), which corresponds to “infinite” time differences $`\|tt^{}\|t^{}`$, FDT and time translation invariance do not apply and the system ages . In this case, following , we take for finite time separations the Ansatz corresponding to equilibrium dynamics, but with $`Q(t^{},t)=q^{}`$. Whereas for the ageing sector we take the AnsatzThe asymptotic equations are obtained neglecting the time derivatives. This has as a consequence that from an asymptotic solution we obtain infinitely many others by re-parameterisation . For the sake of clarity in the following we focus on the particular parameterisation shown in equations (22), (23) and (24) . : $`C(t,t^{})=qC_{ag}(\lambda )`$ $`,`$ $`tR(t,t^{})=R_{ag}(\lambda )`$ (22) $`R_{ag}(\lambda )=xq{\displaystyle \frac{dC}{d\lambda }}`$ $`,`$ $`Q(t,t^{})=q^{}Q_{ag}(\lambda )`$ (23) $`C_{ag}(1)=Q_{ag}(1)=1`$ $`,`$ $`\lambda ={\displaystyle \frac{t^{}}{t}},`$ (24) where $`x`$ parameterises the violation of FDT . The asymptotic solutions arising from the previous Ansätze can be grouped in three classes. #### 1 Equilibrium dynamics. We denote respectively by $`\lambda ^{\mathrm{}}`$ and $`m_i^{\mathrm{}}`$ the asymptotic values of the spherical multiplier and of the local magnetisations. Plugging the equilibrium dynamics Ansatz into the dynamical TAP equations we find that the equations on $`m_i^{\mathrm{}}`$ and $`\lambda ^{\mathrm{}}`$ are the corresponding static TAP equations. In the asymptotic limit the equations (10) and (13) on the correlation and the response functions reduce to: $$\left(\frac{d}{d\tau }+\lambda ^{\mathrm{}}\mu \right)C(\tau )+\mu +1\lambda ^{\mathrm{}}=\mu _0^\tau 𝑑\tau ^{}C(\tau \tau ^{})^{p1}\frac{dC(\tau ^{})}{d\tau ^{}}.$$ (25) The above equation describes the equilibrium dynamics inside the ergodic component associated to a TAP solution $`\{m_i^{\mathrm{}}\}`$. Note that this asymptotic dynamical solution is consistent with the assumption of an equilibrium dynamics only if $`\{m_i^{\mathrm{}}\}`$ is a local minimum of the free energy. Since this asymptotic solution represents the equilibration in a stable TAP state $`\{m_i^{\mathrm{}}\}`$, it is quite natural to associate to this solution an initial condition belonging to this state. This interpretation is suggested by the results of . Indeed in the low temperature dynamics has been studied starting from an initial condition belonging to the TAP states which are the equilibrium states at a temperature $`T^{}`$. In it has been shown that the system relaxes in the TAP states associated to the initial condition. It is easy to show that the equation satisfied by $`C(\tau )`$ in can be written in the form (25). Moreover it is interesting to note that the equations (15) on local magnetisations reduce in the long-time limit to a gradient descent in the free energy landscape with an extra term which vanishes at large time. #### 2 Weak ergodicity breaking. The asymptotic analysis in the time sector corresponding to finite time differences leads to the same equation (25) for the correlation and the response functions. Whereas for infinite time differences we find that the asymptotic equations admit the solution: $`q^{}=0`$, $`q`$ which verifies the equation of the overlap of the threshold states , $`x=\frac{(p2)(1q)}{q}`$ and $`C_{ag}(\lambda )`$ and $`R_{ag}(\lambda )`$, which satisfy the same equations found in . The equation (17) on the spherical multiplier reduces to: $`\lambda ^{\mathrm{}}=(1q)^1+\mu (1q^{p1})`$ and the asymptotic value of the local magnetisations $`m_i^{\mathrm{}}`$ is zero. This is exactly the same asymptotic solution found in for random initial conditions. Therefore it is natural to associate to this solution a random initial condition, which is not correlated with any particular stable TAP state. Note that the difference between $`q`$ and $`q^{}`$ clearly marks that the system does not equilibrate in a single ergodic component. #### 3 Between true and weak ergodicity breaking. In the following we consider the asymptotic solution which corresponds to slow dynamics with $`q=q^{}`$. In this case we find the same solution of section II.C.2 except that $`q^{}=q`$ and $`Q_{ag}(\lambda )=C_{ag}(\lambda )`$. As a consequence the local magnetisations do not vanish in the long-time limit. These results indicate that at very large times the system has almost thermalized within a threshold state. Anyway the slow behaviour of the overlap function $`Q(t,t^{})`$ implies that the local magnetisations evolve forever, even if more and more slowly. In other words if one waits a time $`t_w(>>1)`$ the systems seems to be equilibrated in a certain threshold states on timescales $`\mathrm{\Delta }t<<t_w`$; however on timescales of the same order of $`t_w`$ the system continue to evolve. To understand the slow evolution of $`m_i(t)`$ it is important to recall that the threshold states are characterised by a spectrum of the free energy Hessian which is a semicircle law with minimum eigenvalue equal to zero . As a consequence the free energy landscape around threshold states is characterised by almost flat directions. At large times, the equations satisfied by $`m_i(t)`$ corresponds to a gradient descent in the free energy landscape with an extra term which vanish in the long-time limit. Because of almost flat directions this vanishing term plays a fundamental role and is responsible for ageing. In fact at large times the dynamics takes place only along almost flat directions and this vanishing function of time acts as a vanishing source of drift, so the larger is the time, the weaker is the drift and the slower is the evolution: the system ages. Finally we remark that it seems natural that the initial conditions related to this asymptotic solution are the configurations typically reached in the long-time dynamics (starting from a random initial condition). In fact a way to obtain this asymptotic solution starting from a random initial condition is to introduce fields $`h_i(t)`$ which enforce the condition $`lim_t\mathrm{}1/N_{i=1}^Nm_i(t)^2=q^{}=q_{th}`$ (where $`q_{th}`$ is the overlap of threshold states ). There are many different way to fix the fields $`h_i(t)`$ to enforce this condition; however for each realization of $`h_i(t)`$ it is clear that $`lim_t\mathrm{}h_i(t)=0`$ because the equality between $`q^{}`$ and $`q_{th}`$ is automatically verified in the long-time limit. The vanishing of the local magnetisations is due to the many possible channels that the system can follow in the energy landscape. The role of magnetic fields $`h_i(t)`$ is to bring the system along one of the possible channel. ## III Free energy landscape and long-time dynamics At finite times, the dynamics cannot be represented as an evolution in the free energy landscape because the Onsager reaction term in (15) is non-Markovian. However in the long time regime a connection between the free energy landscape and the dynamical evolution can be established. For initial conditions leading to an equilibrium dynamics, i.e. the equilibration in a stable TAP state $`\{m_i^{\mathrm{}}\}`$, the equations on the local magnetisations imply that the relaxation of $`\{m_i(t)\}`$ toward $`\{m_i^{\mathrm{}}\}`$ coincides with a gradient descent in the free energy landscape with an extra term going to zero at large times. Conversely, in the most interesting and the most physical case of random initial conditions (corresponding to a quench from infinite temperature) the local magnetisations vanish at large times. Anyway a description of the asymptotic dynamics as an evolution in the free energy landscape makes sense also in this case. The local magnetisations vanish asymptotically because the dynamical probability measure at large time tends toward a static probability measure which is broken in separate ergodic components, i.e. the threshold states. One can think at the probability density in configuration space as a wave packet which breaks continuously in sub-packets. Within this picture, the dynamical evolution is characterised by two effects: the cloning of each packet in sub-packets and the slow motion of each single packet. To avoid the spreading of the dynamical measure and to capture only the slow motion, one can take for initial condition a configuration typically reached in the long-time dynamics (starting from random initial conditions). This procedure leads to the asymptotic solution analysed in section II.C.3, in which the correlation and the response functions have the same asymptotic behaviour that for a random initial condition. Moreover $`C(t,t^{})`$ and $`Q(t,t^{})`$ are equal in the ageing time regime. Thus, also the ageing dynamics obtained starting from a random initial condition can be represented in terms of the equation on $`m_i(t)`$, i.e. as a motion in the flat directions of the free energy landscape. In conclusion, through the dynamical TAP approach we have shown that the long-time dynamics can be represented as a gradient descent in the free energy landscape with an extra term going to zero at large time. This result allow one to make a straightforward connection between static free energy landscape and long-time dynamics and to give to the former a meaningful dynamical interpretation. In fact, consider all the stationary and stable dynamical probability distributions $`P_\alpha (\{s_i(t)\})`$. We have shown that the local magnetisations $`m_i^\alpha =s_i_\alpha `$ calculated with the probability law $`P_\alpha (\{s_i(t)\})`$ are the local minima of the TAP free energy. This gives to static TAP solutions a dynamical interpretation in which the properties of stationarity and stability in the free energy landscape are directly related to the properties of stationarity and stability of dynamical distributions $`P_\alpha (\{s_i(t)\})`$. Moreover the relationship, that we have elucidated, between ageing and flat directions in the free energy landscape allows one to clarify the important role played by the threshold states in the slow dynamics: they are stable states having flat directions in the free energy landscape and as a consequence they are related to ageing dynamics. What is missing to a complete dynamical interpretation of the free energy landscape is the comprehension of the role played by the free energy barriers in the activated dynamics, i.e. to go below the threshold energy starting from random initial conditions. Recent progress in this direction has been done in . ## IV Conclusions In summary we have found that for the p-spin spherical model the representation of the long-time dynamics as an evolution in the free energy landscape is correct. This evolution consists in a gradient descent in the free energy landscape with an extra term going to zero at large time. This vanishing source of drift depends on the history of the system and is crucial for slow dynamics. Our results explicitly show that the scenario for slow dynamics found at zero temperature remains valid also at finite temperature: ageing is due to the motion in the flat directions of the free energy landscape in presence of a vanishing source of drift. Finally, the relationship between long-time dynamical behaviour and local properties of the free energy landscape, which was already found in , shows up explicitly by the study of the dynamical TAP equations. This relationship is very important not only from a theoretical point of view, but also from a technical one. Indeed it allows to obtain information about the long-time dynamics by a pure static computation . For these reasons it would be very interesting to generalise the study performed in this article to finite dimensional systems. In this case the free energy landscape cannot be computed exactly and the long-time dynamics cannot be solved; however, the formal analogies (due to superspace notation) between static and dynamic theory let us hope that one can obtain results on the relationship between long-time dynamics and free energy landscape only using the symmetry properties of the asymptotic solution . ## V Acknowledgements I am deeply indebted to L. F. Cugliandolo, J. Kurchan and R. Monasson for numerous, helpful and thorough discussions on this work. I wish also to thank S. Franz and M. A. Virasoro for many interesting discussions on the asymptotic solution analysed in section 2.3.3. I am particularly grateful to R. Monasson for his constant support and for a critical reading of the manuscript. ## VI References
warning/0003/cs0003041.html	ar5iv	text	# Coherence, Belief Expansion and Bayesian Networks ## Introduction Suppose that one receives information $`\{\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}\}`$ from $`n`$ independent but less than fully reliable sources. Is it rational to believe this information? Following a tradition in epistemology that goes back to John Locke, we let belief correspond to a sufficiently high degree of confidence (Foley 1992, Hawthorne and Bovens 1999). There are three factors that determine this degree of confidence: (i) How surprising is the information? (ii) How reliable are the sources? (iii) How coherent is the information? First, suppose that the sources are halfway reliable and the information is halfway coherent. Then certainly the degree of confidence will be greater when the reported information is less rather than more surprising. Second, suppose that the information is halfway surprising and is halfway coherent. Let truth-tellers provide fully reliable information and let randomizers flip a coin for each proposition to determine whether they will affirm or deny it. Then certainly the degree of confidence will be greater when the sources are more like truth-tellers than when they are more like randomizers. Third, consider the following story: a scientist runs two independent tests to determine the locus of a genetic disease on the human genome. In the first case, the tests respectively point to two fairly narrow regions that just about overlap in a particular region. In the second case, the tests respectively point to fairly broad regions that have minimal overlap in the very same region. Suppose that the tests are halfway reliable and that this region is a somewhat surprising locus for the disease. Then certainly the degree of confidence that the locus of the disease is in this region is greater in the former case, in which the information is more coherent, than in the latter case, in which the information is less coherent. We define measures for each of these determinants of the degree of confidence in a probabilistic framework. The real challenge lies in developing a measure of coherence (cf. Lewis 1946, Bonjour 1985). This measure defines a partial ordering over information sets. Subsequently, we argue that belief expansion is a function of the reliability of the sources and the coherence of the new information with the information that we already believe. We construct an acceptance measure which determines whether newly acquired information can be added to our beliefs under alternative suppositions about the reliability of the sources. Our calculations rest on some results in the theory of Bayesian Networks. Throughout we have made some strong idealizations. We show how these idealizations can be relaxed by directly invoking Bayesian Networks. ## The Model For each proposition $`\mathrm{R}_\mathrm{i}`$ (in roman script) in the information set, let us define a propositional variable $`R_i`$ (in italic script) which can take on two values, viz. $`\mathrm{R}_\mathrm{i}`$ and $`\overline{\mathrm{R}}_\mathrm{i}`$ (i.e. not-$`\mathrm{R}_\mathrm{i}`$) for $`\mathrm{i}=1,\mathrm{},\mathrm{n}`$. Let $`REPR_i`$ be a propositional variable which can take on two values, viz. $`\mathrm{REPR}_\mathrm{i}`$, i.e. there is a report from the proper source to the effect that $`\mathrm{R}_\mathrm{i}`$ is true, and $`\overline{\mathrm{REPR}}_\mathrm{i}`$, i.e. there is a report to the effect that $`\mathrm{R}_\mathrm{i}`$ is false. We construct a joint probability distribution $`P`$ over $`R_1,\mathrm{},R_n,REPR_1,\mathrm{},REPR_n`$ satisfying the constraint that the sources are independent and less than fully reliable. We model independence by stipulating that $`P`$ respects the following conditional independences: $$REPR_iR_j,REPR_j\|R_i\mathrm{for}ij;i,j=1,2,\mathrm{},n$$ (1) or, in words, $`REPR_i`$ is probabilistically independent of $`R_j,REPR_j`$, given $`R_i`$, for $`ij`$ and $`i,j=1,2,\mathrm{},n`$. What this means is that the probability that I will receive a report that $`\mathrm{R}_\mathrm{i}`$ given that $`\mathrm{R}_\mathrm{i}`$ is the case or given that $`\mathrm{R}_\mathrm{i}`$ is not the case, is not affected by any additional information about whether $`\mathrm{R}_\mathrm{j}`$ is the case or whether there is a report to the effect that $`\mathrm{R}_\mathrm{j}`$ is the case. Each source tunes in on the item of information that it is meant to report on: it may not always provide an accurate report, but its report is not affected by what other sources have to report or by other items of information than the one it reports on (Lewis 1946, Bovens and Olsson 1999). We define a less-than-fully-reliable source as a source that is better than a randomizer, but short of being a truth-teller and make the simplifying idealization that the information sources are equally reliable. We specify the following two parameters: $`P(\mathrm{REPR}_\mathrm{i}\|\mathrm{R}_\mathrm{i})=p`$ and $`P(\mathrm{REPR}_\mathrm{i}\|\overline{\mathrm{R}}_\mathrm{i})=q`$ for $`\mathrm{i}=1,\mathrm{},\mathrm{n}`$. If the information sources are truth-tellers, then $`p=1`$ and $`q=0`$, while if they are randomizers, then $`p=q>0`$. We model less-than-full-reliability by imposing the following constraint on $`P`$: $$P(\mathrm{REPR}_\mathrm{i}\|\mathrm{R}_\mathrm{i})=p>q=P(\mathrm{REPR}_\mathrm{i}\|\overline{\mathrm{R}}_\mathrm{i})>0$$ (2) The degree of confidence in the content of the information set is the posterior joint probability after all the reports have come in: $$P^{}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})=P(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}\|\mathrm{REPR}_1,\mathrm{},\mathrm{REPR}_\mathrm{n})$$ (3) The motivation for the definition of less-than-full reliability is that we are interested in cases in which incoming information raises our confidence in the content of the information set to different levels. When the sources are randomizers, our confidence will be unaffected (Huemer 1997, Bovens and Olsson 1999), i.e. $`P^{}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})=\mathrm{P}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$; when they are truth-tellers, our confidence will be raised to certainty, i.e. $`P^{}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})=1`$; and when they are worse than randomizers, our confidence will drop, i.e. $`P^{}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})<\mathrm{P}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$. ## Expectation, Reliability and Coherence It can be shown by the probability calculus, that, given the constraints on $`P`$ in (1) and (2), $$P^{}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})=\frac{a_0}{_{i=0}^na_ix^i},$$ (4) in which the likelihood ratio $`x=q/p`$ (note that $`0<x<1`$ for $`p>q>0`$) and $`a_i`$ is the sum of the joint probabilities of all combinations of values of the variables $`R_1,\mathrm{},R_n`$ that have $`i`$ negative values and $`ni`$ positive values: e. g. for $`n=3`$, $`a_2=P(\mathrm{R}_1,\overline{\mathrm{R}}_2,\overline{\mathrm{R}}_3)+P(\overline{\mathrm{R}}_1,\mathrm{R}_2,\overline{\mathrm{R}}_3)+P(\overline{\mathrm{R}}_1,\overline{\mathrm{R}}_2,\mathrm{R}_3)`$. Note that $`_{i=0}^na_i=1`$. We can directly identify the first determinant of the degree of confidence in the information set. Note that $`a_0=P(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$ is the prior joint probability of the propositions in the information set, i.e. the probability before any information was received. This prior probability is lower for more surprising information and higher for less surprising information. Since more surprising information is tantamount to less expected information, let us call this prior probability the expectation measure. It is easy to see that $`P^{}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$ is a monotonically increasing function of $`a_0`$. We can also directly identify the second determinant, i.e. the reliability of the sources. Note that $`P^{}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$ is a monotonically decreasing function of $`x=q/p`$. Hence, let us call $`r:=1x`$ the reliability measure, since $`P^{}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$ is a monotonically increasing function of $`r`$ and this measure ranges from $`0`$ for sources that are randomizers to $`1`$ for sources that are truth-tellers. It is more difficult to construct a coherence measure. Consider the following analogy: to assess the impact of a training program, we consider the rate of the student’s actual performance level over the performance level that he would have reached in an ideal training program, all other things equal. Similarly, to assess the impact of coherence, we consider the rate of the present degree of confidence over the degree of confidence that would have been obtained had the information set been maximally coherent, all other things equal. The information set would have been maximally coherent if and only if $`\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}`$ had all been coextensive. Let $`P`$ be the actual joint probability distribution. Construct a joint probability distribution $`P^{max}`$ with the same expectation measure and the same reliability measure as $`P`$, but $`\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}`$ are all coextensive, i.e., on $`P^{max}`$, $`a_0`$ is the same as on $`P`$, but $`a_n=1a_0=:\overline{a}_0`$, so that $`a_i=0`$, for all $`i0,n`$. It follows from (4) that, $$P^{max}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})=\frac{a_0}{a_0+\overline{a}_0x^n}.$$ (5) Hence, for $`a_00`$, the ratio $`c_x(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$ $`=`$ $`{\displaystyle \frac{P^{}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})}{P^{max}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})}}`$ (6) $`=`$ $`{\displaystyle \frac{a_0+\overline{a}_0x^n}{_{i=0}^na_ix^i}}`$ is a measure of the impact of the coherence of the information set on the degree of confidence in the content of the information set. But note that this measure is contingent on the value of the reliability measure: (6) only provides us with a reliability-relative coherence measure. This is unwelcome: there is a pretheoretical notion of the coherence of an information set which has nothing to do with the reliability of the sources that provides us with their content. On the other hand, this pretheoretical notion seems to be an ordinal rather than a cardinal notion. And furthermore, it seems to require a partial rather than a complete ordering over information sets: for certain, though not for all pairs of information sets, we are prepared to pass a judgment that one set in the pair is more or less coherent than the other. It turns out that the reliability-relative coherence measure indeed induces a partial ordering over informations sets which is not contingent on the reliability of the sources. Consider two information sets of size $`n`$. These sets can be represented by the marginal probability distributions $`P`$ and $`P^{}`$ over $`R_1,\mathrm{},R_n`$. It can be shown that for some $`P`$ with $`a_0,\mathrm{},a_n`$ and $`P^{}`$ with $`a_0^{},\mathrm{},a_n^{}`$, the difference $`c_x(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})c_x^{}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$ has the same sign for any value of $`x`$ ranging from $`0`$ to $`1`$. Hence, the reliability-relative coherence measure $`c_x(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$ induces a partial coherence ordering over information sets that is not contingent on the reliability of the sources. For information pairs, i.e. for information sets containing exactly two propositions, it can be shown that the following is a necessary and sufficient condition for inclusion in the partial coherence ordering: $`P`$ and $`P^{}`$ are such that (i) $`a_0/a_0^{}a_1/a_1^{}`$ and $`a_1a_1^{}`$, or, (ii) $`a_0/a_0^{}a_1/a_1^{}`$ and $`a_1a_1^{}`$. For information sets in general, it can be shown that the following is a sufficient condition for inclusion in the partial coherence ordering: $`P`$ and $`P^{}`$ are such that (i) $`a_i/a_i^{}<a_0/a_0^{}<1`$, or, (ii) $`a_i/a_i^{}>a_0/a_0^{}>1`$, for $`i=1,\mathrm{},n1`$. We provide an example of this condition for information pairs. Suppose that we are trying to locate a corpse of a murder somewhere in Tokyo. We draw a grid of $`100`$ squares over the map of the city so that it is equally probable that the murder occurred in each grid. We interview two independent less-than-fully-reliable sources. Source $`1`$ reports that the corpse is somewhere in squares $`41`$ to $`60`$ and source $`2`$ reports that the corpse is somewhere in squares $`51`$ to $`70`$. In this case, $`a_0=.10`$ and $`a_1=.20`$. This our base case. Now consider alternate case $`A`$ in which source $`1`$ reports squares $`50`$ to $`60`$ and source $`2`$ reports squares $`51`$ to $`61`$. In this case, $`a_0^{}=.10`$ and $`a_1^{}=.02`$. The information set in alternate case $`A`$ is clearly more coherent than in the base case. Notice that the condition for a partial ordering is indeed satisfied. But now consider alternate case $`B`$: source $`1`$ reports squares $`26`$ to $`60`$ and source $`2`$ reports squares $`41`$ to $`75`$. In this case $`a_0^{\prime \prime }=.20`$ and $`a_1^{\prime \prime }=.30`$. Is the information set in alternate case $`B`$ more coherent than in the base case? The proportion of the reported squares that overlap in each report is greater in the alternate case, which suggests that there is more coherence. But on the other hand, the price of getting more proportional overlap is that the overlapping area is less precise and that both sources make a much broader sweep over the map, suggesting less coherence. Indeed, in this case, we cannot pass judgment whether the information set in alternate case $`B`$ is more coherent than in the base case. Notice that the condition for a partial ordering is indeed not justified. ## Belief Expansion Suppose that we acquire various items of background information from various sources and that our degree of confidence in the content of the information set is sufficiently high to believe the information. Now a new item of information is being presented. Are we justified to add this new item of information to what we already believe? The answer to this question has something to do (i) with the reliability of the information source as well as (ii) with the plausibility of the new information, given what we already believe, or in other words, with how well the new information coheres with the background information. The more reliable the source is, the less plausible the new information needs to be, given what we already believe, to be justified to add the new information. The more plausible the new information is, given what we already believe, the less reliable the source needs to be, to be justified to add the new information. The challenge is: can a precise account of this relationship be provided? Our approach is markedly different from AGM belief revision. In the AGM approach, the question is not whether to accept new information or not, but rather, once we have made the decision to accept the new information, how we should revise our beliefs in the face of inconsistency (Makinson 1997, Olsson 1997). Our approach shares a common motivation with the program of non-prioritized belief revision. According to Hansson (1997), we may not be willing to accept the new information because “it may be less reliable (…) than conflicting old information.” Makinson (1997) writes that “we may not want to give top priority to new information (…) we may wish to weigh it against old material, and if it is really just too far-fetched or incredible, we may not wish to accept it.” However, whereas the program of non-prioritized belief revision operates within a logicist framework, we construct a probabilistic model. The cost of this approach is that it is informationally more demanding. The benefit is that it is empirically more adequate, because it is sensitive to degrees of reliability and coherence and to their interplay in belief acceptance. In non-prioritized belief revision, the reliability of the sources does not enter into the model itself and the lack of coherence of an information set is understood in terms of logical inconsistency, which is only a limiting case in our model. To introduce the approach, we address the question of belief expansion. We believe that our model also carries a promise to handle belief revision in general, but this project is beyond the scope of this paper. We need to make some simplifying assumptions about the origin of the background information and the new information: (a) the propositions in the background information are provided by independent sources, which are (b) less than fully reliable, (c) equally reliable as the new source, and (d) independent of the new source. Our background information is contained in $`\{\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}\}`$. Now suppose that we have a certain threshold level for belief and that the degree of confidence for the background information after having received a report to this effect from independent less than fully reliable sources is right at this level. (This stipulation is not required if we model actual cases by means of Bayesian Networks.) Now we are handed a new item of information $`\mathrm{R}_{\mathrm{n}+1}`$ by an independent less than fully reliable source. Then we will expand our belief set from $`\{\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}\}`$ to $`\{\mathrm{R}_1,\mathrm{},\mathrm{R}_{\mathrm{n}+1}\}`$ if and only if $`P(\mathrm{R}_1,\mathrm{},\mathrm{R}_{\mathrm{n}+1}\|\mathrm{REPR}_1,\mathrm{},\mathrm{REPR}_{\mathrm{n}+1})`$ $`P(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}\|\mathrm{REPR}_1,\mathrm{},\mathrm{REPR}_\mathrm{n}).`$ (7) Our sources are independent: $$REPR_iR_j,REPR_j\|R_i\mathrm{for}ij;i,j=1,\mathrm{},n+1$$ (8) (4) defines an acceptance measure for an information set: $$e_x(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{m})=P^{}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{m})=\frac{a_0}{_{i=0}^ma_ix^i}$$ (9) Considering (6) and (9), we can define this acceptance measure in terms of the reliability-relative coherence measure $`c_x`$, provided that $`a_00`$: $$e_x(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{m})=\frac{a_0}{a_0+\overline{a}_0x^m}c_x(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{m})$$ (10) ¿From (Belief Expansion) and (9), it follows that we can expand our belief set with a new item of information if and only if $$e_x(\mathrm{R}_1,\mathrm{},\mathrm{R}_{\mathrm{n}+1})e_x(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}).$$ (11) We can make the following two observations: 1. From (9) and (11), it is clear that whether we can expand our beliefs or not, is a complex function of the reliability of the sources and the dependence of new on earlier information as expressed in the probability distribution over the variables $`R_1,\mathrm{},R_{n+1}`$. The reliability of the sources is reflected in the likelihood ratio $`x`$ and the dependence of new on earlier information is reflected in the series $`a_0,\mathrm{},a_n`$ for $`e_x(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$ and in the series $`a_0^{},\mathrm{},a_{n+1}^{}`$ for $`e_x(\mathrm{R}_1,\mathrm{},\mathrm{R}_{\mathrm{n}+1})`$. 2. From (10), it is clear that the acceptance measure is a weighted reliability-relative coherence measure. The weight tends to 1 for smaller values of $`x`$, i.e. for more reliable sources, and for greater values of $`n`$, i.e. for larger information sets, so that the acceptance measure will coincide with $`c_x`$. We have shown that this measure lets us construct a coherence ordering over a pair of information $`n`$-tuples, if certain conditions are met. We conjecture that such an ordering can also be constructed over pairs containing an information $`n`$-tuple and an expansion of this $`n`$-tuple, i.e. over pairs of the form $`\{\{\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}\},\{\mathrm{R}_1,\mathrm{},\mathrm{R}_{\mathrm{n}+1}\}\}`$, if certain conditions are met. Contingent on this conjecture, we can make a substantial point: if there exists a determinate answer to the relative coherence of the old and the new information sets, then the more reliable the sources are and the larger the information set is, the more the question of belief expansion is determined by whether the new information set is or is not more coherent than the old information set, and not by the reliability of the sources. The acceptance measure depends, at least to some extent, on the value of the likelihood ratio $`x`$. But what, one might ask, should we do when we have no clue whatsoever about the reliability of the sources, except that they are better than mere randomizers and yet less than fully reliable? Let us model our limited knowledge as a uniform distribution over the values $`p`$ and $`q`$ under the constraint that $`p>q`$. Then we can construct the following averaged acceptance measure: $`E(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{m})`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle _0^p}e_{q/p}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{m})𝑑q𝑑p`$ (12) $`=`$ $`{\displaystyle _0^1}e_x(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{m})𝑑x`$ We can formulate a general criterion for belief acceptance: when we have limited knowledge about the reliability of our information sources, we can expand our belief set from $`\{\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}\}`$ to $`\{\mathrm{R}_1,\mathrm{},\mathrm{R}_{\mathrm{n}+1}\}`$ if and only if $$E(\mathrm{R}_1,\mathrm{},\mathrm{R}_{\mathrm{n}+1})E(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}).$$ (13) ## Bayesian Networks Bayesian Networks represent (conditional) independences between variables and when implemented on a computer they perform complex probabilistic calculations at the touch of a keystroke. We are assuming here that the reader has some familiarity with Bayesian Networks (Cowell et. al. 1999, Jensen 1996, Neapolitan 1990, Pearl 1988). We construct a Bayesian Network that permits us to read off the reliability-relative coherence measure of an information set $`\{\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}\}`$ in Figure $`1`$. First, we construct a Bayesian Network with nodes for the variables $`R_1,\mathrm{},R_n`$ which represents the marginal probability distribution over these variables. Then we add nodes for the variables $`REPR_1,\mathrm{},REPR_n`$ and draw in an arrow from each node for the variable $`R_i`$ to the node for the variable $`REPR_i`$ and specify the conditional probabilities in (2) for each arrow. By the standard criterion of $`d`$-separation, we can now read off the conditional independences in (1) from the network. Subsequently, we construct a node for the variable $`R_1\&\mathrm{}\&R_n`$: we draw in the arrows and specify conditional probabilities such that $`\mathrm{R}_1\&\mathrm{}\&\mathrm{R}_\mathrm{n}`$ holds if and only if $`\mathrm{R}_1,\mathrm{},`$ and $`\mathrm{R}_\mathrm{n}`$ hold. We can now read off $`P^{}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$: it is the probability of $`\mathrm{R}_1\&\mathrm{}\&\mathrm{R}_\mathrm{n}`$ after instantiating $`\mathrm{REPR}_1,\mathrm{},\mathrm{REPR}_\mathrm{n}`$. To read off $`P^{max}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$, more construction is needed. Notice that $`P^{max}(\mathrm{R}_\mathrm{i})=P^{max}(\mathrm{R}_1\&\mathrm{}\&\mathrm{R}_\mathrm{n})`$ for $`\mathrm{i}=1,\mathrm{},\mathrm{n}`$ in the counterfactual case of maximal coherence, is equal to $`P(\mathrm{R}_1\&\mathrm{}\&\mathrm{R}_\mathrm{n})`$ in the actual case where the information set may not be maximally coherent. Hence $`P^{max}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$ is the posterior joint probability of $`\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}`$, had we been informed in the actual case by $`n`$ less than fully reliable independent sources that $`\mathrm{R}_1\&\mathrm{}\&\mathrm{R}_\mathrm{n}`$. So we add nodes for the variables $`REP_i\&R`$ (whose positive values states that the $`i`$-th source informs us that $`\mathrm{R}_1\&\mathrm{}\&\mathrm{R}_\mathrm{n}`$), draw in the proper arrows and specify the proper conditional probabilities. We can now read off $`P^{max}(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$: it is the probability of $`\mathrm{R}_1\&\mathrm{}\&\mathrm{R}_\mathrm{n}`$ after instantiating $`\mathrm{REP}_1\&\mathrm{R},\mathrm{},\mathrm{REP}_\mathrm{n}\&\mathrm{R}`$. The measure $`c_x(\mathrm{R}_1\&\mathrm{}\&\mathrm{R}_\mathrm{n})`$ follows by (6). We construct a Bayesian Network in Figure $`2`$ to determine whether belief expansion is warranted or not. The construction of the nodes for the variables $`R_1,\mathrm{},R_{n+1}`$ and $`REPR_1,\mathrm{},REPR_{n+1}`$ should be clear from our construction of the Bayesian Network in Figure $`1`$. This part of the Bayesian Network respects the conditional independences in (8). Now we add a node for the variable $`R_1\&\mathrm{}\&R_n`$ and a node for the variable $`R_1\&\mathrm{}\&R_{n+1}`$ and specify the conditional probabilities so that $`\mathrm{R}_1\&\mathrm{}\&\mathrm{R}_\mathrm{n}`$ holds if and only if $`\mathrm{R}_1,\mathrm{},`$ and $`\mathrm{R}_\mathrm{n}`$ hold and $`\mathrm{R}_1\&\mathrm{}\&\mathrm{R}_{\mathrm{n}+1}`$ holds if and only if $`\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n}`$ and $`\mathrm{R}_{\mathrm{n}+1}`$ hold. We instantiate $`\mathrm{REPR}_1,\mathrm{},\mathrm{REPR}_\mathrm{n}`$ and propagate the evidence throughout the network. We can now read off the acceptance measure $`e_x(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$ which is the posterior probability of $`\mathrm{R}_1\&\mathrm{}\&\mathrm{R}_\mathrm{n}`$. To raise the question of belief expansion, this value should be greater than or equal to our threshold value for belief. Subsequently, we instantiate $`\mathrm{REPR}_{\mathrm{n}+1}`$ and propagate the evidence throughout the network. We can now read off the acceptance measure $`e_x(\mathrm{R}_1,\mathrm{},\mathrm{R}_{\mathrm{n}+1})`$ which is the posterior probability of $`\mathrm{R}_1\&\mathrm{}\&\mathrm{R}_{\mathrm{n}+1}`$. Depending on our treshold value for belief, we can determine whether we are justified to expand our beliefs with the proposition $`\mathrm{R}_{\mathrm{n}+1}`$. It is easy to see how the idealizations can be relaxed in the networks. We can stipulate alternative reliability parameters for the sources. We can let one source report on two items of information. We can add arrows between the $`REPR_i`$ variables or between some $`REPR_i`$ and $`R_j`$ variables (for $`ij`$) to model certain types of dependence between the sources. It suffices that $`P(\mathrm{R}_1,\mathrm{},\mathrm{R}_\mathrm{n})`$ is equal to or greater than the threshold value for belief. Furthermore, even if $`P(\mathrm{R}_1,\mathrm{},\mathrm{R}_{\mathrm{n}+1})`$ is below the threshold value for belief, the model yields a marginal probability distribution over $`R_1,\mathrm{},R_n`$. Hence, the general question of belief revision becomes a question of defining a function which maps joint probability distributions over a set of propositional variables into sets of propositions that are values of a subset of these variables and that can reasonably be believed. Defining such a function is beyond the scope of this paper. ## Conclusion (i) We have designed a procedure to determine a partial coherence ordering over a set of information sets of size $`n`$. If one information set is more coherent than another on this ordering, then our degree of confidence in the content of the former set will be greater than in the content of the latter set, after having been informed by independent and less than fully reliable sources, ceteris paribus. (ii) We have designed a probabilistic criterion for (non-prioritized) belief expansion, which determines whether it is rational to believe new information, considering how reliable the sources are and how well the new information coheres with the old information. (iii) If either the sources are sufficiently reliable or the information set sufficiently large, then the question of belief expansion is largely determined by whether the expanded information set is more coherent than the original information set (provided that there exists an ordering of this pair of information sets), and only marginally by the reliability of the sources. (iv) We have shown how a coherence ordering over information sets can be constructed by means of Bayesian Networks and how belief expansion can be modeled by means of Bayesian Networks in an empirically adequate manner. Acknowledgments. Thanks to Erik Olsson and the NMR referees for discussion, comments or suggestions. The research was supported by the Alexander-von-Humboldt Foundation and the German-American Council Foundation. References. Bonjour, L. 1985. The Structure of Empirical Knowledge. Cambridge, Mass.: Harvard University Press. Bovens, L., and Olsson, E. 1999. Coherentism, Reliability and Bayesian Networks. Technical Report, Logik in der Philosophie - 36, Department of Philosophy, University of Konstanz. Hawthorne, J., and Bovens, L. 1999. The Preface, the Lottery and the Logic of Belief. Mind 108: 241-264. Huemer, M. 1997. Probability and Coherence Justification. Southern Journal of Philosophy 35: 463-472. Cowell, R. G., Dawid, A. P., Lauritzen, S. L., and Spiegelhafter, D. J. 1999. Probabilistic Networks and Expert Systems. New York: Springer. Foley, R. 1992. The Epistemology of Belief and the Epistemology of Degrees of Belief. American Philosophical Quarterly 29: 111-121. Jensen, F. V. 1996. An Introduction to Bayesian Networks. Berlin: Springer. Lewis, C. I. 1946. An Analysis of Knowledge and Valuation. LaSalle, Ill.: Open Court. Makinson, D. 1997. Screened Revision. Theoria 63: 14-23. Neapolitan, R. E. 1990. Probabilistic Reasoning in Expert Systems. New York: Wiley. Olsson, E. 1997. A Coherence Interpretation of Semi-Revision. Theoria 63:105-133. Pearl, J. 1988. Probabilistic Reasoning in Intelligent Systems. San Mateo, Calif.: Morgan Kaufmann.
warning/0003/quant-ph0003028.html	ar5iv	text	# Quantum theory of self-action of ultrashort light pulses in an inertial nonlinear medium ## I INTRODUCTION During the past years the study of formation of nonclassical light pulses in nonlinear media has been the focus of a considerable attention. The present article is devoted to the development of the consecutive theory of formation of nonclassical short light pulses in nonlinear media with inertial Kerr nonlinearity. It is well known, that in nonlinear media in the presence of the self-action effect the squeezing of quantum fluctuations of one quadrature component of a field with conservation of the photon statistics take place . At present there are two basic directions of research in the quantum theory of self-action of ultrashort light pulses (USPs). In the first approach the calculations of the nonclassical light formation at the self-action of light pulses assume that the nonlinear response of nonlinearity is instantaneous and that the relative fluctuations are small. This latter assumption is valid for the intensive USPs frequently used in experiments. However, a finite relaxation time of the nonlinearity has a principled importance as the relaxation time determines the region of the spectrum of the quantum fluctuation below the standard noise level. For the first time, in was noted that for the correct quantum solution of self-action, it is necessary to take into account the presence of quantum noise. The presence of quantum noise was anticipated in as thermal addition to the interaction Hamiltonian. This addition was necessary in order to satisfy the commutation relation for time-dependent Bose-operators. If the interaction Hamiltonian has the normally ordered form then it is not necessary to have deal with thermal fluctuations. This approach allows us to develop an algebra of time-dependent Bose-operators and to investigate the spectrum of quantum fluctuations of the quadrature components. The results of the quantum theory of self-action for USPs in the medium with the relaxation Kerr nonlinearity based on the normally ordered interaction Hamiltonian and the developed algebra of time-dependent Bose-operators are presented below. ## II THE QUANTUM EQUATION OF SELF-ACTION OF USPs For the monochromatic radiation the quantum equation of self-action can be found, for example, in . Up to now, the proper quantum equation of the self-action for USPs with the account of relaxation behavior of the nonlinearity is absent in the literature. It is necessary to mention that the deduction of the quantum equation is based on the interaction Hamiltonian. However, in this case we obtain the time-evolution equation for the Bose-operators. For monochromatic radiation the conversion of time-evolution equation into space-evolution one is done using the $`tz/u`$ replacement, where $`z`$ is time variable and $`u`$ is the speed of pulse in nonlinear media. If we deal with the propagation of pulse in nonlinear media then the ”impulse operator” of a pulse field should be used . We begin with the analyse of the self-action of UPSs in non-inertial nonlinear media. ### A THE QUANTUM EQUATION OF SELF-ACTION IN THE NON-INERTIAL NONLINEAR MEDIA In nonlinear media with non-inertial behavior the self-action process is described using the impulse operator (quantity of movement) $`\widehat{G}_{int}(z)`$ $$\widehat{G}_{int}(z)=\beta \mathrm{}_{\mathrm{}}^{\mathrm{}}\widehat{𝐍}[\widehat{n}^2(t,z)]𝑑t,$$ (1) where $`\widehat{𝐍}`$ is the operator of normal ordering, factor $`\beta `$ is defined by the cubic nonlinearity of the medium . In consequence, in the Heisenberg representation the quantum space-evolution equation for the annihilation photons Bose-operator in a given cross section $`z`$ ($`\widehat{A}(t,z)`$) has the form $$i\mathrm{}\frac{\widehat{A}(t,z)}{z}=[\widehat{A}(t,z),\widehat{G}_{int}(z)].$$ (2) In agreement with (1) the quantum equation of self-action for a light pulse follows from (2) $$\frac{\widehat{A}(t,z)}{z}i\beta \widehat{A}^+(t,z)\widehat{A}^2(t,z)=0.$$ (3) Eq.(3) is written in the moving coordinate system: $`z=z^{}`$ and $`t=t^{}z/u`$. It is important to note that in comparison with the so-called nonlinear Heisenberg equation, used in the quantum theory of optical solitons, in (3) the dissipation of light pulse in the nonlinearity is not taken into account. This approach corresponds to the first approximation of dissipation theory. In fact, the traditional way to introduce the quantum equation of self-action is based on the interaction Hamiltonian. In this case one gets the time-evolution equation. The transition to the space-evolution, as already mentioned, is realized using replacement $`tz/u`$. This approach is enough reasonable in case the radiation is monochromatic. If we deal with the nonlinear propagation of a pulse, we use the impulse operator of a pulse field (1) which is connected with the evolution of field in space. Eq.(3) has the solution $$\widehat{A}(t,z)=e^{i\gamma \widehat{n}_0(t)}\widehat{A}_0(t),$$ (4) where $`\gamma =\beta z`$ and, as usual, $`\widehat{A}_0(t)`$ is the value of the operator at input of nonlinear media ($`\widehat{A}_0(t)=\widehat{A}(t,z=0)`$), $`\widehat{n}_0(t)=\widehat{A}_0^+(t)\widehat{A}_0(t)`$ is the photon number “density” operator. For its hermitian conjugated operator we find $$\widehat{A}^+(t,z)=\widehat{A}_0^+(t)e^{i\gamma \widehat{n}_0(t)}.$$ (5) In agreement with (4) and (5) the operator $`\widehat{n}(t,z)=\widehat{A}^+(t,z)\widehat{A}(t,z)`$ does not change itself in nonlinear medium: $$\widehat{n}(t,z)=\widehat{n}(t,z=0)=\widehat{n}_0(t),$$ (6) where $`z=0`$ corresponds to the input of the nonlinear medium. In fact, (6) means that the photon statistic in media remains unchanged. The commutation relation at the input ($`z=0`$) of nonlinearity $`[\widehat{A}_0(t_1),\widehat{A}_0^+(t_2)]=\delta (t_1t_2)`$, should be satisfied for any coordinate $`z`$ in nonlinear media: $$[\widehat{A}(t_1,z),\widehat{A}^+(t_2,z)]=\delta (t_1t_2).$$ (7) The solutions (4) and (5) do not permit to verify the commutation relation (7). Besides, the analyse of the statistical characteristics of the pulse is accompanied by the necessity of the reduction to the normally ordered form of the expression $`e^{i\gamma \widehat{n}_0(t)}`$. In this case, the solutions (4) and (5) are accompanied by the singularity of the function $`\delta (t)`$ at $`t=0`$. The specified circumstances represent the main deficiency of the quantum theory of self-action of USPs in non-inertial nonlinear media. ### B THE QUANTUM EQUATION OF SELF-ACTION IN THE INERTIAL NONLINEAR MEDIA In the classical theory, the self-action process in inertial nonlinear media is described by the equation (in the first approximation of the dispersion theory) $`{\displaystyle \frac{B(t,z)}{z}}+{\displaystyle \frac{1}{u}}{\displaystyle \frac{B(t,z)}{t}}i{\displaystyle \frac{k_0}{n_0}}\mathrm{\Delta }n(\|`$ $`B`$ $`(t,z)\|^2)`$ (9) $`\times B(t,z)=0,`$ where: $`B(t,z)`$\- the complex amplitude of a pulse, $`z`$\- the distance in the nonlinear media, $`u`$\- the group velocity, $`\mathrm{\Delta }n(\|B(t,z)\|^2)=\mathrm{\Delta }n((t,z)`$\- the nonlinear addition to the coefficient of refraction. We consider that the last term of (9) is caused by the high-frequency Kerr effect, and its evolution follows from the equation $$\tau _r\frac{\mathrm{\Delta }n(t,z)}{t}+\mathrm{\Delta }n(t,z)=\frac{1}{2}n_2\|B(t,z)\|^2.$$ (10) Here $`\tau _r`$ represents the relaxation time of the nonlinearity and $`n_2`$\- the nonlinear factor. We mention that in general the behaviour of the nonlinear addition differs from the one characterized by (10). However, if the carrying frequency of a pulse is far enough from the resonance and the pulse duration $`\tau _p`$ is greater that the relaxation time $`\tau _r`$, then (10) is correct . The solution of (10) looks like $`\mathrm{\Delta }n(t,z)=(n_2/2){\displaystyle _{\mathrm{}}^t}H(tt_1)\|B(t_1,z)\|^2𝑑t_1.`$ (11) The function of nonlinear response $`H(t)`$ is entered in such a way that in the limit $`\tau _p\tau _r`$ the nonlinear addition becomes $`\mathrm{\Delta }n(t,z)=(n_2/2)\|B(t,z)\|^2`$. In the moving system of coordinates $`(t^{}=tz/u,z^{}=z)`$, taking into account (11), eq. (9) takes the form $$\frac{B(t,z)}{z}=i\gamma ^{}_0^{\mathrm{}}H(t_1)\|B(tt_1,z)\|^2𝑑t_1B(t,z),$$ (12) where $`\gamma ^{}=k_0n_2/2n_0`$ and the quotation-marks ( ) in new system of coordinates further will be lowered for simplicity. The transition to the quantum equation usually is carried out in the spectral representation. However, in the considered case it is more natural to use time representation. We make in (12) replacement of the complex amplitudes with the operators, entered in the previous sections, $$B(t,z)iC\widehat{A}(t,z),B^{}(t,z)iC\widehat{A}^+(t,z).$$ (13) The right part of the equation we have gotten in this way, will be written below in normally ordered form. In order to take into account the presence of the vacuum fluctuations, existing up to the moment of arrival of the pulse, we will replace in (12) the bottom limit of integration by $`\mathrm{}`$. As a result we get ($`C=(\mathrm{}\omega _0/2V)^{1/2}`$, $`\beta =\gamma ^{}C^2/2`$) $`{\displaystyle \frac{\widehat{A}(t,z)}{z}}=i\beta {\displaystyle _{\mathrm{}}^{\mathrm{}}}H(\|t_1\|)\widehat{A}^+(tt_1,z)`$ $`\widehat{A}(tt_1,z)`$ (15) $`\times \widehat{A}(t,z)dt_1.`$ Eq.(15) represents the correct quantum equation and can be obtained from space evolution equation for the operator $`\widehat{A}(t,z)`$ in interaction representation (2). Taking into consideration the inertial behaviour of the nonlinearity, the impulse operator of a pulse should be introduced as: $`\widehat{G}_{int}(z)=\mathrm{}\beta {\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t`$ $`{\displaystyle _{\mathrm{}}^t}H(tt_1)`$ (17) $`\times \widehat{𝐍}\left[\widehat{n}(t,z)\widehat{n}(t_1,z)\right]dt_1,`$ where $`H(t)`$ is the nonlinear response of the medium (see (11)) ($`H(t)0`$ at $`t0`$ and $`H(t)=0`$ at $`t<0`$). We note that the integral expression in (17) at the moment of time $`t`$ depends only on the previous ones. Therefore, in this case the causality principle relative to measured physical value is not broken. An important condition that must be satisfied by the impulse operator is $$[\widehat{G}_{int}(z),\widehat{n}(t,z)]=[\widehat{G}_{int}(z),\widehat{n}_0(t)]=0,$$ (18) which means that the photon number operator remains unchanged in medium (photon number is a constant of motion ). The annihilation Bose-operator which agrees with the (17) verifies the equation of self-action $$\frac{\widehat{A}(t,z)}{z}i\beta q[\widehat{n}_0(t)]\widehat{A}(t,z)=0,$$ (19) where $$q[\widehat{n}_0(t)]=_0^{\mathrm{}}H(t_1)\left[\widehat{n}_0(tt_1)+\widehat{n}_0(t+t_1)\right]𝑑t_1.$$ (20) If $`H(t)=\delta (t)`$ then (19) is converted into (3). Eq.(3) can be obtained in the same form if we consider that the response of the nonlinearity has relaxation behaviour $`H(t)`$ in accordance with (11) and the relaxation time $`\tau _r`$ is significantly less than the duration of a pulse $`\tau _p`$. At the same time, as will be shown further, in this limited case also, the account of finite relaxation time plays the main role in formation of the nonclassical light. It is necessary to note, that by replacement in (19) the operators on complex values we obtain the equation, which not completely coincides with the classical equation of self-action in presence of the non-stationary nonlinear response. There is no second term, included in (20) as $`\widehat{n}_0(t+t_1)`$. The presence of this term in quantum theory, which at the first sight is in contradiction with the causality principle, is connected with the quantum description that even in absence of a pulse the vacuum fluctuations always are present. A response function was introduced already in the model for the Kerr effect by Blow et al. . Although the need for an attendant noise source was anticipated by these authors, they did not indicate where it should be inserted. In the quantum noise as thermal fluctuations was additive inserted in the interaction Hamiltonian but this procedure did not allow to develop the consistent quantum theory of self-action of USPs. According to (17), the operator $`\widehat{n}(t,z)`$ remains unchanged in nonlinear media (see (6)). Solving the spatial evolution equation (19), the annihilation (creation) photons Bose-operators in nonlinear medium have the following form: $`\widehat{A}(t,z)`$ $`=`$ $`e^{i\gamma q[\widehat{n}_0(t)]}\widehat{A}_0(t).`$ (21) $`\widehat{A}^+(t,z)`$ $`=`$ $`\widehat{A}_0^+(t)e^{i\gamma q[\widehat{n}_0(t)]}.`$ (22) In (21) and (22) $`\widehat{A}_0(t)=\widehat{A}(t,0)`$, $`\gamma =\beta z`$. The expression $`q[\widehat{n}_0(t)]`$ (see (20)) it is convenient to be written as: $$q[\widehat{n}_0(t)]=_{\mathrm{}}^{\mathrm{}}h(t_1)\widehat{n}_0(tt_1)dt_1,(h(t)=H(\|t\|).$$ (23) If we consider $`\widehat{n}`$ to be time-independent, then the expressions (21) and (22) give results for the monochromatic field. In the case that the nonlinear response in (21), (22) has the form $`h(t)=\delta (t)`$, the results for non-inertial nonlinear media can be obtained (see (5)). To find the statistical characteristics of a pulse at the output of the nonlinear medium it is necessary to estimate the averages of the operators $`\widehat{A}(t,z)`$, $`\widehat{A}^+(t,z)`$ and their combinations. They can be estimated, if the operator expressions are given in the normally ordered form, when the creation photons Bose-operators are placed at the left of the annihilation photons Bose-operators. The use of the expressions (21), (22) involves the development of a special mathematical device. ## III THE ALGEBRA OF THE TIME-DEPENDENT BOSE-OPERATORS For the beginning, in order to simplify some expressions we introduce the operators: $$\widehat{O}(t)=i\gamma q[\widehat{n}_0(t)],\widehat{O}^+(t)=i\gamma q[\widehat{n}_0(t)],$$ (24) where $`\widehat{O}^+(t)=\widehat{O}(t)`$. Hence, the equations of self-action (21), (22) can be represented as: $`\widehat{A}(t,z)`$ $`=`$ $`e^{\widehat{O}(t)}\widehat{A}_0(t),`$ (25) $`\widehat{A}^+(t,z)`$ $`=`$ $`\widehat{A}_0^+(t)e^{\widehat{O}^+(t)}.`$ (26) Taking into account (23), is it easy to remark that $$[\widehat{O}(t_1),\widehat{O}(t_2)]=0.$$ (27) In consequence we have $`e^{\widehat{O}(t_1)}e^{\widehat{O}(t_2)}`$ $`=`$ $`e^{\widehat{O}(t_2)}e^{\widehat{O}(t_1)}`$ (28) $`=`$ $`e^{\widehat{O}(t_1)+\widehat{O}(t_2)}=e^{\widehat{O}(t_2)+\widehat{O}(t_1)}.`$ (29) ### A THE PERMUTATION OPERATOR RELATIONS The following operator permutation relations hold: $`\widehat{A}_0(t_1)\widehat{O}(t_2)`$ $`=`$ $`[\widehat{O}(t_2)+i\gamma h(t_2t_1)]\widehat{A}_0(t_1),`$ (30) $`\widehat{A}_0(t_1)\widehat{O^+}(t_2)`$ $`=`$ $`[\widehat{O}^+(t_2)i\gamma h(t_2t_1)]\widehat{A}_0(t_1),`$ (31) $`\widehat{A}_0^+(t_1)\widehat{O}(t_2)`$ $`=`$ $`[\widehat{O}(t_2)i\gamma h(t_2t_1)]\widehat{A}_0^+(t_1),`$ (32) $`\widehat{A}_0^+(t_1)\widehat{O^+}(t_2)`$ $`=`$ $`[\widehat{O}^+(t_2)+i\gamma h(t_2t_1)]\widehat{A}_0^+(t_1).`$ (33) Using the mathematical induction it is possible to demonstrate the validity of the formulae ($`mN`$): $`\widehat{A}_0(t_1)\widehat{O}^m(t_2)`$ $`=`$ $`[\widehat{O}(t_2)+i\gamma h(t_2t_1)]^m\widehat{A}_0(t_1),`$ (34) $`\widehat{A}_0(t_1)\widehat{O^+}^m(t_2)`$ $`=`$ $`[\widehat{O}^+(t_2)i\gamma h(t_2t_1)]^m\widehat{A}_0(t_1),`$ (35) $`\widehat{A}_0^+(t_1)\widehat{O}^m(t_2)`$ $`=`$ $`[\widehat{O}(t_2)i\gamma h(t_2t_1)]^m\widehat{A}_0^+(t_1),`$ (36) $`\widehat{A}_0^+(t_1)\widehat{O^+}^m(t_2)`$ $`=`$ $`[\widehat{O}^+(t_2)+i\gamma h(t_2t_1)]^m\widehat{A}_0^+(t_1).`$ (37) To simplify the operator algebra is useful to redefine $$\text{G}(t_2t_1)=i\gamma h(t_2t_1),\text{G}^{}(t_2t_1)=i\gamma h(t_2t_1).$$ (38) Decomposing $`e^{\widehat{O}(t)}`$ and $`e^{\widehat{O}^+(t)}`$ in Taylor series we obtain the operator permutation relations which play an important role at the estimation of the statistical characteristics of a pulse. Hence, finally we have: $`\widehat{A}_0(t_1)e^{\widehat{O}(t_2)}`$ $`=`$ $`e^{\widehat{O}(t_2)+\text{G}(t_2t_1)}\widehat{A}_0(t_1),`$ (39) $`\widehat{A}_0(t_1)e^{\widehat{O}^+(t_2)}`$ $`=`$ $`e^{\widehat{O}^+(t_2)+\text{G}^{}(t_2t_1)}\widehat{A}_0(t_1),`$ (40) $`e^{\widehat{O}(t_1)}\widehat{A}_0^+(t_2)`$ $`=`$ $`\widehat{A}_0^+(t_2)e^{\widehat{O}(t_1)+\text{G}(t_2t_1)},`$ (41) $`e^{\widehat{O}^+(t_1)}\widehat{A}_0^+(t_2)`$ $`=`$ $`\widehat{A}_0^+(t_2)e^{\widehat{O}^+(t_1)+\text{G}^{}(t_2t_1)}.`$ (42) Using the permutation relations (41) it is possible to verify the commutation relation (7) for the operators $`\widehat{A}(t,z)`$ and $`\widehat{A}^+(t,z)`$. ### B THE NORMAL ORDERING THEOREM As pointed out previously, in the considered analyse another important question is represented by the reduction to normally ordered form of the operators $`\widehat{A}(t,z)`$ and $`\widehat{A}^+(t,z)`$. In (23) we proceed to normalized time $`\theta =t_1/\tau _r`$ and then (24) can be presented like: $`\widehat{O}(t)`$ $`=`$ $`i\gamma {\displaystyle _{\mathrm{}}^{\mathrm{}}}\stackrel{~}{h}(\theta )\widehat{n}_0(t\theta \tau _r)𝑑\theta `$ (43) $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{G}(\theta )\widehat{n}_0(t\theta \tau _r)𝑑\theta ,`$ (44) where $`\stackrel{~}{h}(\theta )=\tau _rh(\theta \tau _r)`$ and $`\text{G}(\theta )=i\gamma \stackrel{~}{h}(\theta )`$. For the reduction of the expressions (24) to the normally ordered form the following theorem can be used: Theorem: Bose-operator $`e^{\widehat{O}(t)}`$ can be represented in the normally ordered form this way: $$e^{\widehat{O}(t)}=\widehat{𝐍}\mathrm{exp}\left\{_{\mathrm{}}^{\mathrm{}}\left[e^{G(\theta )}1\right]\widehat{n}_0(t\theta \tau _r)𝑑\theta \right\}.$$ (45) The operators in the integral expression may be understood as the $`c`$-numbers. In the similar theorem in the spectral representation was formulated and demonstrated and we mention only that the similar demonstration can be done also in this case. In fact in the integration limits are not defined and the theorem has not a obvious applicability. The demonstration of (45) does not represent the central objective of this article so we formulate the theorem in the time-representation only. The average value of the $`e^{\widehat{O}(t)}`$ is given by the formula $$e^{\widehat{O}(t)}=\mathrm{exp}\left\{_{\mathrm{}}^{\mathrm{}}\left[e^{G(\theta )}1\right]\overline{n}_0(t\theta \tau _r)𝑑\theta \right\}.$$ (46) ### C THE AVERAGE VALUES OF $`e^{\widehat{O}(t)}`$ AND $`e^{\widehat{O}^+(t)}`$ In most of the experimental situations $`\gamma 1`$ which allows one to decompose the integral expression in the (46) and to limit decomposition at terms having order $`\gamma ^2`$. Using this approach we have: $`e^{\widehat{O}(t)}=\mathrm{exp}`$ $`[{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{G}(\theta )\overline{n}_0(t\theta \tau _r)d\theta `$ (48) $`+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\text{G}^2(\theta )\overline{n}_0(t\theta \tau _r)d\theta ],`$ where $`\overline{n}_0(t)=\widehat{n}_0(t)=\|\alpha _0(t)\|^2`$. It is convenient to enter in further analyse the envelope of a pulse $`\rho (t)`$, so that $`\alpha _0(t)=\rho (t)\alpha _0`$. If the initial pulse has the gaussian form then $`\rho (0)=1`$. For simplicity we denote: $`\psi (t)=\gamma {\displaystyle _{\mathrm{}}^{\mathrm{}}}\stackrel{~}{h}(\theta )\overline{n}_0(t\theta \tau _r)𝑑\theta ,`$ (49) $`\mu (t)={\displaystyle \frac{1}{2}}\gamma ^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}\stackrel{~}{h}^2(\theta )\overline{n}_0(t\theta \tau _r)𝑑\theta .`$ (50) From (48) we have $$e^{\widehat{O}(t)}=e^{i\psi (t)\mu (t)},e^{\widehat{O}^+(t)}=e^{i\psi (t)\mu (t)}.$$ (51) The parameters $`\psi (t)`$ and $`\mu (t)`$ are connected with the self-action effect and $`\psi (t)`$ represents nonlinear phase addition. Then $`\psi (t)=\psi _0{\displaystyle _0^{\mathrm{}}}\stackrel{~}{h}(\theta )\rho ^2(t\theta \tau _r)𝑑\theta ,`$ (52) $`\mu (t)=\mu _0{\displaystyle _0^{\mathrm{}}}\stackrel{~}{h}^2(\theta )\rho ^2(t\theta \tau _r)𝑑\theta ,`$ (53) where $`\psi _0=2\gamma \alpha _0^2=2\gamma \overline{n}_0`$ and $`\mu _0=\gamma ^2\overline{n}_0=\gamma \psi _0/2`$. A special interest is represented by the estimation of the average values of the Bose-operator combinations at coherent initial states (see (24)). Taking into account (24) we have: $$e^{\widehat{O}(t_1)}e^{\widehat{O}(t_2)}=e^{\widehat{O}(t_2)}e^{\widehat{O}(t_1)}=e^{\widehat{O}(t_1)+\widehat{O}(t_2)}=e^{\widehat{Q}(t_1,t_2)}.$$ (54) In consequence we find $$\widehat{𝐍}\left[e^{\widehat{O}(t_1)}e^{\widehat{O}(t_2)}\right]=\widehat{𝐍}\left[e^{\widehat{Q}(t_1,t_2)}\right],$$ (55) where: $`\widehat{Q}(t_1,t_2)`$ $`=`$ $`\widehat{O}(t_1)+\widehat{O}(t_2)`$ (56) $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\stackrel{~}{\text{G}}(t_1,t_2;\theta )\widehat{n}_0(\theta \tau _r)𝑑\theta ,`$ (57) $`\stackrel{~}{\text{G}}(t_1,t_2;\theta )`$ $`=`$ $`i\gamma \stackrel{~}{h}(t_1\theta \tau _r)+i\gamma \stackrel{~}{h}(t_2\theta \tau _r).`$ (58) Using the theorem of normal ordering for (57) we estimate the averages of different combinations of Bose-operators: $`e^{\widehat{O}(t_1)+\widehat{O}(t_2)}`$ $`=`$ $`e^{i[\psi (t_1)+\psi (t_2)]\mu (t_1,t_2)\text{K}(t_1,t_2)},`$ (59) $`e^{\widehat{O}^+(t_1)+\widehat{O}(t_2)}`$ $`=`$ $`e^{i[\psi (t_1)+\psi (t_2)]\mu (t_1,t_2)+\text{K}(t_1,t_2)},`$ (60) $`e^{\widehat{O}(t_1)+\widehat{O}^+(t_2)}`$ $`=`$ $`e^{i[\psi (t_1)\psi (t_2)]\mu (t_1,t_2)+\text{K}(t_1,t_2)},`$ (61) $`e^{\widehat{O}^+(t_1)+\widehat{O}^+(t_2)}`$ $`=`$ $`e^{i[\psi (t_1)+\psi (t_2)]\mu (t_1,t_2)\text{K}(t_1,t_2)},`$ (62) where $`\mu (t_1,t_2)=\mu (t_1)+\mu (t_2)`$ and $`\text{K}(t_1,t_2)`$ represents the temporal correlator $$\text{K}(t_1,t_2)=\mu _0_{\mathrm{}}^{\mathrm{}}\stackrel{~}{h}(t_1\theta \tau _r)\stackrel{~}{h}(t_2\theta \tau _r)\rho ^2(\theta \tau _r)𝑑\theta .$$ (63) In agreement with most of the experimental situations the approximation $`\tau _p\tau _r`$ can be used. In this case in (52),(53) and (63) we can suppose that in time the envelope of a pulse slowly change itself so it practically does not depend on the change of the integration variable. Therefore it is possible to eliminate it from the under integral expression in the essential point $`\theta \tau _r=0`$ in (52,53) and $`\theta \tau _r=t_1+\tau /2`$ in (63) which corresponds to the maximal value of under integral expression $`\stackrel{~}{h}(\theta )=1`$ (52,53) and $`\stackrel{~}{h}(t_1\theta \tau _r)\stackrel{~}{h}(t_2\theta \tau _r)=\stackrel{~}{h}(\tau /2)\stackrel{~}{h}(\tau /2)=\stackrel{~}{h}^2(\tau /2)`$ (63) consequently. Then $`\psi (t)`$ $`=`$ $`\psi _0\rho ^2(t){\displaystyle _0^{\mathrm{}}}\stackrel{~}{h}(\theta )𝑑\theta ,`$ (64) $`\mu (t)`$ $`=`$ $`\mu _0\rho ^2(t){\displaystyle _0^{\mathrm{}}}\stackrel{~}{h}^2(\theta )𝑑\theta ,`$ (65) $`\text{K}(t_1,t_2)`$ $`=`$ $`\mu _0\rho ^2(t_1+\tau /2){\displaystyle \frac{1}{\tau _r}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\stackrel{~}{h}(\theta )\stackrel{~}{h}(\theta +\tau )𝑑\theta .`$ (66) Here $`\tau =t_2t_1`$. One should note that in our previous analyse we did not choose the relaxation function of nonlinearity in a definite form. If the nonlinearity is of a Kerr type, the relaxation function should be introduced as $$H(t)=(1/\tau _r)e^{t/\tau _r}(t0).$$ (67) Then $`\stackrel{~}{h}(\theta )=e^{\|\theta \|}`$ and for the integrals in (64)-(66) we find $$_0^{\mathrm{}}\stackrel{~}{h}(\theta )𝑑\theta =1,_0^{\mathrm{}}\stackrel{~}{h}^2(\theta )𝑑\theta =\frac{1}{2},$$ (68) $$g(\tau )=\frac{1}{\tau _r}_{\mathrm{}}^{\mathrm{}}\stackrel{~}{h}(\theta )\stackrel{~}{h}(\theta +\tau )d\theta =\frac{1}{\tau _r}(1+\frac{\|\tau \|}{\tau _r})\stackrel{~}{h}(\frac{\tau }{\tau _r}).$$ (69) ## IV THE CORRELATION FUNCTION OF QUADRATURES As stated earlier (see (6)) in self-action process the photon statistics remains unchanged. Therefore, we are interested in analysing the quadrature components which are defined as: $`\widehat{X}(t,z)`$ $`=`$ $`\left[\widehat{A}(t,z)+\widehat{A}^+(t,z)\right]/2,`$ (70) $`\widehat{Y}(t,z)`$ $`=`$ $`\left[\widehat{A}(t,z)\widehat{A}^+(t,z)\right]/2i.`$ (71) The averages of the operators $`\widehat{A}(t,z)`$ and $`\widehat{A}^+(t,z)`$ at initial coherent state of a pulse are $`\widehat{A}(t,z)`$ $`=`$ $`\alpha _0(t)e^{\widehat{O}(t)},`$ (72) $`\widehat{A}^+(t,z)`$ $`=`$ $`\alpha _0^{}(t)e^{\widehat{O}^+(t)}.`$ (73) Taking into account (51) and that $`\alpha _0(t)=\|\alpha _0(t)\|e^{i\phi (t)}`$, for average values of quadratures we obtain: $`\widehat{X}(t,z)`$ $`=`$ $`\|\alpha _0(t)\|e^{\mu (t)}\mathrm{cos}\mathrm{\Phi }(t),`$ (74) $`\widehat{Y}(t,z)`$ $`=`$ $`\|\alpha _0(t)\|e^{\mu (t)}\mathrm{sin}\mathrm{\Phi }(t).`$ (75) where $`\mathrm{\Phi }(t)=\psi (t)+\phi (t)`$. Exponential term in (74) and (75) is caused by quantum effects - in the classical theory it is not present. From (74)-(75) it is concluded that the changes of quadratures in time are connected with changes in pulse’s envelope. We introduce correlation functions of quadrature components as $`R_X(t,t+\tau )={\displaystyle \frac{1}{2}}`$ $`[`$ $`\widehat{X}(t,z)\widehat{X}(t+\tau )+\widehat{X}(t+\tau ,z)\widehat{X}(t,z)`$ (76) $``$ $`2\widehat{X}(t,z)\widehat{X}(t+\tau ,z)],`$ (77) $`R_Y(t,t+\tau )={\displaystyle \frac{1}{2}}`$ $`[`$ $`\widehat{Y}(t,z)\widehat{Y}(t+\tau )+\widehat{Y}(t+\tau ,z)\widehat{Y}(t,z)`$ (78) $``$ $`2\widehat{Y}(t,z)\widehat{Y}(t+\tau ,z)].`$ (79) To analyse the correlation functions of quadrature components, it is necessary to evaluate the correlators $`\xi _X(t_1,t_2)=\widehat{X}(t_1)\widehat{X}(t_2)`$ and $`\xi _Y(t_1,t_2)=\widehat{Y}(t_1)\widehat{Y}(t_2)`$. Using permutation relations (41) and (61), we obtain $`\xi _X(t_1,t_2)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\delta (t_2t_1)+{\displaystyle \frac{1}{2}}\|\alpha _0(t_1)\|\|\alpha _0(t_2)\|e^{\mu (t_1,t_2)}`$ (82) $`\times [e^{\mathrm{\Lambda }(t_1,t_2)}\mathrm{cos}[\mathrm{\Phi }(t_1)+\mathrm{\Phi }(t_2)+\gamma \stackrel{~}{h}(t_2t_1)]`$ $`+e^{\mathrm{\Lambda }(t_1,t_2)}\mathrm{cos}[\mathrm{\Phi }(t_1)\mathrm{\Phi }(t_2)]],`$ $`\xi _Y(t_1,t_2)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\delta (t_2t_1){\displaystyle \frac{1}{2}}\|\alpha _0(t_1)\|\|\alpha _0(t_2)\|e^{\mu (t_1,t_2)}`$ (85) $`\times [e^{\mathrm{\Lambda }(t_1,t_2)}\mathrm{cos}[\mathrm{\Phi }(t_1)+\mathrm{\Phi }(t_2)+\gamma \stackrel{~}{h}(t_2t_1)]`$ $`e^{\mathrm{\Lambda }(t_1,t_2)}\mathrm{cos}[\mathrm{\Phi }(t_1)\mathrm{\Phi }(t_2)]],`$ where $`\mathrm{\Lambda }(t_1,t_2)=\mu (t_1,t_2)\stackrel{~}{h}(t_2t_1)`$ (see (58)). As a result for the correlation functions of quadratures we have: $`R_X(t,t`$ $`+`$ $`\tau )={\displaystyle \frac{1}{4}}[\delta (\tau )`$ (86) $``$ $`\psi _0\rho (t)\rho (t+\tau )h(\tau )\mathrm{sin}[\mathrm{\Phi }(t)+\mathrm{\Phi }(t+\tau )]`$ (87) $`+`$ $`\psi _0^2\rho (t)\rho (t+\tau )g(t,\tau )\mathrm{sin}\mathrm{\Phi }(t)\mathrm{sin}\mathrm{\Phi }(t+\tau )],`$ (88) $`R_Y(t,t`$ $`+`$ $`\tau )={\displaystyle \frac{1}{4}}[\delta (\tau )`$ (89) $`+`$ $`\psi _0\rho (t)\rho (t+\tau )h(\tau )\mathrm{sin}[\mathrm{\Phi }(t)+\mathrm{\Phi }(t+\tau )]`$ (90) $`+`$ $`\psi _0^2\rho (t)\rho (t+\tau )g(t,\tau )\mathrm{cos}\mathrm{\Phi }(t)\mathrm{cos}\mathrm{\Phi }(t+\tau )],`$ (91) where $$g(t,\tau )=\rho ^2(t+\tau /2)g(\tau ).$$ (92) To obtain (88) and (91) the $`\gamma 1`$ and $`\tau _r\tau _p`$ approximations have been used. ## V THE SPECTRUM OF QUANTUM FLUCTUATIONS OF QUADRATURE COMPONENTS Spectral densities of fluctuations of the quadratures are defined by the following expressions: $`S_X(\omega ,t)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}R_X(t,t+\tau )e^{i\omega \tau }𝑑\tau ,`$ (93) $`S_Y(\omega ,t)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}R_Y(t,t+\tau )e^{i\omega \tau }𝑑\tau .`$ (94) Taking into account the weak change of the envelope during the relaxation time one obtains: $`S_X(\omega ,t)`$ $`=`$ $`{\displaystyle \frac{1}{4}}[1\psi _0\rho ^2(t)\mathrm{sin}2\mathrm{\Phi }(t){\displaystyle _{\mathrm{}}^{\mathrm{}}}h(\tau )e^{i\omega \tau }d\tau `$ (96) $`+\psi _0^2\rho ^4(t)\mathrm{sin}^2\mathrm{\Phi }(t){\displaystyle _{\mathrm{}}^{\mathrm{}}}g(\tau )e^{i\omega \tau }d\tau ],`$ $`S_Y(\omega ,t)`$ $`=`$ $`{\displaystyle \frac{1}{4}}[1+\psi _0\rho ^2(t)\mathrm{sin}2\mathrm{\Phi }(t){\displaystyle _{\mathrm{}}^{\mathrm{}}}h(\tau )e^{i\omega \tau }d\tau `$ (98) $`+\psi _0^2\rho ^4(t)\mathrm{cos}^2\mathrm{\Phi }(t){\displaystyle _{\mathrm{}}^{\mathrm{}}}g(\tau )e^{i\omega \tau }d\tau ].`$ The estimation of integrals in (96),(98) gives us $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}h(\tau )e^{i\omega \tau }𝑑\tau `$ $`=`$ $`{\displaystyle \frac{2}{1+(\omega \tau _r)^2}}=2L(\mathrm{\Omega }),`$ (99) $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}g(\tau )e^{i\omega \tau }𝑑\tau `$ $`=`$ $`{\displaystyle \frac{4}{[1+(\omega \tau _r)^2]^2}}=4L^2(\mathrm{\Omega }),`$ (100) where $`\mathrm{\Omega }=\omega \tau _r`$. Hence $`S_X(\mathrm{\Omega },t)={\displaystyle \frac{1}{4}}[1`$ $``$ $`2\psi (t)L(\mathrm{\Omega })\mathrm{sin}2\mathrm{\Phi }(t)`$ (102) $`+4\psi ^2(t)L^2(\mathrm{\Omega })\mathrm{sin}^2\mathrm{\Phi }(t)],`$ $`S_Y(\mathrm{\Omega },t)={\displaystyle \frac{1}{4}}[1`$ $`+`$ $`2\psi (t)L(\mathrm{\Omega })\mathrm{sin}2\mathrm{\Phi }(t)`$ (104) $`+4\psi ^2(t)L^2(\mathrm{\Omega })\mathrm{cos}^2\mathrm{\Phi }(t)],`$ where $`\psi (t)=2\gamma \|\alpha _0(t)\|^2`$. From (102) and (104) follows that the choice of the phase $`\mathrm{\Phi }(t)`$ determines the level of quantum fluctuations lower and higher than the shot-noise level $`S_X(\omega )=S_Y(\omega )=1/4`$, corresponding to the coherent state of the initial pulse. In conformity with the Heisenberg relation the behaviour of the spectrum of the $`X`$-quadrature appears to be moved with a phase $`\pi /2`$ in comparison with $`Y`$-quadrature. In case of an optimal phase of the initial pulse $$\phi _0(t)=\frac{1}{2}\mathrm{arctan}\left[\frac{1}{\psi (t)L(\mathrm{\Omega }_0)}\right]\psi (t)$$ (105) chosen for the frequency $`\mathrm{\Omega }_0=\omega _0\tau _r`$, spectral densities $`S_X(\mathrm{\Omega }_0,t)`$ $`=`$ $`{\displaystyle \frac{1}{4}}[\sqrt{1+\psi ^2(t)L^2(\mathrm{\Omega }_0)}\psi (t)L(\mathrm{\Omega }_0)]^2,`$ (106) $`S_Y(\mathrm{\Omega }_0,t)`$ $`=`$ $`{\displaystyle \frac{1}{4}}[\sqrt{1+\psi ^2(t)L^2(\mathrm{\Omega }_0)}+\psi (t)L(\mathrm{\Omega }_0)]^2.`$ (107) Eqs.(106),(107) indicate that when the nonlinear phase addition $`\psi (t)`$ increases $`S_X(\mathrm{\Omega }_0,t)`$ monotonously decreases and $`S_Y(\mathrm{\Omega }_0,t)`$ monotonously increases. At any frequency $`\mathrm{\Omega }`$ we have $`S`$ $`{}_{X}{}^{}(\mathrm{\Omega },t)=S_X(\mathrm{\Omega }_0,t)+{\displaystyle \frac{1}{2}}\psi (t)[L(\mathrm{\Omega })L(\mathrm{\Omega }_0)]`$ (108) $`\times `$ $`\{[L(\mathrm{\Omega })+L(\mathrm{\Omega }_0)]\psi (t)[1+(L(\mathrm{\Omega })+L(\mathrm{\Omega }_0))L(\mathrm{\Omega }_0)\psi ^2(t)]`$ (110) $`\times [1+\psi ^2(t)L^2(\mathrm{\Omega })]^{1/2}\},`$ $`S`$ $`{}_{Y}{}^{}(\mathrm{\Omega },t)=S_Y(\mathrm{\Omega }_0,t)+{\displaystyle \frac{1}{2}}\psi (t)[L(\mathrm{\Omega })L(\mathrm{\Omega }_0)]`$ (111) $`\times `$ $`\{[L(\mathrm{\Omega })+L(\mathrm{\Omega }_0)]\psi (t)+[1+(L(\mathrm{\Omega })+L(\mathrm{\Omega }_0))L(\mathrm{\Omega }_0)\psi ^2(t)]`$ (113) $`\times [1+\psi ^2(t)L^2(\mathrm{\Omega })]^{1/2}\},`$ The spectra of $`X`$-quadrature component, calculated by the formula (110), at $`t=0`$ ($`\psi (0)=\psi _0`$) for the cases $`\mathrm{\Omega }_0=0(\omega _0=0)`$, $`\mathrm{\Omega }_0=1(\omega _0=\tau _r^1)`$ are presented in Figs. 1,2 respectively. On Fig. 1 one can see that for $`\omega _0=0`$ spectral density of $`X`$-quadrature component is minimal on frequency $`\omega =0`$ for any values of phase $`\psi _0`$. For $`\omega _00`$ (Fig. 2) and phases $`\psi _0>1`$ the minimum of the fluctuation spectrum of $`X`$-quadrature component lies at frequencies $`\mathrm{\Omega }=1(\omega =\tau _r^1)`$, and for $`\psi _0<1`$ the minimum lies near $`\mathrm{\Omega }0(\omega 0)`$. ## VI THE WIDTH OF THE SPECTRUM OF SQUEEZED QUADRATURE From Fig. 1 one can conclude that the frequency where spectral density of ”-quadrature fluctuations” is lower than the shot-noise level, depends on the nonlinear phase addition $`\psi _0`$. Width of the spectrum below the shot-noise level $`\mathrm{\Delta }\mathrm{\Omega }=\tau _r\mathrm{\Delta }\omega `$ should be defined from $$S_X(\mathrm{\Delta }\mathrm{\Omega },t)=\frac{1}{2}[\frac{1}{4}+S_X(\mathrm{\Omega }_0,t)].$$ (114) Accounting (106,110) for $`\mathrm{\Omega }_0=0`$ from (114) we have: $`2\psi (t)[\psi (t)`$ $``$ $`\sqrt{1+\psi ^2(t)}]L^2(\mathrm{\Delta }\mathrm{\Omega })+2L(\mathrm{\Delta }\mathrm{\Omega })`$ (115) $`+`$ $`\psi (t)\sqrt{1+\psi ^2(t)}\psi ^2(t)1=0.`$ (116) Eq.(116) in $`L(\mathrm{\Delta }\mathrm{\Omega })`$ (see (99,100)) has two solution of which only one is real. Solving the (116) for $`\mathrm{\Delta }\mathrm{\Omega }`$ we get: $`\mathrm{\Delta }\mathrm{\Omega }`$ $`=`$ $`[{\displaystyle \frac{2\psi (t)[\psi (t)\sqrt{1+\psi ^2(t)}]}{1+\sqrt{12\psi (t)[\psi (t)\sqrt{1+\psi ^2(t)}]}}}`$ (118) $`\times {\displaystyle \frac{1}{\sqrt{[\psi (t)\sqrt{1+\psi ^2(t)}\psi ^2(t)1]}}}1]^{1/2}.`$ From (118) it follows that the change of $`\mathrm{\Delta }\mathrm{\Omega }`$ is connected with changes in pulse’s envelope. The frequency band in which the spectral density of the quadrature fluctuations is lower than the shot-noise level depends on the nonlinear phase shift $`\psi (t)`$. The corresponding dependence at $`t=0`$ for $`\omega =0`$ is displayed in Fig. 3. It may be noted that at $`\psi _01`$ width of the spectrum below the shot noise level is one and a half width of the spectral response of nonlinearity. ## VII PHOTON NUMBER SPECTRAL“DENSITY” OF PULSES WITH SELF-PHASE MODULATION The spectral photon number operator is defined by $$\widehat{n}(\omega ,z)=\widehat{a}^+(\omega ,z)\widehat{a}(\omega ,z),$$ (119) where: $`\widehat{a}(\omega ,z)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }\tau _p}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\widehat{A}(t,z)e^{i\omega t}𝑑t,`$ (120) $`\widehat{a}^+(\omega ,z)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }\tau _p}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\widehat{A}^+(t,z)e^{i\omega t}𝑑t.`$ (121) Taking into account (120,121) and (25,26) for photon number spectral “density” (119) we find: $`\overline{n}(\omega ,z)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \tau _p^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\widehat{A}_0^+(t_1)e^{\widehat{O}^+(t_1)+\widehat{O}(t_2)}\widehat{A}_0(t_2)`$ (123) $`\times e^{i\omega (t_2t_1)}dt_1dt_2{\displaystyle \frac{\overline{n}_0}{2\pi }}{\displaystyle \frac{1}{\tau _p^2}}\|𝐈\|^2,`$ where $$𝐈=_{\mathrm{}}^{\mathrm{}}\rho (t)\mathrm{exp}\left\{i\left[\psi _0\rho ^2(t)+\phi (t)+\omega t\right]\right\}𝑑t.$$ (124) The last expression in (123) is written without the account of $`\mu (t_1,t_2)`$ and $`\text{K}(t_1,t_2)`$ (see (61)). If the initial pulse has gaussian form $`\rho (t)=\mathrm{exp}\{t^2/2\tau _p^2\}`$ then, using paraxial approximation ($`\rho ^2(t)1t^2/\tau _p^2`$) in (124) for spectral density (123) we find $$\overline{n}_{class}(\mathrm{\Omega },z)=\frac{\overline{n}_0}{\sqrt{1+4\psi _0^2}}\mathrm{exp}\left[\frac{\mathrm{\Omega }^2}{1+4\psi _0^2}\right],$$ (125) where $`\mathrm{\Omega }=\omega \tau _p`$. From (125) it follows that the spectral density of a pulse with self-phase modulation (SPM-USP) decreases when nonlinear phase addition increases. To calculate (125) in (123) the terms $`\mu (t_1,t_2)`$ and $`K(t_1,t_2)`$ have not been taken into account. In consequence, the spectral density (125) does not depend on relaxation time. If we take the relaxation function of nonlinearity as $$\text{H}(\theta )=\frac{1}{\tau _r}\mathrm{exp}\left\{\frac{\theta ^2}{2\tau _r^2}\right\}$$ (126) then $`\text{K}(t,t+\tau )=\mu _0\rho ^2(t+\tau /2)\mathrm{exp}\{\tau ^2/4\tau _r^2\}`$. Using the expression $`\text{K}(t,t+\tau )`$ in (123) in paraxial approximation we have $`\overline{n}(\mathrm{\Omega },z)`$ $`=`$ $`{\displaystyle \frac{\overline{n}_0}{\left[\left(1+\gamma \psi _0\nu ^2/4\right)^2+4\psi _0^2\right]^{1/2}}}`$ (128) $`\times \mathrm{exp}\left[{\displaystyle \frac{\mathrm{\Omega }^2\left(1+\gamma \psi _0\nu ^2/4\right)}{\left(1+\gamma \psi _0\nu ^2/4\right)^2+4\psi _0^2}}\right],`$ where $`\nu =\tau _p/\tau _r`$. From (128) one can conclude that the spectral density depends on the nonlinear phase addition and on the relation between the pulse duration and the relaxation time of the nonlinearity. ## VIII THE CORRELATION FUNCTION OF SPECTRAL COMPONENTS OF SPM-USPs We introduce the correlation function of different spectral components in the following symmetric form: $`R(\omega _1,\omega _2,z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\widehat{n}(\omega _1,z)\widehat{n}(\omega _2,z)+\widehat{n}(\omega _2,z)\widehat{n}(\omega _1,z)`$ (130) $`2\widehat{n}(\omega _1,z)\widehat{n}(\omega _2,z)].`$ Leaving out the preliminary accounts for (130) we obtain $$R(\omega _1,\omega _2,z)=I_1\delta (\omega _2\omega _1)\frac{1}{2}\psi _0Im\{I_2^{}I_3+I_2I_3^{}\},$$ (131) $`I_1={\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}`$ $`\stackrel{~}{\rho }(t_1,t_2)\mathrm{exp}\{i[\psi ^2(t_1)\psi ^2(t_2)]\}`$ (134) $`\times \mathrm{exp}i[\omega _1t_1\omega _2t_2]dt_1dt_2,`$ $`I_2={\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}`$ $`\stackrel{~}{\rho }(t_1,t_2)\mathrm{exp}\{i[\psi ^2(t_1)+\psi ^2(t_2)]\}`$ (136) $`\times \mathrm{exp}\{i[\omega _1t_1+\omega _2t_2]\}dt_1dt_2,`$ $`I_3={\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}`$ $`\stackrel{~}{q}(t_1,t_2)\mathrm{exp}\{i[\psi ^2(t_1)+\psi ^2(t_2)]\}`$ (138) $`\times \mathrm{exp}\{i[\omega _1t_1+\omega _2t_2]\}dt_1dt_2.`$ $`\stackrel{~}{\rho }(t_1,t_2)=\rho (t_1)\rho (t_2)`$, $`\stackrel{~}{q}(t_1,t_2)=\stackrel{~}{\rho }(t_1,t_2)\stackrel{~}{h}(t_2t_1)`$. If the initial pulse has gaussian form and the relaxation function has the form (126) then in paraxial approximation for correlation function (131) we find $$R(\mathrm{\Omega }_1,\mathrm{\Omega }_2,z)=I_1\delta (\mathrm{\Omega }_1\mathrm{\Omega }_2)\frac{\psi _0}{2\stackrel{~}{\alpha }\sqrt{\stackrel{~}{\beta }}}Im\{\mathrm{\Gamma }(\mathrm{\Omega }_1,\mathrm{\Omega }_2)\},$$ (139) $$I_1=[n_{class}(\mathrm{\Omega }_1)n_{class}(\mathrm{\Omega }_2)]^{1/2}\mathrm{exp}\{i\psi _0\frac{\mathrm{\Omega }_1^2\mathrm{\Omega }_2^2}{1+4\psi _0^2}\}.$$ (140) In (139) are entered the following designations: $$\mathrm{\Omega }=\omega \tau _p,\stackrel{~}{\alpha }=\left[1+4\psi _0^2\right]^{1/2},$$ (141) $$\stackrel{~}{\beta }=[1+2\nu ^24\psi _0^2)^2+16(1+\nu ^2)^2\psi _0^2]^{1/2},$$ (142) $$\mathrm{\Gamma }(\mathrm{\Omega }_1,\mathrm{\Omega }_2)=\mathrm{exp}\{G+iS\}+\mathrm{exp}\{E+iF\},$$ (143) $$Im\mathrm{\Gamma }(\mathrm{\Omega }_1,\mathrm{\Omega }_2)=e^G\mathrm{sin}S+e^E\mathrm{sin}F,$$ (144) $`G`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_1^2+\mathrm{\Omega }_2^2}{2\stackrel{~}{\alpha }}}\mathrm{cos}ϵ{\displaystyle \frac{\mathrm{\Omega }_2^2}{2\varrho }}\mathrm{cos}\mathrm{\Sigma }{\displaystyle \frac{\varrho \mathrm{\Omega }_1^2}{2\stackrel{~}{\beta }}}\mathrm{cos}(\mathrm{\Sigma }\xi )`$ (147) $`{\displaystyle \frac{\mathrm{\Omega }_1\mathrm{\Omega }_2}{\stackrel{~}{\beta }}}\nu ^2\mathrm{cos}\xi {\displaystyle \frac{\mathrm{\Omega }_2^2\nu ^4}{2\varrho \stackrel{~}{\beta }}}\mathrm{cos}(\mathrm{\Sigma }+\xi ),`$ $`S`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_1^2+\mathrm{\Omega }_2^2}{2\stackrel{~}{\alpha }}}\mathrm{sin}ϵ+{\displaystyle \frac{\mathrm{\Omega }_2^2}{2\varrho }}\mathrm{sin}\mathrm{\Sigma }{\displaystyle \frac{\varrho \mathrm{\Omega }_1^2}{2\stackrel{~}{\beta }}}\mathrm{sin}(\mathrm{\Sigma }\xi )`$ (149) $`+{\displaystyle \frac{\mathrm{\Omega }_1\mathrm{\Omega }_2}{\stackrel{~}{\beta }}}\nu ^2\mathrm{sin}\xi +{\displaystyle \frac{\mathrm{\Omega }_2^2\nu ^4}{2\varrho \stackrel{~}{\beta }}}\mathrm{sin}(\mathrm{\Sigma }+\xi )(ϵ+{\displaystyle \frac{\xi }{2}}),`$ $`E`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_1^2+\mathrm{\Omega }_2^2}{2\stackrel{~}{\alpha }}}\mathrm{cos}ϵ{\displaystyle \frac{\mathrm{\Omega }_1^2}{2\varrho }}\mathrm{cos}\mathrm{\Sigma }{\displaystyle \frac{\varrho \mathrm{\Omega }_2^2}{2\stackrel{~}{\beta }}}\mathrm{cos}(\mathrm{\Sigma }\xi )`$ (151) $`{\displaystyle \frac{\mathrm{\Omega }_1\mathrm{\Omega }_2}{\stackrel{~}{\beta }}}\nu ^2\mathrm{cos}\xi {\displaystyle \frac{\mathrm{\Omega }_1^2\nu ^4}{2\varrho \stackrel{~}{\beta }}}\mathrm{cos}(\mathrm{\Sigma }+\xi ),`$ $`F`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_1^2+\mathrm{\Omega }_2^2}{2\stackrel{~}{\alpha }}}\mathrm{sin}ϵ+{\displaystyle \frac{\mathrm{\Omega }_1^2}{2\varrho }}\mathrm{sin}\mathrm{\Sigma }{\displaystyle \frac{\varrho \mathrm{\Omega }_2^2}{2\stackrel{~}{\beta }}}\mathrm{sin}(\mathrm{\Sigma }\xi )`$ (153) $`+{\displaystyle \frac{\mathrm{\Omega }_1\mathrm{\Omega }_2}{\stackrel{~}{\beta }}}\nu ^2\mathrm{sin}\xi +{\displaystyle \frac{\mathrm{\Omega }_1^2\nu ^4}{2\varrho \stackrel{~}{\beta }}}\mathrm{sin}(\mathrm{\Sigma }+\xi )(ϵ+{\displaystyle \frac{\xi }{2}}),`$ $`\varrho `$ $`=`$ $`\left[(1+\nu ^2)^2+4\psi _0^2\right]^{1/2},`$ (154) $`ϵ`$ $`=`$ $`\mathrm{arctan}\left\{2\psi _0\right\},`$ (155) $`\xi `$ $`=`$ $`\mathrm{arctan}\left\{{\displaystyle \frac{4(1+\nu ^2)\psi _0}{1+2\nu ^24\psi _0^2}}\right\},`$ (156) $`\mathrm{\Sigma }`$ $`=`$ $`\mathrm{arctan}\left\{{\displaystyle \frac{2\psi _0}{1+\nu ^2}}\right\}.`$ (157) We define the spectral correlation function of the photons with frequency $`\mathrm{\Omega }`$ in the spectral band $`\mathrm{\Delta }\mathrm{\Omega }`$ $`\stackrel{~}{R}_{\mathrm{\Delta }\mathrm{\Omega }}(\mathrm{\Omega },z)`$ $`=`$ $`{\displaystyle _{\mathrm{\Omega }\mathrm{\Delta }\mathrm{\Omega }/2}^{\mathrm{\Omega }+\mathrm{\Delta }\mathrm{\Omega }/2}}{\displaystyle _{\mathrm{\Omega }\mathrm{\Delta }\mathrm{\Omega }/2}^{\mathrm{\Omega }+\mathrm{\Delta }\mathrm{\Omega }/2}}R(\mathrm{\Omega }_1,\mathrm{\Omega }_2,z)𝑑\mathrm{\Omega }_1𝑑\mathrm{\Omega }_2`$ (159) $`\overline{n}_0.`$ As a consequence, the conclusion which one can make is that for $`\stackrel{~}{R}(\mathrm{\Omega },z)<0`$ take place the photon antibunching and for $`\stackrel{~}{R}(\mathrm{\Omega },z)>0`$ the photon bunching. The graphic dependence of the spectral correlation function (159) on $`\psi _0`$ at $`\mathrm{\Omega }=0.04`$, $`\mathrm{\Delta }\mathrm{\Omega }=2.510^3`$ and $`\nu =10`$ is displayed in Fig.4, whence it follows that the photon bunching or antibunching can take place, and for phases $`\psi _0>1`$ it becomes significant. At frequency $`\mathrm{\Omega }=0`$, the correlation function (see (159)) has the following simplified form: $$\stackrel{~}{R}_{\mathrm{\Delta }\mathrm{\Omega }}(0,z)=\frac{\psi _0}{2\stackrel{~}{\alpha }\sqrt{\stackrel{~}{\beta }}}(\mathrm{\Delta }\mathrm{\Omega })^2\mathrm{sin}\left(ϵ+\frac{\xi }{2}\right),$$ (160) and its dependence on $`\psi _0`$ at $`\mathrm{\Delta }\mathrm{\Omega }=0.75`$ is shown in Fig. 5, whence it follows that the minimum of the spectral correlation function lies near $`\psi _01`$. In this case the photon antibunching takes place for all phases $`\psi _0>0`$ and it is maximal for $`\psi _01`$. It may be mentioned that the greater is the spectral band of measurement the stronger is the photon bunching or antibunching. ## DISCUSSION AND CONCLUSIONS The results presented in the present paper can be used for the correct interpretation of the results of experiments , in which the laser pulses with the duration of the order $`100`$ ps and quartz optical fibres were used and the maximal meaning of nonlinear phase shift $`\psi _0`$ was greater than $`1`$. Certainly, in the measurement of the quadrature spectrum the suppression of quantum fluctuations of a pulse will be smoothed out (see (110)). This time over which the “smoothing out” occurs in the case of balanced homodyne detection is determined by the duration of the heterodyne pulse. The developed theory enables the choice of the optimal strategy at producing and registration of ultrashort pulses in a squeezed states. The measurement of quantum fluctuation of short pulses take place at high frequencies of the order of several tens MHz in order to avoid any effects due to technical fluctuation concentrated at low frequencies. However, in this area the suppression of quantum fluctuations is greatest. The presented results show that by adjusting the phase of the signal pulse (or the phase of a heterodyne pulse), maximal suppression of the quantum fluctuations can be realized at the spectral component of interest for us. This spectral component of interest can lie on the wing of the spectral response of nonlinearity (Fig. 2). This means, that for obtaining squeezed-light pulses the nonlinear media with a longer relaxation time and consequently with the greater nonlinearity can be used . Our results suggest that in the spectral measurements the photon antibunching can be observed. Usually, in the experiments spectral devices with confined spectral bands are used, thus limiting the amplitude of vacuum fluctuations which participate in the measurements. The final results show that the spectral correlation function depends on the nonlinear phase addition, the relation between the pulse duration and the relaxation time of the nonlinearity and also on the spectral band of the measurement. In consequence the choice of the width of the spectral band of measurement can represent an effective method of control of the photon bunching or antibunching. The obtained results indicate that the photon antibunching can be observed at any value of nonlinear phase addition in the low frequency measurements. At high frequency measurements the photon bunching or antibunching strongly depends on the nonlinear phase additions. We note that the approach developed in the present article can be used to analyse the formation of polarization-squeezed light in media with a cubic nonlinearity. This will be treated in a future publication. ## ACKNOWLEDGMENTS F.P. is grateful to S. Codoban (JINR, Dubna) for useful discussions and rendered help. The work has been performed with partial financial support from Programme ”Fundamental Metrology”.
warning/0003/astro-ph0003285.html	ar5iv	text	# Magellanic-Cloud-Type Interstellar Dust Along Low Density Sightlines in the Galaxy ## 1. Introduction There is an average Milky Way extinction relation, A($`\lambda `$)/A(V), over the wavelength range 0.125 $`\mu `$m to 3.5 $`\mu `$m, which is applicable to a wide range of interstellar dust environments, including lines of sight through diffuse dust, and dark cloud dust, as well as dust associated with star formation (Cardelli, Clayton, & Mathis 1989 (CCM); Cardelli & Clayton 1991; Mathis & Cardelli 1992; Fitzpatrick 1999). The existence of this relation, valid over a large wavelength interval, suggests that the environmental processes which modify the grains are efficient and affect all grains. The CCM relation depends on only one parameter, the ratio of total-to-selective extinction, R<sub>V</sub> which is a crude measure of the size distribution of interstellar dust grains. However, the CCM relation does not appear to apply beyond the Milky Way. It does not always fit the observed extinction along sightlines observed in the Magellanic Clouds and M31 (e.g., Clayton & Martin 1985; Fitzpatrick 1985, 1986; Clayton et al. 1996; Bianchi et al. 1996; Gordon & Clayton 1998; Misselt, Clayton, & Gordon 1999). The 2175 Å bump is weaker and the far-UV extinction is steeper in many of the Magellanic cloud sightlines but there are also sightlines in both the LMC and SMC where the dust extinction does follow CCM. The few lines of sight studied in M31 seem to show a CCM-like far-UV extinction and a weak 2175 Å bump (Bianchi et al. 1996). On the other hand, the starburst nucleus of M33 appears to be associated with Milky-Way-type dust (Gordon et al. 1999). The variations in extinction properties seen in the Magellanic Clouds and M31 may be due to several factors. Different environments, such as star formation regions where large amounts of UV radiation and shocks are present, may play a large role in processing dust. Evidence for this can be seen in the LMC where two distinct wavelength dependences of UV extinction have been found for dust inside and outside the supergiant shell, LMC 2, which lies on the southeast side of 30 Dor. This structure was formed by the combined stellar winds and supernovae explosions from the stellar association at its center (Misselt et al. 1999). In the SMC, the dust properties are even more extreme, showing extinction curves for three of four sightlines which have virtually no bump and are very steep in the far-UV (Gordon & Clayton 1998). Although the dust responsible for these curves is located near regions of star formation in the SMC, the environment is likely to be less severe than for the LMC 2 dust. The 30 Dor region, where LMC 2 is located, is a much larger star forming region than any in the SMC. The dust environments in starburst galaxies and QSOs, which also show SMC-like extinction, are much more extreme than 30 Dor (e.g., Gordon, Calzetti, & Witt 1997; Pitman, Clayton & Gordon 2000; Gordon, Smith & Clayton 2000). The SMC has star formation occurring at only 1% the rate of a starburst galaxy so other factors such as the known differences in metallicity between galaxies may be important (Fitzpatrick 1986; Gordon & Clayton 1998; Misselt et al. 1999). Setting aside global metallicity differences, are there sightlines in the Galaxy where the dust environment is similar to those seen in the Magellanic clouds? Real deviations from CCM are seen in the Galaxy but deviations of the kind seen in the Magellanic clouds have been seen only rarely (Cardelli & Clayton 1991; Mathis & Cardelli 1992). A few sightlines (e.g., 62542, 204827, and 210121) show weak bumps and anomalously strong far-UV extinction for their measured values of $`R_V`$. Their extinction curves are plotted in Figure 1. These deviant sightlines represent a variety of dust environments. The Galactic sightline toward HD 62542 is somewhat similar to LMC 2. Its dust was swept up by bubbles blown by two nearby O stars (Cardelli & Savage 1988). HD 204827 is also in a star formation region where the dust has been subject to shocks (Clayton & Fitzpatrick 1987). HD 210121 lies behind a single cloud in the halo. There is no present activity near this cloud although it was ejected into the halo at some time in the past. There are some important differences between these Galactic extinction curves and those in the Magellanic clouds. The bump seen for HD 62542 is not just weak but it is very broad and shifted to the blue (Cardelli & Savage 1988). Mantles on the bump grains has been suggested as the reason for the weak, broad, and shifted Galactic bumps (Mathis & Cardelli 1992; Mathis 1994). These sightlines show that dust in a variety of environments with a range of $`R_V`$ values can have extinction curves similar to those in the LMC. However, none of the anomalous Galactic sightlines, seen in Figure 1, approach the SMC extinction properties. The SMC dust has weaker bumps and steeper far-UV extinction than any known Galactic or LMC sightline. Most of the Galactic sightlines, that have been studied previously, differ in one respect from the LMC and SMC sightlines. They are significantly more reddened than the Magellanic cloud sightlines. In particular, those sightlines showing the greatest deviations from CCM, those near the supershell LMC 2 and those in the SMC, all have E(B-V) $`<`$ 0.25. Of the twenty-nine CCM sightlines only two have E(B-V) $`<`$ 0.30. The others range up to E(B-V) = 1.2. Similarly, the Fitzpatrick & Massa sample of eighty stars includes only seven with E(B-V) $`<`$ 0.30 (Fitzpatrick & Massa 1990; Fitzpatrick 1999). Therefore, the dust along the Magellanic cloud sightlines is more diffuse and more representative of the warm intercloud medium than the cold cloud medium which is better represented in the Galactic samples. Kiszkurno-Koziej & Lequeux (1987) suggest from ANS extinction measurements of 1200 stars in the Galaxy that there may be a correlation between UV extinction parameters and distance from the Galactic plane. As $`\|z\|`$ increases, the bump becomes weaker and the far-UV extinction stronger. These sightlines have low reddenings and long sightlines so they are also more diffuse and therefore more like those in the Magellanic clouds. To investigate whether the extinction properties observed in the Magellanic clouds are related to the diffuse nature of the sightlines, a sample of long sightlines with low reddenings in the Galaxy was chosen and UV data were obtained with the International Ultraviolet Explorer (IUE) so that extinction curves could be constructed. ## 2. The Sample Sembach, Danks, & Savage (1993) obtained high resolution Na I D and Ca II K spectra for a sample of distant (d $`>`$ 1 kpc) stars. These sightlines were selected to investigate the distribution and physical conditions of gas located in low density regions of the Galactic disk and halo. The sightlines listed in Table 1 were selected from the Sembach et al. sample for a complementary study of the UV extinction properties of interstellar dust in low density conditions. The Sembach et al. sample is limited to stars with spectral types between O8 and B3 which makes them ideal for extinction studies. Following the definitions of Sembach et al. the sightlines in Table 1 lie outside the Sagittarius spiral arm (IA1), between the Sagittarius and Scutum-Crux spiral arms (IA2), beyond the Scutum-Crux arm toward the Galactic center (GC), and the inner 4 kpc of the Galaxy (IGC). The stars lie at distances ranging from 1.5 to 9.5 kpc and have heights above or below the plane of 0 to 1.5 kpc. Twelve of the sightlines extend into the Galactic halo, defined as $`\|z\|`$ $`>`$ 500 pc . The locations of these stars in the Galaxy are plotted in Figure 1 of Sembach et al. (1993). ## 3. Extinction Curves Low dispersion short and long wavelength IUE spectra were obtained between 1991 and 1994 by Jason Cardelli. The spectra listed in Table 2 were downloaded from the IUE archive. The archive spectra were reduced using NEWSIPS and then were recalibrated using the method developed by Massa & Fitzpatrick (2000). The short and long wavelength spectra for each star were co-added, binned to the instrumental resolution ($``$ 5 Å) and merged at the maximum wavelength of the short wavelength spectrum. Extinction curves were constructed using the standard pair method (e.g., Massa, Savage & Fitzpatrick 1983). Uncertainties in the extinction curves contain terms that depend both on the broadband photometric uncertainties as well as those in the IUE fluxes, which are calculated directly in NEWSIPS. Our error analysis is described in detail in Gordon & Clayton (1998). The sample includes early type supergiants which may be used with the same accuracy as main sequence stars in calculating extinction (Cardelli, Sembach, & Mathis 1992). We required $`\mathrm{\Delta }(BV)0.14`$ between the reddened and comparison stars to minimize the uncertainties. The comparison stars have been dereddened as described by Cardelli et al. (1992). Table 1 lists a value of E(B-V) for each star from Sembach et al. (1993). Table 2 lists the $`\mathrm{\Delta }(BV)`$ between the measured (B-V) of the reddened star and the $`(BV)_o`$ of the best-match dereddened comparison star. The extinction curves for the sample stars are shown in Figure 2. The extinction curves have been fitted using the Fitzpatrick & Massa (1990, hereafter FM) parameterization. They have developed an analytical representation of the shape of the extinction curves using a small number of parameters. This was done using linear combinations of a Drude bump profile, $`D(x;\gamma ,x_o)`$, a linear background and a far-UV curvature function, $`F(x)`$, where $`x=\lambda ^1`$. There are 6 parameters determined in the fit: The strength, central wavelength, and width of the bump, $`c_3`$, $`x_o`$, and $`\gamma `$, the slope and intercept of the linear background, $`c_1`$ and $`c_2`$, and the strength of the far-UV curvature, $`c_4`$. The FM fits to individual extinction curves are plotted in Figure 2 and the best fit parameters for each curve are given in Table 3. Near-infrared photometry exists for a few of the reddened stars in our sample. For these stars, using JHK photometry, we calculated values of R<sub>V</sub>. Due the small values of E(B-V) for our sample stars, the uncertainties in R<sub>V</sub> are relatively large. Within these uncertainties, most are consistent with the typical diffuse dust value of 3.1. There is no trend with position on the sky discernible with the small number of sightlines having measured $`R_V`$ values. ## 4. Discussion The sightlines in our sample cover very long distances and have relatively low reddenings so the average densities are amongst the lowest known (Sembach et al. 1993). The measured values for $`n_o`$(H I) (= N(H I)sin$`\|b\|`$/h(H I) where h(H I) is the scale height of H I) are listed in Table 1. The parameter, $`n_o`$(H I), is a measure of average density along a sightline (Sembach et al. 1993). Typically, when $`n_o`$(H I) $`<`$ 0.42 $`cm^3`$, no large cold clouds are present along the line of sight (Sembach et al. 1993). All the stars in Table 1 satisfy this criterion. In addition, the warm intercloud medium dominates over the diffuse cold cloud medium if $`n_o`$(H I) $`<`$ 0.2 $`cm^3`$. Most of the stars in Table 1 satisfy this criterion or come close to it. The ratio of the column density of Ca II to Na I, also listed in Table 1, is another measure of the relative contributions of cloud and intercloud medium. Na I is relatively stronger in clouds while Ca II is relatively strong in the diffuse ISM. This is due to the strong variation in the calcium depletion from the gas phase into dust grains (Sembach & Danks 1994; Crinklaw, Federman, & Joseph 1994). The depletion is higher inside clouds and lower outside where the harsher environment including sputtering and grain collisions will return calcium to the gas phase. The wide range in N(Ca II)/N(Na I) indicates that the cloud/inter-cloud fraction varies strongly from one sightline to another. The absorption features along these lines of sight show multiple components indicating that the distribution of gas is patchy. The average number of components or clouds is 1.5 kpc<sup>-1</sup> for Na I and 2.0 kpc<sup>-1</sup> for Ca II (Sembach & Danks 1994). Using these values and the measured reddening, the average E(B-V) per cloud is 0.05 mag which is typical for standard diffuse clouds (Spitzer 1978). Similarly, the average H I column density per cloud is 2.3 x 10<sup>20</sup> cm <sup>-2</sup>. It is likely that most of these sightlines are dominated by warm intercloud medium and have little contribution from the cold cloud medium. Figure 2 shows that even for extremely diffuse sightlines such as those in our sample, most extinction curves still follow CCM with $`R_V`$ = 3.1. However, there is a subsample of these sightlines whose extinction curves show weak bumps and very steep far-UV extinction similar to the Magellanic clouds. These sightlines all lie in one region of the sky in the direction of the Galactic center. This region (Sembach & Danks 1994 (hereafter the SD region)) coincides with an area of Galactic longitude, l=325<sup>o</sup> to 0<sup>o</sup>, where large forbidden velocities have been observed in the gas. Figure 3 shows bump strength plotted against far-UV steepness for our sample extinction curves. The bump strength can be quantified as the area under the bump using the FM parameters to be $`\pi c_3/2\gamma `$ (Fitzpatrick & Massa 1986). The steepness of the far-UV extinction can be characterized using the FM parameter, $`c_2`$, which represents the slope of the linear background. Seven of the nine sightlines lying in the SD region have the largest values of $`c_2`$ in our sample. Six of these sightlines also have bump strengths below the Galactic average. The sightlines in our sample outside the SD region have values of bump strength and $`c_2`$ grouped around the Galactic average. As can be seen in Figure 3, the SD region extinction parameters fall roughly along a line running from the average Galactic values through the LMC and LMC 2 averages to the those seen in the SMC. This is a promising result as it indicates that whatever factors are affecting the dust properties in the Magellanic clouds may also be affecting the low density dust in the Milky Way. As discussed below, the two SD-region stars which appear in the upper left-hand corner of Figure 3 do not share the extinction properties of the rest of the SD region sample. The locations of the stars in our sample are plotted in Figure 4. We have separated sightlines according to the steepness of the far-UV extinction as measured by $`c_2`$ and also the ratio, N(Ca II)/N(Na I), which measures the relative fraction of cloud and intercloud medium along a sightline. A high value indicates a high fraction of intercloud medium. All of the sightlines in Table 1 with both N(Ca II)/N(Na I) $`>`$ 0.62 and $`c_2>`$ 1.02 lie in the SD region. Figure 5 shows height above or below the Galactic plane plotted against distance in the plane for the stars in the SD region. Figures 4 and 5 clearly show the anomalous extinction is seen only for stars in one particular direction. The two stars in the SD region, that do not share the anomalous extinction of the other sightlines, provide information on the location of this dust. HD 151805 has b= 1.59<sup>o</sup> so its sightline does not go below the Galactic plane as do the other stars in the sample. HD 161653 lies at a distance of 1.8 kpc and is, by far, the closest star in the l=325<sup>o</sup> to 0<sup>o</sup> region. So the dust responsible for the Magellanic-cloud-like extinction lies further than about 2 kpc and in a direction defined by 325$`{}_{}{}^{o}l0^o`$ and -5$`{}_{}{}^{o}b11^o`$. In fact, the four sightlines with the steepest extinction, HD 158243, 160993, 163522, and 164340, most resembling the SMC extinction, lie in an even smaller region bounded by 337$`{}_{}{}^{o}l352^o`$ and -8$`{}_{}{}^{o}b11^o`$. These four stars also have the lowest reddenings. So the remaining three less-extreme sightlines are likely a combination of clouds including more CCM-like dust as well. Kennedy, Bates, & Kemp (1998) give a nice analysis of the absorption components in this direction of the sky including the sightline to HD 163522 (l = 349.6, b = -9.1, d= 9.4 kpc). In their picture, nearby clouds at 50 pc, 100 pc to 1 kpc, the Sagittarius (1.5-2 kpc) and Scutum-Crux (3 kpc) spiral arms all provide absorption components with negative or small positive velocities. They identify peculiar velocity gas corresponding to the forbidden velocities of +10-50 km s<sup>-1</sup> found by Sembach & Danks (1994). These forbidden velocities differ significantly from those expected purely from Galactic rotation. The origin and distance of this gas is not well known. Higher ionization ions are associated with this gas indicating that it may be associated with a Galactic fountain or worm (Savage, Massa, & Sembach 1990; Savage, Sembach, & Cardelli 1994). The extinction curves for the sightlines inside and outside of the SD region are plotted together in Figure 6. Two sightlines, toward HD 151805 and 161653, in the SD region show normal CCM ($`R_V`$ = 3.1) extinction. As discussed above, the dust along those sightlines is not located in the same volume of space with the SD type dust so they are plotted with the non-SD-region sample in Figure 6. The average SD-region curve (again not including HD 151805 and 161653) is plotted in Figure 7 along with the average Milky Way and Magellanic cloud curves for comparison. The average SD-region curve most resembles the HD 62542 and 210121 curves, seen in Figure 1, as well as the LMC 2 curve. The four steepest SD-region curves resemble the SMC Bar curve in the far-UV but still have stronger bump strengths. These curves show little or no slope change from the near to the far-UV while most other curves seen in Figures 1 and 7 tend to turn up steeply to the blue of the bump. The environments of the LMC 2 and HD 62542 dust may be quite similar to the SD region dust. All are in diffuse regions subject to shocks and strong UV radiation fields. The sightline to HD 210121 contains one quiescent cloud with E(B-V) =0.40. However, this cloud is located in the halo and shows UV extinction quite similar to the that seen in the SD region. Calcium is heavily depleted indicating that grain destruction has not been an important mechanism in producing the unusual extinction (Welty & Fowler 1992). The low optical depth of the Magellanic sightlines implies that the dust is not well shielded from these environmental pressures. The typical molecular cloud in the Magellanic clouds is bigger but more diffuse than in the Galaxy (Pak et al. 1998). Then, the small size of the dust grains in a cloud could be the result of the lack of a very dense environment necessary for the grains to grow through coagulation. This has been suggested as the cause of the anomalous extinction seen along the HD 210121 sightline (Larson et al. 1996; 2000). This kind of low density sightline may mimic the conditions in the SMC where extinction properties have been measured over long sightlines with low values of E(B-V) = 0.15-0.24 (Gordon & Clayton 1998). The value of N(Ca II)/N(Na I) is not known for the SMC sightlines but the Ca II abundance in the gas phase is much higher than in the Galaxy for a given reddening (Cohen 1984). The gas-to-dust ratio in general in the SMC is ten times that of the Galaxy (Koornneef 1983). The gas-to-dust ratio in our low density Galactic sample is not significantly different from the average Galactic value (Sembach & Danks 1994). This implies that the gas-to-dust ratio is not well correlated with dust extinction properties. The forbidden velocities seen in our sample are associated with warm intercloud material and the turbulent ISM (Sembach & Danks 1994) and may indicate that the dust may have been subject to shocks. Sembach & Savage (1996) investigated the gas and dust abundances in the halo, finding that they are consistent with progressively more severe processing of grains from the disk into the halo. In addition, they find that while some material is returned to the gas phase, the grain cores seem resistant to destruction. Dust models indicate that the far-UV rise in the extinction curve becomes steeper with increased frequency of exposure to shocks which produces more small dust grains (O’Donnell & Mathis 1997). However, the frequently shocked dust models that produce steeper far-UV extinction also result in stronger bumps. So, producing an SMC-like extinction curve is not as simple as placing dust in a diffuse environment and waiting for a supernova shock. The next logical step is to attempt to directly connect grain properties to their respective sightline environments. Consequently, we have included the average SD region extinction curve in a comprehensive study of dust in the Local Group where our goal is to explicitly examine the correlation of grain size distributions to sightline characteristics such as depletion patterns, and radiation environment (Wolff et al. 2000). ## 5. Conclusions $``$ Magellanic-cloud-like extinction has now been found in the Milky Way. $``$ Large values of N(Ca II)/N(Na I) indicating low depletion are associated with steep far-UV extinction as measured by $`c_2`$. $``$ Global metallicity seems not to be a direct factor. Local environmental conditions seem to be the most important factor in determining dust properties. $``$ Similar UV dust extinction properties have now been seen in the Milky Way, the Magellanic clouds, starburst galaxies and in high redshift star-forming galaxies. $``$ There may be at least two ways to achieve similar extinction properties. A lack of dust coagulation has been suggested for HD 210121 to explain the observed extinction (Larson et al. 1996). The Galactic SD-region properties are closely tied to forbidden velocities indicating that processing of the grains in the diffuse ISM resulted in their observed properties. $``$ There seems to be a correlation between decreasing bump strength and far-UV steepness that includes the Galaxy and the Magellanic Clouds. $``$ All the sightlines contained in CCM lie within 1 kpc of the Sun. As this study shows, dust properties are not well mapped even in our own Galaxy. There are a larger range of UV extinction parameters seen in the Milky Way than implied by CCM. Thanks to Ed Fitzpatrick for providing his IUE IDL procedure. This project was originally envisioned by the late Jason Cardelli.
warning/0003/astro-ph0003416.html	ar5iv	text	# 1 Introduction ## 1 Introduction The evolution of galaxies on cosmological timescales has 3 aspects: the chemical, spectral, and dynamical evolution. In reality, all of them are intimately coupled as shown by the existence of several observational relations involving quantities from any two of these 3 domains. Examples are the Faber – Jackson, Fundamental Plane, or luminosity – metallicity relations for ellipticals or the Tully – Fisher relation and trends of characteristic HII region abundances, average colours, emission line strengths, luminosities, and mass - to - light ratios with spiral type, i.e. with the transition from bulge-dominated to disk-dominated systems. While over short timescales or lookback times and, at least, for giant galaxies, the 3 aspects of galaxy evolution might be treated independently, any modelling of galaxy evolution over cosmological timescales definitely requires a consistent treatment of all 3 aspects. Unfortunately, this is still too complex today, even for numerical simulations on powerful machines. Independently, all 3 aspects of galaxy evolution modelling have quite some tradition already. We are, however, still at the very beginning of a consistent understanding of galaxy formation and evolution over cosmological timescales including all 3 aspects. A consistent modelling of the chemical and dynamical evolution of galaxies is e.g. attempted by the group of G. Hensler in Kiel (‘Chemodynamics’). With the availability of stellar input physics for various metallicities, we developed in Göttingen a consistent description of the chemical and spectral evolution over cosmological timescales, i.e. coupled to a cosmological model. This chemically consistent evolutionary synthesis method and selected applications for the interpretation of galaxy observations over a wide redshift range will be presented here. While still being far from a consistent coupling of our chemically consistent chemical, spectrophotometric and cosmological model with a dynamical galaxy evolution model, we at least started to study the starbursts accompanying strong interactions and mergers involving gas rich galaxies and their impact on the chemical and spectral evolution of galaxies. This brings us to some surprising features accompanying those strong bursts: the formation and evolution both of a new population of bright and massive young star clusters and of a new class of dwarf galaxies forming from recycled material along the tidal features of the merging giant galaxy pair (Tidal Dwarf Galaxies). ## 2 Chemically Consistent Evolutionary Synthesis ### 2.1 Observational Evidence for Broad Metallicity Distributions and Subsolar Metallicities Evidence for extended stellar metallicity distributions and, in particular, for the importance of subsolar abundance contributions has been accumulating over the last years. Stars in elliptical galaxies and bulges, e.g., show a metallicity distribution extending over more than a factor 10 from $`0.06\mathrm{Z}_{}/\mathrm{Z}_{}2.5`$ with average metallicity slightly below solar (Mc William & Rich 1994), G dwarf stars in the solar neighbourhood cover a range from $`0.15\mathrm{Z}_{}/\mathrm{Z}_{}1.5`$ with, again, the average metallicity being subsolar (Rocha-Pinto & Maciel 1996). Characteristic HII region abundances (i.e. measured at 1 $`\mathrm{R}_\mathrm{e}`$) in spirals are observed to be in the range from $`\mathrm{Z}_{}`$ for early type spirals Sa/b through $`\mathrm{Z}_{}`$ for late type spirals Sd (Oey et al. 1993, Zaritsky et al. 1994, Ferguson et al. 1998). Dwarf galaxies of all classes (dEs, dSphs, dIs, BCDs) show – sometimes substantially – subsolar metallicities, both in their stellar populations and in their gas phase abundances (e.g. Richer & McCall 1995). So, while already in the local Universe, subsolar abundances are clearly prevailing on global galaxy scales – with the exception of some central regions of massive luminous galaxies – this is, of course, even more so the case in the early Universe. From restframe UV stellar wind lines, redshifted into the optical for star forming galaxies at redshifts 3 – 4, stellar abundances around 1/10 $`\mathrm{Z}_{}`$ are derived (cf. Lowenthal et al. 1997, Trager et al. 1997). Damped Ly$`\alpha `$ absorbers – most probably the progenitors of present day galactic disks – feature ISM abundances in the range $`10^3<\mathrm{Z}_{\mathrm{ISM}}/\mathrm{Z}_{}<1`$ over a redshift range $`0.4\mathrm{}<\mathrm{z}\mathrm{}<4.5`$ (e.g. Pettini et al. 1997, 1999, see also Sect. 5). These are only a few selected examples from a much longer and still increasing list of direct observational hints to the existence of significant metallicity distributions in galaxies and to the importance of subsolar metallicities. ### 2.2 Method and Input Physics From a very basic consideration it is immediately clear that any stellar system with star formation (SF) extending over more than the lifetime of the most massive stars ($`10^6`$ yr) is composite both in terms of age and metallicity. While star clusters are Simple Stellar Populations (SSPs) consisting of one stellar generation with one age and one metallicity, any galaxy is a Composite Stellar Population (CSP) with its stars spanning finite ranges both in age and metallicity. The Göttingen Evolutionary Synthesis Model starts from a gas cloud of given mass, assumes some Star Formation History (SFH) and Initial Mass Function (IMF) – the basic parameters of this kind of approach – and then, with some book-keeping algorithm, calculates the time evolution of the stellar population across the HRD from a set of stellar evolutionary tracks. With photometric calibrations for colours or absorption indices along all stellar evolutionary tracks the photometric evolution in terms of luminosities, colours, and absorption indices is obtained. Assigning stellar spectra to any point along the tracks and weighting them with the numbers of stars calculated at these points at any given time yields the time evolution of a synthetic galaxy spectrum. Integrating over the synthetic galaxy spectrum with response function for any filter system (e.g. Johnson UBVRIJHK, HST F300W, F450W, …, F814W), luminosities in the respective passbands are obtained in the same way as for an observed galaxy spectrum. In this approach, evolutionary consistency is guaranteed and the time evolution can explicitly be studied. This opens the possibility to directly couple to a cosmological model and study the spectrophotometric evolution as a function of redshift over cosmological timescales (cf. F.-v.A. 1989). Evolutionary synthesis modelling per se accounts for the age distribution of the stars within a galaxy and its time evolution (for a review on synthesis methods see F.-v.A. 1994). To properly account for the metallicity distributions in composite systems like galaxies, the abundance evolution in the gas has to be followed in order to describe successive generations of stars forming out of this gas with stellar evolutionary tracks, yields, and spectra appropriate for their respective initial metallicities, i.e. the gas metallicity at their birth. We call this approach chemically consistent (= cc). A modified form of Tinsley’s equations (Tinsley 1980) with stellar yields for SNII, SNI, PN, and stellar mass loss is solved to obtain the gas content G, the global metallicity Z of the ISM, as well as individual element abundances \[$`\mathrm{X}_\mathrm{i}/\mathrm{H}`$\], and abundance ratios \[$`\mathrm{X}_\mathrm{i}/\mathrm{X}_\mathrm{j}`$\]. We use sets of input physics for 5 different stellar metallicities in the range $`2.3[\mathrm{Fe}/\mathrm{H}]+0.3`$ with logarithmic element abundances with respect to solar defined by \[X<sub>i</sub>/H\] $`:=log(\mathrm{X}_\mathrm{i}/\mathrm{X}_\mathrm{H})log(\mathrm{X}_\mathrm{i}^{}/\mathrm{X}_\mathrm{H}^{}`$). These sets of input physics for the description of both the spectral and chemical evolution of model galaxies comprise stellar evolutionary tracks – covering all relevant evolutionary phases from ZAMS through the PN phase or SN explosion – lifetimes, and remnant masses from the Padova group (Bressan et al. 1993, Fagotto et al. 1994a, b, c). Evolution of low mass stars is taken from Chabrier & Baraffe 1997. Stellar yields for a series of individual elements from <sup>12</sup>C through <sup>56</sup>Fe for massive stars ($`>8\mathrm{M}_{}`$) are from Woosley & Weaver 1995 and yields for intermediate mass stars from van den Hoek & Groenewegen 1997. SNIa contributions to Fe, C, …, are included for the carbon deflagration white dwarf binary scenario as outlined by Matteucci & Greggio 1986. We caution that metallicity dependent stellar yields depend on $`\frac{\mathrm{\Delta }\mathrm{Y}}{\mathrm{\Delta }\mathrm{Z}}`$, explosion energies, remnant masses, etc. and the metallicity dependence of these factors is still poorly understood. SNIa yields are only available for $`\mathrm{Z}_{}`$ (Nomoto et al. 1997, model W7). No important metallicity dependence is expected for SNIa yields except for a possible lower metallicity limit to the explosion (Kobayashi et al. 1998). Model atmosphere spectra, colour and absorption index calibrations are from Lejeune et al. 1997, 1998, and Worthey et al. 1994, respectively. Although this cc approach goes significantly beyond what was possible before, it still is a simplification in the sense that all the input physics is only available for scaled solar abundance ratios. Abundance ratios in galaxies are determined by the metallicity dependent stellar yields, the SFH and IMF, and, in general, will be non-solar. I.e., stellar evolution and galaxy evolution themselves are implicitly coupled and the coupling depends on the SFH and IMF (cf. F.-v.A. 1998b). For a review on chemically consistent evolutionary synthesis see F.-v.A. 1999b. The spectrophotometric evolution of the stellar component of galaxies and the chemical evolution of ISM abundances both are studied not only as a function of time, but also – for any cosmological model as given by $`\mathrm{H}_\mathrm{o},\mathrm{\Omega }_\mathrm{o},\mathrm{\Lambda }_\mathrm{o},`$ and a redshift of galaxy formation $`\mathrm{z}_\mathrm{f}`$ – as a function of redshift. ### 2.3 Model Parameters and Local Templates We use an IMF in the form given by Scalo 1986 with lower and upper mass limits of 0.08 and 85 $`\mathrm{M}_{}`$, respectively, and SFHs, $`\mathrm{\Psi }(\mathrm{t})`$, appropriate for the various spectral types of galaxies. For ellipticals we use a standard $`\mathrm{\Psi }(\mathrm{t})\mathrm{e}^{\mathrm{t}/\mathrm{t}_{}}`$. For spiral types Sa … Sc the Star Formation Rate (SFR) at any timestep is tied to the gas-to-total mass ratio, $`\mathrm{\Psi }(\mathrm{t})\frac{\mathrm{G}}{\mathrm{M}}(\mathrm{t})`$, and for Sd galaxies $`\mathrm{\Psi }(\mathrm{t})=const.`$ is assumed. Characteristic timescales for SF $`\mathrm{t}_{}`$ (for spirals defined via $`_0^\mathrm{t}_{}\mathrm{\Psi }\mathrm{dt}=0.63\mathrm{G}_{\mathrm{t}=0}`$) thus range from 1 Gyr for ellipticals to 2, 3, 10, and 16 Gyr for Sa, Sb, Sc, and Sd galaxies, respectively. The SFHs, together with the IMF, have been chosen as to provide agreement of our model galaxies after a Hubble time of evolution with integrated colours, luminosities, absorption features (E/S0s), emission line strengths (spirals), and template spectra from the UV through NIR (Kennicutt 1992, Kinney et al. 1996, cf. Möller et al. 1998), as well as with characteristic HII region abundances ($`:=`$ measured at the effective radius) typical for the respective galaxy types (Zaritsky et al. 1994, Oey & Kennicutt 1993, van Zee et al. 1998, Ferguson et al. 1998). ### 2.4 Advantages and Shortcomings The combined approach to study the chemical evolution of ISM abundances and the spectrophotometric evolution of the stellar population allows – with the same number of parameters (IMF and SFH) as any single aspect model – for a much larger number of model observables (spectra, luminosities and colours from UV through NIR, emission and absorption line strengths, gas content and metallicity) to be compared to observations. The combination with a cosmological model provides a long redshift or time baseline for comparison with galaxy data, allowing not only to test the models, constrain the parameters, but also to make numerous predictions testable by future observations (F.-v.A. PhD Thesis 1989). The Göttingen Evolutionary Synthesis code has a remarkable analytical potential. It is possible to trace back – at any time and in its time or redshift evolution – the luminosity contribution to any wavelength band of every single stellar mass, of the various spectral types, luminosity classes, and metallicity subpopulations. Ejection rates of every individual element, as well, are monitored for every stellar mass, nucleosynthetic origin (PN, SNII, SNIa), and metallicity subpopulation. The SFHs we use for different spectral types of galaxies are similar to those used by other groups (Bruzual & Charlot 1993, Rocca – Volmerange & Guiderdoni 1988) and meant to be global SFRs – averaged over the entire galaxy and reasonable intervals of time. SFRs fluctuating around our smooth SFHs – either locally or on short timescales – could not be discriminated from their smooth idealisations in integrated galaxy properties after sufficiently long lookback times. Our models are simple 1-zone descriptions without any dynamics or spatial resolution, meant to describe global average quantities like integrated spectra, luminosities, colours, emission and absorption line strengths, or characteristic HII region abundances as measured around the effective radius. While the finite lifetimes of individual stars before they restore their partly enriched material to the ISM are properly accounted for in the chemical modelling (we do NOT use an Instantaneous Recycling Approximation), the mixing of the gas is assumed to happen instantaneously. If not indicated otherwise, models are closed boxes without any inflow (or outflow). While, clearly, real galaxies may not be closed boxes over cosmological timescales, this simplification reflects our ignorance of the time or redshift evolution of gas infall rates and infall abundances. Those cannot yet be taken from hierarchical galaxy formation models that are restricted to Dark Matter. The only consistent way to account for mass (and energy) exchange beween a (proto-)galaxy and its environment seems to be a coupling of our chemo-spectrophotometric evolution with models for cosmological structure and galaxy formation (see Contardo, F.-v. A. & Steinmetz 1998 for a first attempt). ### 2.5 Results for Nearby Galaxies Our cc spectrophotometric evolution models were first presented in Möller et al. 1997 where we discuss the comparison between the (mass-weighted) ISM metallicities for various galaxy types and the luminosity-weighted metallicities of their stellar populations as seen in different wavelength bands. At late stages, stars in models with $`const.`$ SFR (Sd) show a stellar metallicity distribution strongly peaked at $`\frac{1}{2}\mathrm{Z}_{}`$ at all wavelengths and close to the ISM metallicity. Elliptical models, in agreement with observations, show broad stellar metallicity distributions extending from $`\mathrm{Z}=10^4\mathrm{to}\mathrm{Z}=0.05`$ in all bands (cf. Fig. 1 in F.-v. A. 1999b). In collaboration with D. Calzetti (STScI) we are currently working on a consistent inclusion of dust, tying the amount of dust to the evolution of both the gas content and the metallicity in our models, as well as including different spatial distributions of dust and stars in different galaxy types (cf. Möller et al. astro-ph/9906328 for first results). ## 3 CC Spectro – Cosmological Evolution: <br>Selected Results Before the spectrophotometric evolution as a function of redshift can be studied for any given set of cosmological parameters, evolutionary corrections $`\mathrm{e}_\lambda `$ (since a galaxy at $`\mathrm{z}>0`$ is younger) and cosmological corrections – also called k-corrections – $`\mathrm{k}_\lambda `$ (since its spectrum is redshifted) are required. For redshifts $`\mathrm{z}1`$ the attenuation of galaxy light by the cumulative effect of intervening hydrogen has to be taken into account. Intergalactic neutral hydrogen HI largely comes in the form of Ly$`\alpha `$ clouds causing a forest of narrow low column density absorption lines in the featureless continua of background QSOs. The cumulative effect of intergalactic HI distributed stochastically along the line of sight to galaxies at $`\mathrm{z}\mathrm{}>2.5`$ has been shown by Madau (1995) to significantly attenuate the spectra of distant galaxies at restframe wavelengths below 1216 Å. The average attenuation obtained from a large sample of sightlines given by Madau et al. 1996 is applied to our redshifted model galaxy spectra. It additionally weakens their fluxes below restframe $`\lambda =1216`$ Å and its effect is included in our cosmological corrections. The age of a galaxy at redshift $`\mathrm{z}=0`$ is given by $`\mathrm{t}_\mathrm{o}:=\mathrm{t}_{\mathrm{gal}}(\mathrm{z}=0):=\mathrm{t}_{\mathrm{Hubble}}(\mathrm{z}=0)\mathrm{t}_{\mathrm{Hubble}}(\mathrm{z}_\mathrm{f})`$ The red colours of present-day elliptical galaxies require galaxy ages in the range 12 – 15 Gyr. Reasonable combinations of the cosmological parameters are therefore restricted by this minimum age requirement. Model galaxy spectra convolved with filter response functions yield absolute fluxes and magnitudes. Luminosities at $`\mathrm{z}=0`$ are normalised to the average absolute B - band luminosities of the respective galaxy types observed in Virgo (cf. Sandage et al. 1985) before redshifting the synthetic spectra. Apparent magnitudes $`\mathrm{m}_\lambda `$ of high redshift galaxies in any filter $`\lambda `$ are obtained from model absolute magnitudes $`\mathrm{M}_\lambda `$ via the bolometric distance modulus BDM($`\mathrm{H}_\mathrm{o},\mathrm{\Omega }_\mathrm{o},\mathrm{\Lambda }_\mathrm{o})`$, evolutionary and cosmological corrections: $`\mathrm{m}_\lambda (\mathrm{z})=\mathrm{M}_\lambda (\mathrm{z}=0,\mathrm{t}_\mathrm{o})+\mathrm{BDM}(\mathrm{z})+\mathrm{e}_\lambda (\mathrm{z})+\mathrm{k}_\lambda (\mathrm{z})`$ The cosmological or k - correction $`\mathrm{k}_\lambda `$ in any wavelength band $`\lambda `$ accounts for the magnitude difference between a galaxy of age $`\mathrm{t}_\mathrm{o}`$ locally and the same galaxy spectrum redshifted to z $`\mathrm{k}_\lambda (\mathrm{z}):=\mathrm{M}_\lambda (\mathrm{z},\mathrm{t}_\mathrm{o})\mathrm{M}_\lambda (0,\mathrm{t}_\mathrm{o})`$ and can also be calculated for observed galaxy spectra. In this case, the maximum redshift to which this is possible depends on how far into the UV the observed spectrum extends. Our model galaxy spectra at $`\mathrm{z}=0`$ extend from 90 Å through 160 $`\mu `$m and, hence, allow for cosmological corrections in optical bands up to $`\mathrm{z}10`$. Evolutionary corrections $`\mathrm{e}_\lambda `$ account for the age difference between a galaxy of today’s age $`\mathrm{t}_\mathrm{o}`$ and the same galaxy at a younger age $`\mathrm{t}_{\mathrm{gal}}(\mathrm{z})`$ corresponding to the redshift z when its light was emitted $`\mathrm{e}_\lambda (\mathrm{z}):=\mathrm{M}_\lambda (\mathrm{z},\mathrm{t}_{\mathrm{gal}}(\mathrm{z}))\mathrm{M}_\lambda (\mathrm{z},\mathrm{t}_\mathrm{o})`$ Evolutionary corrections, of course, cannot be given without an evolutionary synthesis model. It is important to stress that both the evolutionary and the cosmological corrections do not only depend on the cosmological parameters but also on the SFH or spectral type of the galaxy. In Möller et al. 2000a (in prep.) the redshift evolution of cc galaxy spectra, evolutionary and cosmological corrections, apparent magnitudes, and colours will be presented for two different cosmologies. ### 3.1 CC Models Compared to Solar Metallicity Models A first comparison of the spectrophotometric evolution of model galaxies described in a cc way with those described using only solar metallicity input physics is presented in Möller et al. 1997. The most important result is that cc spiral models, as compared to $`\mathrm{Z}_{}`$ models, are brighter by 1 – 2 mag in B and $`\mathrm{}<1.5`$ mag in K at redshifts $`\mathrm{z}\mathrm{}>1`$. This has significant implications for the interpretation of high redshift galaxy data, e.g. for our understanding of the Faint Blue Galaxy Excess. It means that – contrary to earlier expectations – normal spiral galaxies at $`\mathrm{z}1...\mathrm{}>2`$ can still contribute to galaxy counts around $`\mathrm{B}27`$ mag, and if counts are compared to cc models, less of an excess is expected. The impact of cc modelling on the Luminosity Functions (LFs) in various wavelength bands and on the redshift evolution of LFs of specific galaxy types is currently being explored. ### 3.2 Comparing with Lyman Break Galaxies In any type of galaxy at any time, the flux is close to zero below the Lyman break at 912 Å. While galaxies with passively evolving stellar populations do not have important UV fluxes longward of 912 Å either, those with active SF do show significant UV fluxes longward, while being self-absorbed shortward of $`912`$ Å. At $`\mathrm{z}\mathrm{}>2.5`$ the Lyman break has moved beyond the U - band, causing star forming galaxies to drop out of deep U images while visible in B, V, etc. images. By $`\mathrm{z}\mathrm{}>3.5`$ the Lyman break has moved beyond the V - band, causing star forming galaxies at $`\mathrm{z}\mathrm{}>3.5`$ to also drop out of the V - band while visible at longer wavelengths. Per definitionem only galaxies with active SF are detected with this technique. A passive galaxy at redshift $`\mathrm{z}=3`$, if it existed, would neither be detected in U nor in V or any other optical band. Its restframe flux steeply increasing around 4000 Å only, it would first appear in the NIR H-band. Specific colour criteria have been developed to isolate galaxy candidates at $`2.5\mathrm{z}3.5`$ and at $`3.5\mathrm{z}4.5`$, called Lyman Break Galaxies (LBGs) (see e.g. Steidel et al. 1995, 1999). Follow-up spectroscopy has proven the efficiency of this technique and provided fairly large samples of spectroscopically confirmed galaxies by today: $`>560`$ galaxies at $`\mathrm{z}3`$ and $`>46`$ galaxies at $`\mathrm{z}4`$ both from the Hubble Deep Field and ground based deep imaging surveys. The FORS Deep Field currently observed at the ESO-VLT will soon increase the number of very high redshift galaxies by a significant factor. The nature of the LBGs is controversial. Appearing quite compact with scale radii of 1 – 3 kpc they were suspected to be the progenitors of present-day spheroidal galaxies or bulges (Steidel et al. 1996, Giavalisco et al. 1996, Friaça & Terlevich 1999) or subgalactic fragments (Somerville et al. 1998). Caution seems required, however, comparing rest frame UV scale lengths of LBGs with optical scale lengths of local galaxies, since local galaxies on UV images look very different in overall morphology – not to mention scale lengths – from how they look in the optical (cf. Hibbard & Vacca 1997). The high surface density of LBGs, their SFRs – estimated from UV fluxes to be in the range 3 – 60 $`\mathrm{M}_{}\mathrm{yr}^1`$ –, and luminosities are easier to understand if they are the progenitors of spirals. Comparing the redshift evolution of luminosities and colours of our cc galaxy models with the first set of spectroscopically confirmed HDF dropout galaxies (Lowenthal et al. 1997) in Fig. 1, we find that all these LBGs at $`2\mathrm{}<\mathrm{z}\mathrm{}<3.5`$ are well compatible with normal spiral galaxy progenitors. SFRs derived for LBGs from KECK spectroscopy – corrected for moderate dust extinction and consistent with their non-detection with SCUBA (Chapman et al. 1999) – also agree well with spiral model SFRs in this redshift range. These same results are obtained for two different cosmological models $`(\mathrm{H}_\mathrm{o}=50,\mathrm{q}_\mathrm{o}=0.5)`$ and $`(\mathrm{H}_\mathrm{o}=75,\mathrm{q}_\mathrm{o}=0.05)`$, with $`\mathrm{z}_\mathrm{f}=5`$. The redshift of galaxy formation does not have any visible effect on model luminosities or colours in this redshift range unless $`\mathrm{z}_\mathrm{f}<5`$ for these galaxies. If formed around $`\mathrm{z}5`$, classical initial collapse elliptical galaxy models are significantly brighter and have SFRs $`>20`$ times higher than LBGs. An extensive comparison including all presently available LBGs is in preparation. Clearly, our assumption of one single value for the redshift of galaxy formation is an oversimplification as cosmological structure formation models rather show protracted epochs, extending to very low redshifts for the formation of some galaxy types or masses. No generally accepted specific prediction for the visible parts of galaxies, however, seems to have emerged so far. ## 4 SSPs for Various Metallicities The time evolution of SSPs of different metallicities, i.e. of stellar systems formed in one short burst of star formation ($`\mathrm{t}_{\mathrm{SF}}=10^5`$ yr) with one age and one metallicity, is useful not only for the interpretation of star clusters, but also for combination with various kinds of dynamical galaxy evolution or cosmological galaxy formation models, once a SF criterium is specified. Any kind of extended SFH in a galaxy or some part of it – as complicated as it might be – is readily expanded into a series of single bursts or SSPs. In SSPs (as opposed to galaxies with smooth SFHs), the discreteness of the stellar mass spectrum with evolutionary tracks available causes strong fluctuations in luminosities and colours, that require a posteriori smoothing. To avoid this without interpolating tracks or using isochrones we developed a Monte Carlo Approach. In principle, this method is simple. The HRD population at any timestep is obtained by chosing at random a large number of stellar masses, interpolating their lifetimes, and using for each mass not contained in the original track table appropriate proportions of the two adjacent tracks with each time interval along those two tracks increased or decreased by factors calculated such that the lifetime of the “artificial star” equals the interpolated lifetime. In Kurth et al. 1999, we present results from our Monte Carlo SSPs. Evolution of luminosities UBVRIJHK, colours, and absorption indices is shown for SSPs of 6 different metallicities $`2.3[\mathrm{Fe}/\mathrm{H}]+0.5`$ over age ranges from $`10^7`$ yr to 16 Gyr. While early in the evolution of a star cluster, changes in broad band luminosities and colours are generally rapid, they get weaker and weaker with increasing age. SSP models are compared to observations of globular clusters in the Milky Way and M32 and theoretical calibrations for various indices in terms of \[Fe/H\] are presented in their time evolution. Results are available at http://www.uni-sw.gwdg.de/$`^{}`$ufritze/okurth/ssp.html. Work on the detailed spectral evolution of SSPs from the UV through the NIR is currently in progress. ## 5 Chemically Consistent Chemical Evolution Timmes et al. 1995 and Portinari et al. 1998 used stellar yields for a range of metallicities to model the chemical evolution of the Milky Way and the solar neighbourhood. In the following, we will present some of the results from a comparison of our cc chemical evolution models with observed abundances in the ISM of Damped Ly$`\alpha `$ Absorbers (DLAs). ### 5.1 CC Spiral Models compared to DLAs DLAs show radiation damped Ly$`\alpha `$ absorption lines in the spectra of background QSOs. These damped Ly$`\alpha `$ lines are due to high column density gas ($`\mathrm{log}\mathrm{N}(\mathrm{HI})[\mathrm{cm}^2]20.3`$) and usually accompanied by a large number of low ionisation lines of Al, Si, S, Cr, Mn, Fe, Ni, Zn, … with the same absorption redshift. Only if the often complex velocity structure in the lines can be fully resolved, as e.g. in high-resolution spectroscopy on KECK and WHT with $`\mathrm{\Delta }\lambda /\lambda =\mathrm{60\hspace{0.17em}000}`$, precise element abundances can be derived. These are becoming available for a large number of DLAs over the redshift range 0 … $`4`$ (Boissé et al. 1998, Lu et al. 1993, 1996, Pettini et al. 1994, 1999, Prochaska & Wolfe 1997, …). Based on similarities between their HI column densities and those of local spiral disks, between their comoving gas densities at high z and densities of gas $`+`$ stars in local galaxies, and based on line asymmetries indicative of rotation, DLAs are though to be (proto-)galactic disks along the line of sight to distant QSOs (e.g. Wolfe 1995, Prochaska & Wolfe 1997, 1998, Wolfe & Prochaska 1998). Alternatively, Matteucci et al. 1997 propose starbursting dwarf galaxies, Jimenez et al. 1999 LSB galaxies, and Haehnelt et al. 1998 subgalactic fragments to explain DLA galaxies at low and high redshifts, respectively. For any cosmological model the time evolution of metallicity Z, element abundances \[$`\mathrm{X}_\mathrm{i}/\mathrm{H}`$\], gas content G, and SFR directly transforms into the corresponding redshift evolution. We stress that for a given SFH and IMF our models yield absolute abundances that do not require any scaling or normalisation. After referring all observed DLA abundances to one homogeneous set of oscillator strengths and solar reference values, we compare them with our cc chemo-cosmological models in Lindner et al. 1999. For a standard cosmological model ($`\mathrm{H}_\mathrm{o}=50,\mathrm{\Omega }_\mathrm{o}=1,\mathrm{\Lambda }_\mathrm{o}=0,\mathrm{z}_\mathrm{f}=5`$) it is seen on the example of Zn, which locally is known not to be depleted on dust grains, that our Sa and Sd models bracket the redshift evolution of DLA abundances from $`\mathrm{z}4.4`$ to $`\mathrm{z}0.4`$. The entirety of Zn abundances observed in DLAs fall between an upper envelope provided by the rapidly enriching Sa model and a lower envelope made up by the slowly enriching Sd model. The enrichment evolution of Sb and Sc models run between those for Sa and Sd (cf Fig. 2a). Similar agreement is found for all 8 elements with a reasonable number of DLA abundances, i.e. for Zn, Fe, Si, Cr, Ni, S, Al, Mn. Phillipps & Edmunds 1996 and Edmunds & Phillipps 1997 argue that the probability for an arbitrary QSO sightline to cut through an intervening gas disk and produce DLA absorption is highest around the effective radius where our models were shown to match average HII abundances in the respective spiral types by $`\mathrm{z}=0`$. Hence, our models bridge the gap from high-z DLAs to nearby spirals. We conclude that from the point of view of abundance evolution, DLA galaxies may well be the progenitors of normal spirals Sa – Sd, although we cannot exclude that some starbursting dwarf or LSB galaxies may also be among the DLA galaxy sample. If, at the highest redshifts, damped Ly$`\alpha `$ absorption were caused by subgalactic fragments still bound to merge, our simplified 1-zone models are meant to describe the total SFR of all fragments ending up in one galaxy by low redshift. The weak redshift evolution of DLA abundances is a natural result of the long SF timescales for disk galaxies. The range of SF timescales $`\mathrm{t}_{}`$ for spirals from Sa through Sd fully explains the abundance scatter observed among DLAs at any redshift. Note that our comparison of spiral galaxy models with DLA abundance data extends to redshifts $`\mathrm{z}\mathrm{}>4.4`$, i.e. over lookback times of $`>90\%`$ of the age of the Universe. In comparison with solar metallicity models the influence of the metallicity dependent yields is seen to vary from element to element. Whenever a significant difference is seen, the cc models give better agreement with the data than $`\mathrm{Z}_{}`$ models. Changing the SN explosion energy (yields for model C from Woosley & Weaver) has a minor effects (see also Fig. 7a-h in Lindner et al. 1999). Somewhat surprisingly, elements which locally deplete strongly onto dust grains (like Fe or Cr) are as well described by our models as are non-refractory elements like Zn (cf. Fig. 2b). We hesitate, however, to draw conclusions about the importance of dust in DLAs in view of the uncertainties in the stellar yields and the simplicity of our closed-box models. We also tried models with constant infall rates and primordial infall abundances. In this case, SFHs have to be adjusted as to still give agreement at $`\mathrm{z}=0`$ with average type-specific HII region abundances, colours, etc. This is the reason for moderate constant infall rates to only marginally change the redshift evolution of model abundances and to not affect any of our conclusions. A redshift dependent infall rate, however, with or without evolving infall abundance, might affect our results. The difficulty is to constrain the additional parameter(s) without embedding the galaxy into a cosmological context. Mass estimates from column densities, linear dimensions (DLAs seen in close QSO pairs), and rotation velocities (Prochaska & Wolfe 1997) indicate that – at least some of – these systems at $`\mathrm{z}23`$ already have the masses of local spirals and, hence, may not require important infall. ### 5.2 A Change with Redshift in the DLA Population ? While at high redshift all spiral types Sa through Sd seem to give rise to DLA absorption, no more data points at $`\mathrm{z}\mathrm{}<1.5`$ reach close to our early type spiral models. The boundary to the observed abundances coincides with the 50 % gas-to-total-mass ratio in our models. I.e., by low redshift, the gas poor early type spirals seem to drop out of DLA samples. A deficiency of high N(HI) DLA systems at low z had been noted before (Lanzetta et al. 1997) and attributed to their high metallicity and dust content (Steidel et al. 1997). We give an additional reason: as the global gas content drops, the probability for an arbitrary line of sight to a QSO to cut through a high column density part of the galaxy decreases, i.e. the cross section for damped Ly$`\alpha `$ absorption gets reduced. ### 5.3 Spectroscopic Predictions for DLA Galaxies <br>and Future Prospects Despite considerable efforts and large amounts of telescope time devoted to the optical identification of galaxies responsible for damped Ly$`\alpha `$ absorption, only less than a handful have been found by today. Contrary to earlier expectations, the success rate has not increased when low redshift DLAs were found in UV spectra of QSOs. If, indeed, there were a shift in the DLA galaxy population towards later spiral types at lower redshift, then this is what our models predict since, locally on average, Sd galaxies are fainter by $`2`$ mag in B than Sa’s. Hence, the low-z late type DLA galaxies should be about as faint in B, $``$, and K as the brightest members (i.e. the early spiral types) of the high-z population. Both, an average Sd model at $`\mathrm{z}0.5`$ and an average Sa model at $`\mathrm{z}23`$ have $`\mathrm{B}25`$, $``$$`24.5`$, K $`22`$ mag. Luminosities of the few optically identified DLA galaxies (and candidates) are in good agreement with our predictions. In particular, if there were early type spirals among the low redshift DLA galaxy population they would have $`\mathrm{B}22.5`$, $``$$`21`$, K $`18.5`$ mag at $`\mathrm{z}0.5`$ and they would have been detected. Their non-detection is consistent with our finding of a change with redshift in the DLA galaxy population (see F.-v.A. et al. 1999a, b, c for details). DLA galaxies are within the reach of 10m-class telescopes up to redshifts z $`>3`$ and trace the normal galaxy population to these high redshifts without any bias as to high luminosity, radio power, or the like. With information about ISM abundances from the metal absorption lines and spectrophotometric properties of the stellar population, they will allow for powerful constraints on model parameters, ages, and cosmological parameters (cf. Lindner et al. 1996). Accurate abundance data in very low metallicity DLAs may provide clues for the nucleosynthesis at low metallicity. The knowledge about the abundance evolution of gaseous spiral disks as a function of redshift will allow to predict abundances of stars, star clusters, and Tidal Dwarf Galaxies that form in the powerful starbursts accompanying spiral galaxy mergers in the local universe and in the past (cf. Sect. 6.3 and 6.4). As will be shown, knowing these abundances is important for properly interpreting observed colours and luminosities of star clusters and Tidal Dwarf Galaxies in terms of ages, masses, etc. Spectroscopy of those objects, in turn, provides an independent cross check of our abundance predictions. ## 6 Interactions and Starbursts ### 6.1 Importance of Galaxy Interactions and Starbursts Evidence for the importance of galaxy interactions and merging is coming both from theoretical galaxy formation modelling and observations. In a variety of cosmological contexts (CDM with and without $`\mathrm{\Lambda }`$, mixed DM), semianalytic as well as numerical structure formation models predict galaxies to build up hierarchically from a series of mergers of successively larger masses (see e.g. Carlberg 1990, Lacey & Cole 1993, Kauffmann & White 1993). With increasing redshift, observations indicate increasing numbers of physically close galaxy pairs (e.g. Zepf & Koo 1989), galaxies seen in interaction, and disturbed galaxy morphologies indicative of a recent interaction (e.g. Conselice & Bershady 1999). All strong galaxy interactions or mergers are observed to be accompanied by bursts of star formation if one or both of the galaxies are gas rich. When falling into a cluster, the spiral rich field galaxy population must somehow be transformed into the E/S0/dE rich cluster population. These transformation processes also involve interactions – among individual galaxies or of an infalling galaxy with the cluster environment, potential, or its hot intracluster medium. In gas rich galaxies some of these processes, as well, may trigger starbursts – sometimes powerful and global, as reported for E$`+`$A galaxies by Poggianti & Barbaro 1996. In dissipationless stellardynamical mergers (of disks or spheroidal systems), the remnant after partially complete violent relaxation usually is a spherical system. Its central phase space density cannot significantly increase beyond that of the progenitors. Therefore, it does not seem possible to produce the high central density cores of massive ellipticals by dissipationless mergers of less massive subcomponents or dwarf galaxies which intrinsically have lower central densities. Gaseous mergers, however, are highly dissipative, the gas is driven to the center very efficiently, and gas densities comparable to the central stellar densities in massive ellipticals are reached in high-resolution simulations. In Ultraluminous Infrared Galaxies (ULIRGs), all of which have been shown to be mergers of gas rich spirals, central molecular gas densities of the order of 1000 $`\mathrm{M}_{}\mathrm{pc}^3`$ are indeed observed. It has not been possible yet to treat the full dynamics of mergers involving stars, DM, and gas in at least 3 phases (Hot X gas, HI, cold molecular gas), and to consistently include SF and feedback (under extreme conditions) into high-resolution stellar $`+`$ gaseous dynamical models of interacting gas rich galaxies. It is clear, however, that strong bursts of SF or/and efficient fuelling of a central AGN may occur. Which alternative is chosen or dominates may depend on the properties of the galaxies involved, the geometry of the encounter, and possibly on the stage of the interaction (see e.g. F.-v.A. 1996 for a review in Galaxy Interactions). By studying interaction-triggered starbursts with spectrophotometric and chemical evolutionary synthesis models we hope to learn about the SF process, SF efficiencies, etc. under the violent conditions in mergers/strong interactions that seem to significantly differ from the comparatively “peaceful” SF situation in our Milky Way. ### 6.2 Selected Results We calculated an extensive grid of starburst models. In the idealisation of equal-type mergers Sa-Sa, Sb-Sb, Sc-Sc, Sd-Sd, occuring at a variety of evolutionary stages of the parent galaxies, we model starbursts of various strengths and durations on top of the pre-existing stellar population. The age and type of the progenitor galaxies determine the luminosities and colours of the stellar population and the ISM metallicity at the onset of the burst. The gas content of the spirals at merging sets an upper limit to the gas reservoir available for SF in the burst. Burst strengths being defined as the increase of the stellar mass S during the burst, $`\mathrm{b}:=\frac{\mathrm{\Delta }\mathrm{S}_{\mathrm{burst}}}{\mathrm{S}_{\mathrm{pre}\mathrm{burst}}}`$, this means that e.g. Sa-Sa mergers in the local universe cannot have bursts as strong as Sd-Sd mergers today or Sa-Sa mergers earlier in their evolution when the galaxies contained more gas. The evolution of galaxies before, during, and after the burst was studied in terms of two-colour diagrams, as well as in their luminosity and metallicity evolution for solar metallicity models in F.-v.A. & Gerhard 1994a. The models show that after a burst spiral-spiral mergers of galaxy types Sa, Sb, Sc occuring 3 – 4 Gyr ago readily develop typical elliptical galaxy colours by today – provided SF ceases completely. Sd mergers remain too blue for $`>4`$ Gyr. In detail, the timescale until a merger remnant with its postburst reaches the observed colour range of E/S0s depends on the progenitor types, the burst strength, and on the wavelength range considered. Using our grid of models we analysed the starburst in the prototypical merger remnent NGC 7252, for which a wealth of additional information is available (e.g. spectra from UV – optical, HI-, HII-, CO-maps, dynamical modelling). Despite its enormous tidal tails it already features an azimutally averaged $`\mathrm{r}^{1/4}`$ – profile in its inner part, and it is the oldest from Toomre’s (1977) dynamical age sequence of interacting galaxies. Comparing broad band colours UBVR with our model grid, we were left with some range of possible combinations of burst strengths and ages. Within this cell of parameter space, we in detail compared our model spectra with the spectrum of NGC 7252 showing deep Balmer absorption lines and a small H<sub>β</sub> emission component. This allowed us to identify a very strong starburst that increased the stellar population by $`\mathrm{}<50`$ % and started $`1`$ Gyr ago, i.e. around the time of $`1^{\mathrm{st}}`$ pericenter as obtained from dynamical modelling (cf. Hibbard & Mihos 1995). Both in terms of morphology and colours, NGC 7252 could evolve into an elliptical galaxy within the next 1 – 3 Gyr, provided its SFR $`0`$. Its present centrally concentrated SFR as obtained from an IUE spectrum (F.-v.A. et al. in prep.) is $`3\mathrm{M}_{}\mathrm{yr}^1`$, consistent with the weak H<sub>β</sub> emission. HI falling back from the tidal tails (Hibbard & Mihos 1995) together with gas restored by dying burst stars supports the present SFR and may do so for another few Gyrs, allowing NGC 7252 to rebuild itself a small disk of stars and gas. In this case, both in terms of morphology and colours, it might rather come to resemble an S0 or even an Sa galaxy. In an attempt to identify typical values and ranges for burst parameters and study their possible dependence on galaxy and interaction properties, Liu & Kennicutt’s (1995) sample of interacting galaxies was analysed by O. Kurth in his diploma thesis (Göttingen 1996). Only on the basis of those typical values and relations can interaction induced starbursts consistently be included into the statistical evolution of high redshift galaxies, as e.g. into models for the redshift evolution of the luminosity function or into cosmological galaxy formation scenarios. In general, burst durations can only weakly be constrained to be of the order of a few $`10^8`$ yr for spiral – spiral mergers, i.e. definitely longer than for Blue Compact Dwarf Galaxies (BCDs), where they are of the order of $`10^6`$ yr. In both cases, burst durations reflect the dynamical timescales of the systems involved. One of the most important results is that starbursts in galaxy mergers can be very strong, increasing the stellar mass by as much as 30 % or even more. In all massive interacting systems, burst strengths $`\mathrm{}>10\%`$ were found, larger by a factor $`10100`$ than the burst strengths obtained for BCDs (Krüger et al. 1995). For NGC 7252, the peak SFR during the burst must have been several hundred $`\mathrm{M}_{}\mathrm{yr}^1`$, comparable to those of ULIRGs. Assuming a gas content from the high end of what is typical for Sc galaxies, the SF efficiency (SFE) – defined as the mass of stars formed in the burst relative to the mass of gas available $`\eta :=\frac{\mathrm{\Delta }\mathrm{S}_{\mathrm{burst}}}{\mathrm{G}_{\mathrm{pre}\mathrm{burst}}}`$ – must have been very high, of order 30 – 45 %, a factor $`10`$ larger than in the Milky Way (F.-v.A. & Gerhard 1994b). CO, HCN, and CS observations of ULIRGs tracing molecular gas at densities $`\mathrm{n}500,10^4,10^5\mathrm{cm}^3`$, respectively, explain their high SFEs as a consequence of their high fractions of molecular gas at the highest densities (Solomon et al. 1992). While for BCDs a trend for burst strengths to decrease with increasing galaxy mass was found (Krüger et al. 1995), burst strengths and SFEs in interacting galaxies with masses 1 – 2 orders of magnitude higher are generally higher by the same 1 – 2 orders of magnitude. Does this mean, that the SF process is intrinsically different in different environments ? (See Sect. 5 in F.-v.A. 1996 for a discussion of “violent vs peaceful” SF and the possible relation with differences in the molecular cloud structure). SFEs that high are usually assumed for the collapse of protogalactic clouds in the early Universe that gives rise to globular cluster formation. We speculated that globular cluster formation might have been possible in the starburst of NGC 7252 and predicted the metallicity of both stars and clusters forming in the burst to be $`\mathrm{Z}_{}\mathrm{Z}_{\mathrm{ISM}}^{\mathrm{Sc}(\mathrm{t}_\mathrm{o}1\mathrm{Gyr})}\mathrm{}>\frac{1}{2}\mathrm{Z}_{}`$. ### 6.3 Young Star Clusters in Mergers Within a few years, bright Young Star Cluster (YSC) populations were detected with HST in NGC 7252 and a number of other interacting galaxies and merger remnants. They are interesting for many reasons. E.g., star clusters are better suited than the integrated light to study the age and duration of the starburst. Being most probably observed around the same galactocentric radii at which they formed, they thus allow to study the spatial extent of the starburst and – for sufficiently large cluster populations – its time evolution. YSCs offer the possibility to study star and cluster formation processes in a violent environment, in particular in comparison with observations of the molecular gas content or cloud structure. Using SSP models for the interpretation of YSC systems in the old merger remnant NGC 7252 (F.-v.A. & Burkert 1995) and in the young interacting Antennae galaxies NGC 4038/39 (F.-v.A. 1998a), we show that metallicities of YSCs are crucial to age-date them on the basis of observed broad band colours and to predict their future luminosity and colour evolution. For a few of the brightest YSCs in NGC 7252, our metallicity prediction from spiral models was confirmed by spectroscopy (Schweizer & Seitzer 1993, 1998). The question if some, many, or most of these YSCs could be young Globular Clusters (GCs) is of eminent importance. If a secondary GC population, comparable in number to the original one, can be formed in a merger induced starburst, this would invalidate the last surviving argument against the formation of – at least some – elliptical galaxies from spiral-spiral mergers based on the difference in specific GC frequencies between elliptical galaxies and spirals. Analysing WFPC1 data for some 40 YSCs in NGC 7252 and $`>500`$ in NGC 4038/39 we tentatively conclude that the bulk of the YSC populations in both systems may well evolve into “normal” GC populations in terms of colour distributions, luminosity and mass functions (F.-v.A. 1999a). Before definite conclusions can be reached, we are currently reanalysing WFPC2 data and studying YSCs in an age sequence of interacting, merging and merged galaxies to assess the impact of dynamical cluster destruction processes. Because of its enhanced metallicity, a secondary GC population might eternally testify to a merger origin, still at times when tidal tails, kinematic peculiarities, and fine structure will long have disappeared, and colours will no longer reveal a past starburst. Analyses of GC systems in a number of elliptical galaxies – from low-luminosity ellipticals all through cD galaxies – have revealed bimodal colour distributions in many bright ellipticals, including 2 S0s, broad or multi-peak colour distributions in all cD galaxies investigated, and single-peak distributions in a few low-luminosity ellipticals (Kissler – Patig this volume, Kissler – Patig et al. 1998, Gebhardt & Kissler – Patig 1999). With 10 m telescopes in combination with HST imaging, MOS becomes feasible for YSCs and even old GCs out to Virgo distances. In comparison with SSP models, it will allow to decompose the colour distribution of star cluster systems into metallicity and age distributions and, thereby, give information about the formation of their parent galaxies. ### 6.4 Tidal Dwarf Galaxies Not only secondary populations of YSCs are formed during mergers of gas-rich galaxies, but also a new galaxy formation mechanism has been detected in these systems a few years ago. In the extended tidal tails of several interacting systems bright blue star forming knots, often associated with large HI concentrations, are observed with masses and luminosities typical of dwarf galaxies. Two of these Tidal Dwarf Galaxies (TDGs) are detected in NGC 7252, one in NGC 4038/39, two in Arp 105, several in the “Superantennae”, and other systems (see Duc et al. 1998 for a recent review). TDGs form from “recycled” material torn out from a spiral disk. They deviate from the dwarf galaxy luminosity – metallicity relation in the sense that they have enhanced metallicities for their luminosities as compared to “non-recycled” dwarf galaxies. Abundance determinations from their HII region-like spectra agree well with predictions from our spiral models. While TDGs forming in the lower parts of tidal features are expected to fall back into the merger remnant on timescales of $`10^8`$ to few $`10^9`$ yr, those at the tips of the tails will probably escape and might survive as independent entities. Kinematic independence from the parent galaxy has been demonstrated in a few cases (Duc et al. 1997, Duc & Mirabel 1998). Independently, both N-body simulations for the stars and hydrodynamical models for the gas in merging galaxies feature condensations along their respective tidal tails, leading to two competing scenarios for the formation of TDGs (cf. Barnes & Hernquist 1992, Elmegreen et al. 1993). Evolutionary synthesis models involving fractions of the stellar populations of spiral progenitors plus starbursts of various strengths are used to analyse the stellar populations and the evolutionary states of TDGs. Gaseous continuum and line emission are included. Analysing a small sample of TDGs observed by P.-A. Duc with a coarse grid of models (F.-v.A. et al. 1998) we find evidence for strong bursts, $`\mathrm{b}:=\frac{\mathrm{S}^{\mathrm{young}}}{\mathrm{S}^{\mathrm{old}}}=0.10.4`$, on top of ‘old’ stellar populations, i.e. populations with the age distribution of stars in their parent spiral ($`\mathrm{S}^{\mathrm{young}}\mathrm{and}\mathrm{S}^{\mathrm{old}}`$ being the masses of young and ’old’ stars, respectively). Since these strong bursts completely dominate the light in the optical, NIR data are required to constrain the mass contribution of the ‘old’ population. The latter is important for the fading that TDGs will experience in the future, as well as for their dynamical evolution. Without a significant old population the strongest bursts might disrupt the TDG. Both, the fading and future dynamical evolution, in turn, are relevant for a possible cosmological significance of TDGs. In the past, the merger rate was higher, galaxies were more gas-rich and probably less stable so that the production of a significant population of TDGs might be expected. To explain that part of the faint blue galaxy excess that is due to dwarf galaxies, Babul & Ferguson 1996 invoke a hypothetic population of dwarf galaxies, the formation of which is delayed until $`\mathrm{z}1`$ by the intergalactic radiation field. Properties they require for their dwarf galaxies to explain the faint blue galaxy excess are very similar to those of our TDGs. Depending on their fading, TDGs might well explain part of the faint blue galaxy excess. Moreover, the remnant problem faced by Babul & Ferguson and Ferguson & Babul 1998 would be alleviated if part of the TDGs would spiral back into the merger remnant on timescales of a $`10^810^9`$ yr. In Weilbacher et al. 2000 (submitted) we analyse a sample of 10 interacting galaxies from the AM catalogue with imaging in B, V, R. Comparing colours of knots along tidal structures with an extensive grid of TDG model calculations, a number of promising TDG candidates are identified. Knots with colours not explained by models most probably are background galaxies. TDG candidates in this first sample have young burst ages of $`7`$ Myr, burst strengths $`\mathrm{b}\mathrm{}<0.2`$, and will fade by up to 2.5 mag in B, on average, within 200 Myr after the burst. Follow-up spectroscopy with VLT and HET is underway, models are currently being calculated for their detailed spectral evolution. In any case, this recently discovered mode of ongoing galaxy formation in the local Universe from recycled gas and stars is an interesting field of study just at the beginning of being explored. ## 7 Conclusion and Outlook I presented a versatile model for the evolution of stellar populations and gas that offers a variety of applications from star clusters to nearby and distant galaxies. Only a few of them have been presented here. With a minimum number of parameters our combined chemical, spectrophotometric and cosmological evolution models describe a large number of observables and provide a long evolutionary baseline to compare with high-redshift galaxy abundances and spectra and understand the evolution of various galaxy types from the earliest phases to the present. The chemically consistent treatment is a first attempt to consistently combine 2 out of 3 aspects of galaxy evolution that nature, too, has coupled. Interactions play an important role for the evolution of galaxies over cosmological timescales. While in its present state, the model does not include any dynamical aspect nor spatial resolution, we tried and studied the effects of starbursts accompanying galaxy interactions if gas is involved. Several interesting phenomena were observed in this context, as e.g. the formation of large populations of bright star clusters and of “recycled” Tidal Dwarf Galaxies. Application of our models provided a first step to understand the nature and properties of these systems as well as their possible future evolution. The model has allowed for a series of precise observational predictions, part of which became verified already while others keep standing for a test. Our first attempt to also include the $`3^{\mathrm{rd}}`$ aspect of galaxy evolution, the formation and dynamical evolution of galaxies in their cosmological environments, could not be discussed here (cf. Contardo et al. 1998). Over the years this model was developed, extended and refined, with its input physics continuously updated, observational extragalactic research has seen a tremendous progress. The amount of information from HST and large ground-based surveys is enormous and several quantum steps have been performed, e.g. concerning image resolution with HST, spectral resolution with KECK and WHT, and the number of high-redshift galaxies by the Lyman break technique. With 10 m telescopes, like VLT and HET, observational progress is challenging theory to keep path. The particularly close interplay between observations and the conceptually simple galaxy evolution modelling presented here, has proven very fruitful and stimulating for both sides. Acknowledgements My thanks go to K. Fricke for his encouragement and to all the present and former members of our Galaxy Evolution Group, who – over the years – have contributed to various aspects of the work presented here, i.e. to Harald Krüger, Christian Einsel, Johannes Loxen, Claudia Möller, Oliver Kurth, Ulrich Lindner, Peter Weilbacher, Jens Bicker, and Jochen Schulz. Partial financial support from the Deutsche Forschungsgemeinschaft and the Verbundforschung Astronomie for various aspects of this work is gratefully acknowledged.
warning/0003/cond-mat0003079.html	ar5iv	text	# External fluctuations in front dynamics with inertia: The overdamped limit ## 1 Introduction Front propagation is being the subject of very active research in the last few years: Indeed, the problem of the selection of the front velocity is a paradigm of the dynamical selection mechanisms arising in a large number of physical, chemical and biological systems with a certain kind of instability (see ch and references therein). One of the important questions into which interest has been focused is how the deterministic front scenario is modified by the presence of noise. In this context, the effect of stochastic fluctuations on front dynamics and the modification of its deterministic features have been considered by several authors. A detailed summary of those results, which were mostly devoted to the changes of the front velocity and the spreading of the front position, can be found in book . The present work addresses a related problem which, on the other hand, arises naturally from the above line of research: Is the influence of (external) noise on front propagation the same if inertial effects (as far as we know, neglected in previous work) are taken into account? As is well known, including inertia leads to a description in terms of a damped, hyperbolic partial differential equation. It is important to note that first, a model like this arises when one considers a more realistic jump process for the individuals whose probability density is described by the partial differential equation; and, second, that hyperbolic equations of this kind describe many actual physical phenomena, such as, e.g., population dynamics, nonlinear transmission lines, cell motion, branching random walks, dynamics of ferroelectric domains, and others struc ; do ; merlin ; vm ; gallay ; erzan ; scott . The rôle of inertia in the scenario of (deterministic) front propagation has been studied recently sancho99 (see also vm ; gallay ; erzan , and specially preprint for a detailed study of the underdamped dynamics restricted to the linear regime). In this paper, it was proven that the different dynamical regimes of front propagation using deterministic parabolic models do not change, i.e., the values of the control parameter separating the different regimes do not depend on the inertia parameter (“mass”). However, it was also shown in sancho99 that the values of the front velocities corresponding to every interval of the control parameter and the spatial shape of the propagating fronts do depend on the inertia parameter. As stated above, our explicit objective will now be the study of the rôle of multiplicative noise and its comparison to the results in Joan1 ; Joan2 , in order to understand the interplay of inertia and external stochastic perturbations. Accordingly, we undertake the study of a hyperbolic partial diferential equation with a multiplicative noise term, used to model front dynamics subjected to both inertia and external fluctuations. Opposite to the deterministic case, in which for very small mass or inertia a naive adiabatic elimination procedure (i.e., leaving the second time derivative term out) gives a parabolic equation which describes very accurately the front dynamics sancho99 , we will prove that this is not at all the case when noise is present. As a matter of fact, the starting hyperbolic equation with noise in the Stratonovič interpretation transforms, upon adiabatic elimination, into a parabolic equation with an extra term coming from the noise. Furthermore, the so obtained equation turns out to be equivalent to the usual parabolic equation if interpreted in the Itô sense. In addition, we include numerical simulation results confirming this unexpected result. We report on these results along the following scheme: We begin by presenting our model and by briefly summarizing what is known about the noise influence in the overdamped case. Subsequently, we concern ourselves with the study of the externally perturbed (stochastic) case. We conclude by summarizing our main findings and discussing their implications. ## 2 Model definition and notations Let us begin by introducing the purely deterministic problem. Generally speaking, the situation which we are interested in is generally modeled by the hyperbolic equation $$\varphi _{tt}+\alpha \varphi _t=D\varphi _{xx}+\tau ^1f(\varphi ,a)$$ (1) where $`\alpha `$ is the friction (dissipation), $`D`$ is the diffusion coefficient, $`\tau `$ is the characteristic time of the reaction term and $`a`$ is the external control parameter of the nonlinear reaction term $`f(\varphi ,a)`$. The first step will be the reduction of the number of parameters by introducing the change of variables $`t\tau \alpha t`$, $`x\sqrt{\tau D}x`$; our initial model, Eq. (1), reduces then to $$ϵ\varphi _{tt}+\varphi _t=\varphi _{xx}+f(\varphi ,a),$$ (2) where a new parameter (the “mass”), $`ϵ=(\tau \alpha ^2)^1`$, appears. We note that the information regarding both the characteristic reaction time and the dissipation coefficient is contained now in $`ϵ`$. With this new notation, the parabolic or overdamped limit is obtained by letting $`ϵ0`$ gallay , which leads to (note that it is a singular limit) $$\varphi _t=\varphi _{xx}+f(\varphi ,a).$$ (3) We will refer to this procedure along the paper as naive adiabatic elimination. For the sake of definiteness, we take as a representative example of nonlinear reaction term $$f(\varphi ,a)=\varphi (1\varphi )(a+\varphi ).$$ (4) Such a term can be obtained from a local (bistable) potential, $`f(\varphi ,a)=V^{}(\varphi )`$, with $$V(\varphi )=\frac{a}{2}\varphi ^2\frac{1a}{3}\varphi ^3+\frac{1}{4}\varphi ^4.$$ (5) It is then straightforward to show that the steady states are, $`\varphi _1=0,\varphi _2=1`$ and $`\varphi _3=a`$. We are interested in those solutions which are front-like (kinks) conecting the (unstable if $`a>0`$ and metastable otherwise) state $`\varphi _1=0`$ with the globally stable state $`\varphi _2=1`$. Consequently, we supplement Eqs. (2) and (3) with boundary conditions $`\varphi (\mathrm{},t)=\varphi _2`$, $`\varphi (\mathrm{},t)=\varphi _1`$. Let us now move on to the stochastic version of the problem. As is well known, thermal (additive) noise is not expected to be relevant in actual, experimental contexts, as its ratio to other terms in the governing equations can be estimated to be $`10^9`$ ch . However, in addition to thermal noise we must expect ch multiplicative noise sources arising from fluctuating control parameters hl . Examples of this case are recent experiments on the Belousov-Zabotinsky reaction in a light-sensitive medium showalter ; irene ; wang . The fluctuating light intensity enters as a multiplicative noise in the theoretical modelization of this chemical reaction. According to this, noise is introduced in the system described so far through the parameter $`a`$, which fluctuates according to $$aa+\xi (x,t).$$ (6) The noise $`\xi `$ is gaussian, with zero mean and correlation given by $$\xi (x,t)\xi (x^{},t^{})=2\sigma ^2C(xx^{})\delta (tt^{}),$$ (7) with $`C(xx^{})`$ being the spatial correlation function, normalized by imposing $`C(x)𝑑x=1`$. The noisy parabolic case considered in Joan1 ; Joan2 corresponds to the following stochastic partial differential equation: $$\varphi _t=\varphi _{xx}+f(\varphi ,a)+g(\varphi )\xi (x,t),$$ (8) with $`g(\varphi )=\varphi (1\varphi )`$ in case $`f(\varphi ,a)`$ is given by Eq. (4). Equation (8) with noise statistical properties (7), being well known, will be taken as our reference scenario; its main features are summarized below. Nevertheless, before going into those results, it is important to note that, prior to any other consideration, it is necessary to prescribe a mathematical interpretation of the noise in Eq. (8). Based on physical and mathematical grounds we will follow the Stratonovič interpretation (see stoc1 ; stoc2 for in-depth discussions of the interpretation of stochastic differential equations). This interpretation fulfills two important properties from a physical point of view: First, it corresponds to the white noise limit of a real (nonwhite) noise, and second, when manipulating stochastic terms, usual rules of calculus apply. We note that in an actual experimental situation the noise has necessarily a finite (non zero) characteristic time. However, this time is indeed very small compared with any other characteristic time of the system, and therefore the assumptions of white noise and Stratonovič interpretation seem physically well founded. ## 3 Stochastic results on the overdamped model As already mentioned in the introduction, the overdamped limit with noise, Eq. (8), has been studied recently book ; Joan1 ; Joan2 . It is worth summarizing here the main points in order to compare with our results. In addition, we will refer to the quantities introduced in these section later along the text. In the parabolic model (8), when $`1/2<a<1/2`$ (metastable and nonlinear regimes of the deterministic equation, see vs1 ; vs2 ; vs3 ), starting from a sufficiently localized initial front evolves to the nonlinear solution, and correspondingly the selected velocity of the front is $$v_{nl}(a)=\frac{2a+1}{\sqrt{2(12\sigma ^2C(0))}}.$$ (9) On the other hand, in the linear regime, $`1/2<a<1`$, the velocity is given by $$v_l=2\sqrt{a+\sigma ^2C(0)}.$$ (10) We note that, in the three regimes, this velocity increases monotonously as a function of the noise intensity; we will come back to this result below. ## 4 Stochastic perturbation We now proceed to analyze the case we are interested in, namely the stochastic version of the hyperbolic partial differential equation (2). For convenience, by introducing a new field variable $`\psi `$, we cast it in the form $`\varphi _t`$ $`=`$ $`\psi ,`$ $`ϵ\psi _t`$ $`=`$ $`\psi +\varphi _{xx}+f(\varphi ,a)+g(\varphi )\xi (x,t).`$ (11) At this point, it is important to note that naive adiabatic elimination leads again to Eq. (8). As a starting point, let us recall that in sancho99 we proved that the velocity of the deterministic hyperbolic model can be obtained from the parabolic model by using the transformation $$v_{nl}(ϵ,a)=\frac{1}{\sqrt{ϵ+v_{nl}(a)}}.$$ (12) Therefore, our first aim here is to check whether or not this result applies in the presence of multiplicative noise, i.e., whether or not we can still use the expression (12) substituting for the deterministic velocities the stochastic ones found in Joan1 ; Joan2 . From those papers, we know that the important noise effects come from the fact that the multiplicative noise term has a non zero mean value, provided the Stratonovič interpretation is used. Within that interpretation, it can be found by using only usual stochastic calculus and a little algebra that the mean value of the noisy term in Eq. (11) is book $$g(\varphi )\xi (x,t)=\sigma ^2C(0)\frac{g(\varphi )}{\varphi }\frac{\delta \varphi }{\delta \xi (x^{},t^{})}\|_{x^{}=x,t^{}=t}=0.$$ (13) This result allows us to conjecture that in the hyperbolic case the noise is not as relevant as in the parabolic case. In order to verify this conjecture, we carried out numerical simulations of Eq. (11) (with noise white in space, implying $`C(0)=1/\mathrm{\Delta }x`$ with $`\mathrm{\Delta }x`$ the spatial discretization step) for different choices of the parameters, finding that the velocity is not affected by the perturbation even for very small values of $`ϵ`$. In Fig. 1 we present some of these numerical results. It can be clearly seen the dramatic difference between the effects of noise in the two models. It is most important to stress that the results in Fig. 1 correspond to individual realizations of the noise; of course, we have verified that this is the typical behavior by repeating the simulations very many times. This means that our conjecture that the noise effects and in particular those of the front velocity are not important, is in fact true beyond mean values, i.e., individual fronts do not change their deterministic velocity in the presence of noise. For the parabolic case we see that the velocity does depend on the intensity of the noise, being independent in the hyperbolic case. This is an unexpected result but, as we have just seen, it is in complete agreement with our theoretical analysis, that leads to Eq. (13). Let us now discuss the reasons for this result. It is clear from Eq. (13) that such a null result comes from the fact that the response of the field $`\varphi `$ to the noise is zero at $`t^{}=t`$. This is not the case in the parabolic case (8), for which one finds Joan1 ; Joan2 $`g(\varphi )\xi (x,t)`$ $`=`$ $`\sigma ^2C(0){\displaystyle \frac{g(\varphi )}{\varphi }}{\displaystyle \frac{\delta \varphi }{\delta \xi (x^{},t^{})}}\|_{x^{}=x,t^{}=t}=`$ (14) $`=`$ $`\sigma ^2C(0){\displaystyle \frac{g(\varphi )}{\varphi }}g(\varphi )`$ This non zero mean average is the responsable of the strong effects of noise in the parabolic model (8) which were studied and quantified in Joan1 ; Joan2 . As it turns out, this is not the case in the hyperbolic model (11); therefore, it seems likely that the naive adiabatic elimination procedure, which led to Eq. (8), is not the proper way to take the limit $`ϵ0`$ in Eq. (11), in view of the different behavior of the two models. In order to check this idea, we have followed an alternative, non naive adiabatic elimination procedure sancho82 to see whether an equation different from (8) arises for the overdamped limit. We will outline here the main steps of the calculation following sancho82 (see stoc2 for an alternative presentation). For the sake of simplicity, we rewrite Eq. (11) in a more compact form: $$\varphi _t=\psi ;\psi _t=\frac{1}{ϵ}\left(F(\varphi )+g(\varphi )\xi (x,t)\right),$$ (15) where $`F(\varphi )=f(\varphi ,a)+\varphi _{xx}`$. Formally integrating the second expression in (15) and using the first one, we find an integro-differential equation for the variable $`\varphi `$: $$\psi =\varphi _t=_0^t𝑑t\frac{e^{(tt^{})/ϵ}}{ϵ}\left(F(\varphi )+g(\varphi )\xi (x,t)\right),$$ (16) with the initial condiction $`\psi (0)=0`$. Subsequently, formal integrations by parts are performed in order to obtain a series of terms in powers of $`ϵ`$, which is the situation we are interested in ($`ϵ`$ small). From this formal expansion a Fokker-Planck equation, whose first order term does not depend on $`ϵ`$, is obtained. The calculation is involved, but it does not require any further physical assumptions, it is only (lengthy) algebra. We then skip the details and refer the reader interested in them to sancho82 . The final result is that the corresponding Langevin equation in the Stratonovič interpretation is $$\varphi _t=\varphi _{xx}+f(\varphi ,a)\sigma ^2C(0)\frac{g(\varphi )}{\varphi }g(\varphi )+g(\varphi )\xi (x,t)$$ (17) We have thus obtained a different overdamped limit, for which one can check that, according to Eq. (14), the mean value of the noisy term compensates the new term in Eq. (17), hence rendering the noise contribution null, as in the $`ϵ0`$ case. Therefore, Eq. (17) is the physically consistent overdamped limit of Eq. (11), insofar it exhibits the same behavior as this last one does for any value of the “mass” $`ϵ`$. ## 5 Discussion and conclusions In this work, we have shown analytically and numerically that inertial effects of any magnitude suppress the external noise influence on the velocity of fronts. Whereas the theoretical result has been conjectured by taking averages, our numerical simulations show that the velocity of individual fronts is unchanged by noise. This means that the overdamped (parabolic) equation usually employed to describe front propagation in reaction diffusion model systems is not simply the limit of an underdamped (hyperbolic) version, as in that case it is known Joan1 ; Joan2 that the velocity of the front does depend on the noise strength. In other words, the naive prediction based on deterministic results, Eq. (12), is not correct in the presence of multiplicative noise. We have also shown that, by means of a more involved adiabatic elimination procedure, it is possible to obtain an equation for the overdamped limit which does not show noise effects. However, this equation differs from the usual one by an extra term, arising from elimination, which exactly cancels the noise contribution to the mean velocity of the fronts. From the physical viewpoint, these results are very relevant. Indeed, we have seen that any amount of inertia present in the system will lead to front propagation at the deterministic velocity even in the presence of external noise. Although the calculation has been done for a specific model, the reasons for the vanishing of the noise contribution are generic and do not depend on the specific choice of the reaction term. In this context, it is then clear that, even if a system is considered overdamped, generally speaking there will be some degree of inertia in its dynamics. In that case, the predicted changes in the velocity Joan1 ; Joan2 will not actually occur. Therefore, our results are in fact a criterion to establish whether a system (where front propagation arises) is truly overdamped and correspondingly described by a parabolic equation, or in turn, it is an inertial system with very large damping: As we have seen, the response of the system to external noise would be fundamentally different in both cases. This result is of particular importance in the study of excitable media in noisy environments showalter ; wang . An additional implication of our findings is that, if one is interested in identifying or measuring possible noise effects in hyperbolic problems in the class of Eq. (11), the front velocity is not a good observable and it is necessary to resort to other indicators, such as, e.g., the fluctuations of the velocity. Finally, another remarkable fact is that the naive parabolic limit (8) if interpreted in the Itô sense, it is stochastically equivalent to the correct limit (17) valid for the Stratonovič interpretation. This is but a further indication that the overdamped limit of Eq. (11) is problematic and has to be carefully performed, as in any other instance where, for physical reasons, multiplicative noise has to be considered. On the other hand, the fact that Eq. (8) is the consistent overdamped limit in the Itô interpretation has to do as well with the known result that, in general, perturbative expansions such as the one needed in the adiabatic elimination procedure have many identically zero terms when carried out in Itô sense stoc2 . This result poses questions of a more mathematical character that would be interesting to address in a general framework for stochastic partial differential equations. ## Acknowledgements We thank Esteban Moro for a critical reading of the manuscript. Work at Barcelona was supported by DGES (Spain) through grant PB96-0241. Work at Leganés was supported by DGES (Spain) through grant PB96-0119.
warning/0003/cond-mat0003414.html	ar5iv	text	# Phonon Signature of Charge Inhomogeneity in High Temperature Superconductors YBa2Cu3O6+x ## Abstract Temperature and composition dependences of high-energy longitudinal optical (LO) phonons in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub>, studied by inelastic neutron scattering measurements, provide clear evidence for temperature dependent spatial inhomogeneity of doped charges. The measurements indicate that charge doping increases the volume fraction of the charged regions, while the local charge density remains unchanged. For the optimally doped sample the charge distribution changes sharply near the superconducting transition temperature, with stronger charge inhomogeneity below T<sub>C</sub>. The remarkable magnitude of the phonon response to charge suggests strong involvement of phonons in charge dynamics. PACS No. 74.25.Ke, 63.20.Kr, 74.20.Mn For some time, high-temperature superconductivity in cuprates has been believed to occur in a homogeneous system through a magnetic mechanism. However, these assumptions are seriously challenged by recent experimental observations that suggest spatial charge inhomogeneity and lattice effects. In particular, the observation of the spin/charge stripe structure for non-superconducting La<sub>1.475</sub>Nd<sub>0.4</sub>Sr<sub>0.125</sub>CuO<sub>4</sub> has eloquently explained the observed incommensurate magnetic and lattice diffraction . In superconducting samples, the incommensurate periodicity is observed only in the dynamic spin structure. The dynamic periodicity changes almost linearly with composition and is consistent with the static superlattice periodicity for the non-superconducting composition . Thus, it was postulated that dynamic stripe correlations exist in superconducting cuprates , and various theories have been advanced assuming such correlations . However, direct evidence of a lattice signature of the charge inhomogeneity is incomplete. While our earlier work on the LO phonons in La<sub>1.85</sub>Sr<sub>0.15</sub>CuO<sub>4</sub> (LSCO) implicated charge inhomogeneity , only one composition was studied. Phonon anomalies were recently observed for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub> (YBCO, x = 0.2, 0.35, and 0.6) , but anomalies are rather weak, and temperature dependence was only cursorily characterized. In this report, we describe the results of neutron inelastic scattering measurements of high-energy LO phonons in superconducting YBCO (x = 0.2, 0.35, 0.6, and 0.93) that not only provide strong evidence for temperature dependent charge inhomogeneity, but also clarify important features of the charge inhomogeneity. The results indicate that the local charge density of the charged region remains constant with doping, even in the optimally doped sample, and doping only changes the volume fraction of the charged region. The size of the charged region is microscopic, and for the optimally doped sample the charge separation is significantly enhanced in the superconducting phase. The studied YBCO samples were large single crystal disks with a height of about 1 cm and diameter ranging from 5 cm for the optimally doped sample (x = 0.93) to 4 cm for the strongly underdoped sample (x = 0.2). Measurements were carried out with the HB-3 triple axis spectrometer at the High Flux Isotope Reactor at Oak Ridge National Laboratory. To monochromatize incident neutrons a beryllium (101) reflection was used, while a pyrolitic graphite (002) reflection (Silicon (111) for x = 0.6 sample) was used for an analyzer, set to give a final neutron energy of 14.7 meV. The angular divergence of the beam was 48’-40’-80’-120’ for large x = 0.93 crystal and 48’-60’-80’-240’ for smaller x = 0.2, 0.35 and 0.6 crystals. As in Ref. 7 we focused on the LO mode along the in-plane Cu-O bond direction which is polar at the zone-center and half-breathing at the zone-edge. Inelastic neutron scattering measurements were made at energy transfers from 50 - 80 meV and momentum transfers, Q, along the (100) direction from (3, 0, 0) to (3.5, 0, 0) in the unit of the reciprocal lattice vector (a* = 1.63 Å<sup>-1</sup>). Thus, the measurement only detects longitudinal optical (LO) phonons. Since the samples are twinned, a-axis and b-axis phonons are observed at the same time. The inelastic neutron scattering intensity pattern at T = 10 K for various doping levels is shown in Fig. 1. In this system the undoped (x = 0) sample has a dispersionless LO branch at 75 meV, while doping softens the zone-edge mode down to 55 meV . Thus one might expect continuous softening at the zone-edge as doping level is increased. Instead, Fig. 1 shows that the LO branch is always split into two, the high-energy branch around 75 meV and the low-energy branch around 55 meV. Instead of continuous softening the spectral weight is transferred from the high-energy branch to the low-energy one as doping is increased. This result is most naturally understood in the two-phase picture if we associate the high-energy and low-energy branches with the micro-phases possessing low and high charge densities, respectively. When the charges are segregated into microscopic domains, increasing the doping level does not change the local charge density in the domains, but simply increases their total volume fraction, causing the spectral weight transfer from the high-energy branch to the low-energy one. The characteristic Q-dependence of the two branches indicates that the size of the charged domain is microscopic, since otherwise two branches with a similar Q-dependence in intensity will be observed. For the optimally doped sample the inelastic magnetic neutron scattering is dominated by the resonance at 41 meV, so that the evidence of incommensurate magnetic modulation is weak . The present data, however, indicate very clear and strong charge inhomogeneity for the x = 0.93 sample, suggesting that the magnetic modulation may not be the dominant force to produce charge inhomogeneity. This is in accord with the observation that the charge signature sets in at a higher temperature than the spin signature does for the sample with the 1/8 charge density, implying that charge is driving the stripe ordering . The dispersion shown in Fig. 1 agrees with the earlier result for the optimally doped YBCO , and has a striking resemblance to the one obtained earlier for LSCO . Fig. 2 shows the intensity pattern for the x = 0.93 sample obtained at room temperature, I<sub>300K</sub>, and the difference, I<sub>300K</sub> \- I<sub>10K</sub>. At room temperature the two branches are more connected at the middle, in agreement with the result on LSCO . In order to study the temperature dependence in more detail the scattering intensity as a function of energy transfer was measured at (3.25, 0, 0) at various temperatures. The average intensities over the energy range from 56 to 68 meV (Range 1), I(1), and from 51 to 55 meV (Range 2), I(2), are plotted in Fig. 3(a), and the difference, I(2) - I(1), in Fig. 3(b). It is clear that a sharp change is observed near the superconducting transition temperature of T<sub>C</sub> = 95 K. Above T<sub>C</sub>, I(2) - I(1) is practically constant. This indicates that the spectral weight is shifted from Range 2 to Range 1 as temperature is raised up to T<sub>C</sub>, with no further change above T<sub>C</sub>. Since for the optimally doped YBCO T<sub>C</sub> coincides with the pseudo-gap temperature T, it is not clear whether the change is associated with T<sub>C</sub> or with T. Recent reports on the isotope effect on T lead us to think the latter is more likely. Our preliminary study on x = 0.35 indicates gradual change and saturation around 200 K, which is close to T. An objection was raised recently regarding the temperature dependence in LSCO , this issue involves the problem of focusing of the spectrometer, and will be discussed elsewhere. For YBCO the presence of strong temperature dependence is unquestionable as presented here. The results described here represent the most convincing evidence obtained so far on the spatial inhomogeneity of the charged state in the superconducting cuprates, including the optimally doped sample. While the pattern of inhomogeneity is not directly discernible from the present result, a simulation indicates that a quasi-static stripe pattern is not compatible with the split dispersion, and the cell-doubling model exhibits a better fit . But the shape of the domains is probably stripe-like, since the size of the domain speculated from the shape of the dispersion-less portion of the phonon dispersion is 8 $`\times `$ 20 Å . Thus it is possible to characterize the charged domains as fragments of stripes. The fact that greater charge inhomogeneity was observed below T<sub>C</sub>/T implies that charge inhomogeneity does not compete against superconductivity, but may facilitate its occurrence. As shown in Fig. 4 the half-breathing mode at the zone-edge results in charge transfer between Cu and O, while the polar mode at the zone-center does not. Thus the zone-edge mode is expected to couple strongly to the doped holes in a charge transfer system such as the cuprates . Indeed the two branches are remarkably different in energy (15 - 20 meV), suggesting strong involvement of the low-energy phonon branch in charge dynamics. Such microscopic inhomogeneity was observed for manganites that show colossal magnetoresistance (CMR), even in the metallic phase . Also strong phonon softening similar to the one shown here was observed for manganites , suggesting intimate involvement of phonons in producing charge segregation. At low doping levels charges in manganites are localized as spin/lattice polarons or polaron aggregates, but as the doping level is increased the polaron aggregates percolate to produce an insulator-to-metal transition . In the metallic phase, charges in the percolating charged domains are not localized, but still spatially restricted. It is possible that a similar picture applies to cuprates, probably with a smaller length scale. While the details of the charged domains and the pairing mechanism are still unclear, the present results strongly challenge conventional wisdom concerning the origin of high-temperature superconductivity. The authors are grateful to A. R. Bishop, V. J. Emery, L. P. Gor’kov, A. Bussmann-Holder, and V. Kresin for useful discussions. The research at University of Pennsylvania was supported by the National Science Foundation through DMR96-28136. Measurements were made at Oak Ridge National Laboratory, which is managed by Lockheed Martin Energy Research under contract no. DE-AC05-96OR22464 for the Department of Energy.
warning/0003/math-ph0003012.html	ar5iv	text	# 1 Introduction ## 1 Introduction In the past decade in a series of papers Verbeure and coworkers developed a beautiful and ingeneous framework to study so-called macroscopic fluctuation phenomena in systems and various regimes of quantum statistical mechanics (see the cited literature). The approach is to a large extent based on a quantum variant of the central limit theorem and is mainly performed in real (i.e. configuration) space. Among other things, the general goal is it, to study the limit behavior of correlation functions of so-called fluctuation observables, i.e. appropriately renormalized averages of microscopic observables, averaged over volumes, $`V`$, which approach the whole space, $`^n`$, say. Typically, one arrives, depending on the type of clustering of the microscopic $`l`$-point functions, at certain simple limit algebras as e.g. $`CCR`$. We approach the field from a slightly different angle. In a first step we choose another averaging procedure, which avoids sharp volume cut-offs and, a fortiori, has a very nice and transparent scaling behavior. This is then exploited in the following analysis which systematically develops so-called Fourier-space and energy-momentum spectral methods of observables and correlation functions. We consider it to be an advantage that the calculations turn out to be relatively transparent and lead in a direct way to the desired results. We first treat the case of normal fluctuations and $`L^1`$-clustering. We show that all the truncated $`l`$-point functions vanish for $`l3`$ while they approach a finite, non-trivial limit for $`l=2`$. The analysis is done both for the $`(k=0)`$\- and the $`(k0)`$-modes. We emphasize that the calculations for net-momentum different from zero remain also very simple. A variant of the method is then applied to the case of $`L^2`$-clustering. In the second part of the paper we embark on the analysis of fluctuations in the presence of spontaneous symmetry breaking (ssb). In a first step we prove some general results in the context of $`ssb`$ and the Goldstone phenomenon. We then address the problem of macroscopic fluctuations within this context. Among other things, we give a general and rigorous proof that the limit fluctuations are always classical for temperature states (a phenomenon already observed by Verbeure et al in various simple models). The paper ends with a treatment of extremely poor clustering, which can be controlled by a new method we develop in the last section. To sum up, we think that in our view the two different frameworks seem to neatly complement each other and should lead to further interesting results if being combined. ## 2 The Scenario of Normal Fluctuations The following analysis works for statistical equilibrium states and/or for vacuum states in quantum field theory. To avoid constant mentioning of the respective scenario we are actually working in, we usually treat equilibrium (i.e. KMS-) states, to fix the framework. Now, let $`\mathrm{\Omega }`$ be the vacuum or equilibrium state (rather its GNS-representation; usually we work within a concrete Hilbert space, $``$). As an abstract state we denote it by $`\omega `$. Expectations of observables are written as $$A=\omega (A)=(\mathrm{\Omega },A\mathrm{\Omega })$$ (1) with $`A`$ taken from the local algebra, $`𝒜_0𝒜`$, the latter one being the quasi-local norm closure of $`𝒜_0`$. We assume $`\mathrm{\Omega }`$ to be cyclic with respect to $`𝒜_0`$ or $`𝒜`$. That is, we assume $$\overline{𝒜_0\mathrm{\Omega }}=$$ (2) There are certain differences as to the (assumed) locality properties of the dynamics between (non-)relativistic statistical mechanics and relativistic quantum field theory (RQFT). Denoting the time evolution (acting on the algebra of observables) by $`\alpha _t`$, we are confronted with the following phenomenon. ###### Observation 2.1 In RQFT part of the usual framework is the assumption $$\alpha _t:𝒜_0𝒜_0$$ (3) while in statistical mechanics (due to weaker locality behavior) we have in the generic case only $$\alpha _t:𝒜𝒜$$ (4) whith $`𝒜_0`$ usually not left invariant as the observables will typically develop infinitely extended tails. Furthermore, we assume once for all that our system is in a pure, translation invariant phase, that is $`\mathrm{\Omega }`$ is extremal translation invariant under the space translations ( which can, as in the case of lattice systems, also be a discrete subgroup). There can of course exist several coexisting pure phases at the same external parameters as in the regime below a phase transition threshold. These assumptions imply that we can expect certain cluster properties, i.e. decay of correlations (see e.g. ). ### 2.1 Definition of Ordinary Fluctuation Operators We begin by defining the fluctuation operators in the normal situation as it was done in . We assume, for the time being, $`L^1`$-clustering for the two-point-function, that is $$\|A(x)B^T\|d^nx<\mathrm{}A,B𝒜$$ (5) with $`A(x)`$ the translate of $`A`$ and $$AB^T=ABAB$$ (6) Once for all we assume, to simplify notation, in our particular context that the occurring observables are normalized to $`A=0`$ unless otherwise stated. ###### Definition 2.2 We define the normal (finite volume) fluctuation operators as $$A_V^F:=1/V^{1/2}_VA(x)d^nx=:1/V^{1/2}A_V$$ (7) In a next step one wants to give sense to these objects in the limit $`V\mathrm{}`$. From the $`L^1`$-condition we however infer $$\|(A_V^F\mathrm{\Omega },B\mathrm{\Omega })\|1/V^{1/2}_^n\|(A(x)\mathrm{\Omega },B\mathrm{\Omega })\|d^nx0$$ (8) Hence $`A_V^F\mathrm{\Omega }0`$ on a dense set. Furthermore we have $$\begin{array}{c}(A_V^F\mathrm{\Omega },A_V^F\mathrm{\Omega })=1/V_V_V(A(x)\mathrm{\Omega },A(y)\mathrm{\Omega })𝑑x𝑑y=\hfill \\ \hfill 1/V_V𝑑x\left(_{Vx}A^{}A(yx)d(yx)\right)\end{array}$$ (9) This is less or equal to $$(1/V)V\underset{x}{sup}(_{Vx}\|(\mathrm{})\|)_^n\|F(yx)\|d^n(yx)<\mathrm{}$$ (10) (for convenience we sometimes denote a general two-point function by $`F(xy)`$). This suffices to prove weak convergence to zero for $`A_V^F\mathrm{\Omega }`$ on the total Hilbert space $``$. ###### Remark 2.3 We note that this proves also the well-known normal-fluctuation result $`A_VA_VV`$ in the $`L^1`$-case. Under certain well-specified conditions the fluctuations can even be weaker than normal. If e.g. $`Q_V`$ is the local integral over a conserved quantity we proved a divergence significantly weaker than $`V`$ (cf. ). But in general the local fluctuations will diverge in the limit $`V\mathrm{}`$ in contrast perhaps to ordinary intuition, even if the quantity is globally conserved due to quantum fluctuations (see also the section about spontaneous symmetry breaking) A weaker than normal divergence can occur in the following situation. An asymptotic behavior $`V`$ does only prevail if $`_VF(u)𝑑u0`$ in the limit $`V\mathrm{}`$. On the other side such correlation functions tend to oscillate about zero (for physical reasons; there are e.g. usually preferred relative positions in, say, a quantum liquid). In other words, while $$F(u)𝑑u=0$$ (11) may seem to be rather ungeneric at first glance, it can nevertheless happen in a specific context. The general situation is analyzed in the above reference; certain examples of better than normal fluctuations were also found by Verbeure et al in e.g. (see also ). For the fluctuation operators themselves we have due to locality for $`A,B𝒜_0`$: $$[A_V,B]\text{independent of }V\text{ for }VV_0V_B$$ (12) for some $`V_0`$ which contains the localisation region $`V_B`$ for $`B𝒜_0`$. We then have $$\underset{V}{lim}(A_V^FC\mathrm{\Omega },B\mathrm{\Omega })=\underset{V}{lim}([A_V^F,C]\mathrm{\Omega },B\mathrm{\Omega })+\underset{V}{lim}(A_V^F\mathrm{\Omega },C^{}B\mathrm{\Omega })$$ (13) We have already shown that the second term goes to zero. In the first term the commutator becomes $$[A_V^F,C]=V^{1/2}[A_{V_0},C]$$ (14) and hence the first term goes also to zero. In case we assume only $`A𝒜`$ a further $`L^1`$-condition for the three-point function is needed to arrive at the same result. As $`𝒜_0\mathrm{\Omega }`$ is assumed to be dense in $``$ and $`A_V^F<\mathrm{}`$ uniformly in $`V`$, we have ###### Proposition 2.4 $`L^1`$-clustering implies that $$A_V^F0\text{weakly on },A_V^F\mathrm{\Omega }<\mathrm{}\text{uniformly in }V$$ (15) but $`A_V^F\mathrm{\Omega }`$ bounded away from zero in general. That is, $`A_V^F`$ does not converge strongly to zero and, a fortiori, there is no convergence in norm. This clearly shows that, in order to have non-trivial limit operators, one has to leave the original Hilbert-space of microscopic observables and has to define or construct an entirely new representation living on a different state. ### 2.2 A Smoothed Version of Fluctuation Operators Since we employ in the following so-called Fourier-methods and related calculational tools, it is advantageous to change to a smoother version of fluctuation operators. As everybody knows, sharp volume cut-offs are both a little bit artificial and technically nasty, since they may sometimes lead to non-generic or spurious effects. In other branches of rigorous statistical mechanics or axiomatic quantum field theory volume integrations have therefore frequently been emulated or implemented in a slightly different way (see e.g. ). Two choices have basically been in use with the second version having much nicer properties in several respects as we will explain below. Instead of integrating over a sharp volume, $`V`$, centered e.g. around the coordinate origin, one integrates the shifted observable, $`A(x)`$, over a smooth test function localized basically in $`V`$ but having smooth tails. Remark: As $`V`$ we choose in the following a ball centered at the origin with radius $`R`$ and let $`R`$ go to infinity. ###### Definition 2.5 Two admisssible families of test functions are the following ones: $`f_R(x)0`$ smooth with $$f_R(x):=\{\begin{array}{cc}1\hfill & \text{for }\|x\|R\hfill \\ 0\hfill & \text{for }\|x\|R+h\hfill \end{array}$$ (16) or $$f_R(x):=f(\|x\|/R)\text{with}f(s)=\{\begin{array}{cc}1\hfill & \text{for }\|x\|1\hfill \\ 0\hfill & \text{for }\|x\|2\hfill \end{array}$$ (17) Note that the latter choice has much nicer behavior under Fourier transform while working with the Fourier transforms of the former version or e.g. the indicator function of the volume $`V`$ is quite cumbersome). On the other hand, the latter version has tails which are also scaled. ###### Lemma 2.6 $$\widehat{f}_R(k)=constR^n\widehat{f}(Rk)$$ (18) where here and in the following “ $`const`$” denotes an (in this context) irrelevant numerical factor which, a fortiori, may change in the course of a calculation. With the help of this smearing functions we now define ###### Definition 2.7 (Smooth Volume Integration) We redefine the fluctuation operators in the following way $$A_R^F:=R^{n/2}A(x)f_R(x)d^nx$$ (19) with $`f_R`$, unless otherwise stated, the family given in the second example above (remember $`A:=0`$). ## 3 The Limiting Case for Normal Fluctuations In order to arrive at a rigorous definition of fluctuation operators in a certain limit state we will follow a line of arguments which may complement the treatment of Verbeure et al in several respects. We will study directly the macroscopic limit of the n-point functions with the help of certain momentum space methods. As they are perhaps not so common in statistical physics we will give the technical details below. ### 3.1 Some Generalities Any n-point (correlation) function of the kind $`A_1(x_1)\mathrm{}A_n(x_n)`$ with the $`A_i(x_i)`$ the translates of the observables $`A_i`$ (which may also contain an implicit time variable $`t_i`$ which is however kept fixed in the following) is written as $`W(x_1,\mathrm{},x_n)`$. With the state $`\mathrm{\Omega }`$ being translation invariant we have $$W(x_1,\mathrm{},x_n)=W(x_1x_2,\mathrm{},x_{n1}x_n)$$ (20) To express cluster properties in a clear way, we introduce the so-called truncated correlation functions via the following recursion relation: $$W(x_1,\mathrm{},x_n)=\underset{part}{}\underset{P_i}{}W^T(x_{i_1},\mathrm{},x_{i_k})$$ (21) where the sum extends over all partitions of the set $`\{1,\mathrm{},n\}`$ into subsets $`P_i`$ with the elements in each subset ordered as $`i_1<i_2\mathrm{}<i_k`$. The first elements of the recursion are $$W(x)=W^T(x)=0\text{in our case}$$ (22) $$W^T(x_1,x_2)=W(x_1,x_2)W(x_1)W(x_2)$$ (23) ###### Observation 3.1 In the truncated correlation functions the vacuum state, ground state or equilibrium state, $`\mathrm{\Omega }`$, has been eliminated in a symmetric way, so that we have, in a sense to be specified, $$W^T(x_1,\mathrm{},x_n)0\text{for}sup\|x_ix_j\|\mathrm{}$$ (24) In this section we assume the following cluster property $$W^T(x_1,\mathrm{},x_n)L^1\text{in the variables}\{x_1x_2,\mathrm{},x_{n1}x_n\}$$ (25) From the above we see that the original hierarchy of $`n`$-point functions can be reconstructed from the new hierarchy of truncated $`n`$-point functions, which have more transparent cluster properties. The $`L^1`$-condition allows us to Fourier transform the $`W^T(x_1,\mathrm{},x_l)`$ and we get from translation invariance: $$\begin{array}{c}const\stackrel{~}{W}^T(p_1,\mathrm{},p_l)e^{i{\scriptscriptstyle p_ix_i}}dp_i=\hfill \\ \hfill W^T(x_1,\mathrm{},x_l)=W^T(x_1x_2,\mathrm{},x_{l1}x_l)\\ \hfill =const\widehat{W}^T(p_1,p_1+p_2,\mathrm{},p_1+\mathrm{}p_{l1})\delta (p_1+\mathrm{}p_l)e^{i{\scriptscriptstyle p_ix_i}}dp_i\\ \hfill =const\widehat{W}^T(q_1,\mathrm{},q_{l1})e^{i_{i=1}^{l1}q_iy_i}\underset{i=1}{\overset{l1}{}}dq_i\end{array}$$ (26) with $$y_i:=x_ix_{i+1},q_i=\underset{j=1}{\overset{i}{}}p_ji(l1)$$ (27) The functional determinant $`det(q/p)`$ is one and we can regard $`\widehat{W}^T`$ either as a function of the $`q_i`$’s or the $`p_i`$’s. We hence have ###### Lemma 3.2 As a Fourier transform of a $`L^1`$-function $`\widehat{W}^T(p_1,\mathrm{},p_{l1})=\widehat{W}^T(q_1,\mathrm{},q_{l1})`$ is a continuous and bounded function which decreases at infinity in the $`q`$-variables. ### 3.2 The $`(k=0)`$-Modes We now study the limit of truncated $`l`$-point functions with the entries being fluctuation operators $`A_R^F`$, more precisely their Fourier transforms, i.e. $$\begin{array}{c}A_R^F(1)\mathrm{}A_R^F(l)^T=\hfill \\ \hfill constR^{ln/2}\widehat{f}(Rp_1)\mathrm{}\widehat{f}(R[p_1+\mathrm{}+p_{l1}])\widehat{W}^T(p_1,\mathrm{},p_{l1})dp_i\\ \hfill =constR^{ln/2}R^{(l1)n}\widehat{f}(p_1^{})\mathrm{}\widehat{f}([p_1^{}+\mathrm{}+p_{l1}^{}])\widehat{W}^T(p_1^{}/R,\mathrm{},p_{l1}^{}/R)dp_i^{}\end{array}$$ (28) $`\widehat{W}`$ is continuous and bounded and the $`\widehat{f}`$’s are of rapid decrease. Hence we can perform the limit $`R\mathrm{}`$ under the integral and get ###### Theorem 3.3 The expression $`A_R^F(1)\mathrm{}A_R^F(l)^T`$ scales as $`R^{(2l)n/2}`$. This implies that for $`l>2`$ the above limit is zero, for $`l=2`$ the limit is a finite number bounded away from zero in general. In other words we have $$\underset{R\mathrm{}}{lim}A_R^F(1)\mathrm{}A_R^F(l)^T=0\text{for}l>2$$ (29) and $$\underset{R\mathrm{}}{lim}A_R^F(1)\mathrm{}A_R^F(l)=\underset{R\mathrm{}}{lim}\underset{part}{}\underset{\{ij\}}{}A_R^F(i)A_R^F(j)$$ (30) The relation between the original microscopic system $`(𝒜,\omega )`$ and the coarse-grained system of fluctuation operators is a little bit subtle. Note that $`\omega _F`$, the limit state to be constructed, can no longer be considered as a state or something like that on the original algebra nor can the fluctuation operators be considered as a representation of, say, $`𝒜`$. One aspect of the impending problems can perhaps best be seen by realizing that e.g. $$(AB)_V^FA_V^FB_V^F$$ (31) which pertains also in the limit. That is, in a sense to be defined, we have $$(AB)^FA^FB^F$$ (32) the same holding in general for all the higher products. This is one source of non-uniqueness as there is no invariant discrimination between an observable regarded as a single object to be scaled and as a product of other observables, where now each factor has to be scaled separately. The appropriate point of view has to be a different one (as has also been emphasized by Verbeure et al, cf e.g. , second ref. p.540f and private communication). The picture remains relatively clear for the intermediate scales, $`V<\mathrm{}`$. We have a start system $`(𝒜,\omega )`$, labelled by, say, $`V=0`$. On every scale $`V`$ we have a new algebra, $`𝒜_V^F`$, (actually a subalgebra of $`𝒜`$), generated by the observables $`A_V^F,A𝒜`$ (including arbitrary finite products $`(A_1\mathrm{}A_n)_V^F`$). If we prefer to consider this algebra on scale $`V`$ as a new abstract algebra (i.e. forgetting about the underlying finer algebra $`𝒜`$), we get also a new, coarse-grained state via the identification $$\omega _V^F(\mathrm{\Pi }A_V^{F,i}):=\omega (\mathrm{\Pi }A_V^{F,i})$$ (33) (A related philosophy was expounded by Buchholz and Verch in e.g. within the context of the algebraic analysis of ultra-violet behavior in quantum field theory.) The map $$R_V:𝒜𝒜_V^F$$ (34) can be viewed as kind of a renormalization map, which does however not preserve the algebraic structure (i.e.the algebras are in general not isomorphic). Furthermore one gets a “new” dynamics on this algebra by defining $$\alpha _t^V(A_V^F):=(\alpha _t(A))_V^F$$ (35) ###### Remark 3.4 In our context $`\alpha _t`$ is assumed to commute with the space translations or with a corresponding lattice version, that is, we have $`\alpha _t(A_V^F)=(\alpha _tA)_V^F`$. (Furthermore it may turn out to be reasonable to scale the time variable on the lhs also.) On the other hand, in order to construct the limit theory itself, one can proceed in a slightly different direction. The above limits of n-point functions define a consistent hierarchy of new n-point functions which then allow to define a new limit system via the so-called reconstruction theorem (for a pendant in quantum field theory see e.g. ). Put differently, we define limit objects, $`\{A_i^F\}`$, the so-called fluctuation operators, which live in a new Hilbert space built upon the new state, $`\omega _F`$, defined by the limits: $$\omega _F(A_1^F\mathrm{}A_n^F):=\underset{R\mathrm{}}{lim}A_{1,R}^F\mathrm{}A_{n,R}^F=\underset{part}{}\underset{\{ij\}}{}\omega _F(A_i^FA_j^F)$$ (36) Note however that the so-called Gelfand-ideal, $`I_F`$, is large, that is, there are a lot of elements of $`𝒜`$ which are mapped to zero by this limit with $$I_F:=\{A;\omega _F((A^F)^{}A^F)=0\}$$ (37) This is of course typical for such kind of mean-values, as e.g. all space-translates of $`A`$ yield the same limit element. Shifting one of the observables in the above $`l`$-point functions by, say, $`a_i`$ yields an extra factor $`e^{ip_ia_i}`$ in the Fourier transform which after the above coordinate tranformation goes over into $`e^{ip_i^{}/Ra_i}`$ which goes to one. Summing up we have ###### Conclusion 3.5 With the help of equation (36) we construct a new limit system, consisting of the algebra of fluctuation operators, $`𝒜_F`$, and the limit state $`\omega _F`$. The well-known GNS-construction (see e.g. ) allows to construct the corresponding Hilbert-space representation with $$\omega _F(A_1^F\mathrm{}A_n^F)=(\mathrm{\Omega }_F,A_1^F\mathrm{}A_n^F\mathrm{\Omega }_F)$$ (38) (where, by abuse of notation, we do not discriminate between operators and their equivalence classes on the rhs). As all the $`n`$-point functions decay into a product of $`2`$-point functions all the commutators are $`c`$-numbers: $$[A^F,B^F]=\omega _F([A^F,B^F])$$ (39) The system of fluctuation operators is a quasi-free system (cf. ) Taking now self-adjoint elements one can, as in , represent the new system as a representation of the $`CCR`$ over the real vector space of s.a. operators. Our scalar product, induced by the hierarchy of $`n`$-point functions, can be split in the following way. $$(A^F\mathrm{\Omega }_F,B^F\mathrm{\Omega }_F)=Re(\mathrm{})+iIm(\mathrm{})=:s_F(A^F,B^F)+(i/2)\sigma _F(A^F,B^F)$$ (40) $$\omega _F([A^F,B^F])=\sigma _F(A^F,B^F)$$ (41) where $`\sigma _F`$ defines a symplectic form. The Weyl-operators, $`e^{iA^F}`$ with $`A^F`$ s.a., fulfill the $`CCR`$-relations $$\omega _F(e^{iA^F})=e^{1/2s_F(A^F,A^F)}$$ (42) $$e^{iA^F}e^{iB^F}=e^{i(A^F+B^F)}e^{i/2\sigma _F(A^F,B^F)}$$ (43) In our context the first equation can e.g. be verified as follows: Only the $`2n`$-point functions are different from zero. On the lhs we hence have $$\omega _F(e^{iA^F})=(1)^n/(2n)!\omega _F([A^F]^{2n})$$ (44) It remains to count the number of partitions of an $`2n`$-set into $`2`$-sets. This number is $`(2n)!/2^nn!`$. In (44) we now get for $`A^F`$ s.a. on the rhs $$\underset{n}{}1/n!(1/2\omega _F(A^FA^F))^n=e^{1/2s_F(A^FA^F)}\mathrm{}$$ (45) The above general cluster result of the limit $`n`$-point functions make the study of the limit time evolution relatively straightforward. In a first step it suffices to study the $`2`$-point functions. We define the time evolution in the limit theory by $$\omega _F(A^F(t^{})B^F(t)):=lim\omega (A_V^F(t^{})B_V^F(t))=lim\omega (A(t^{})_V^FB(t)_V^F)$$ (46) On the limiting GNS-Hilbert space constructed above we now get a bounded sesquilinear form $`(x,y(t))`$ which, by standard results, yields a bounded operator $`U^F(t)`$ implementing the time evolution. Here we use that the limit n-point functions are products of $`2`$-point functions. Furthermore we infer with the help of the above limit process that $$(U_t^Fx,U_t^Fy)=\omega _F(\mathrm{})=lim\omega (\mathrm{})=(x,y)$$ (47) In other words, we arrive at the following conclusion ###### Theorem 3.6 The preceding construction yields a strongly continuous unitary time evolution on the limiting GNS-Hilbert space. Another point worth to be mentioned (since it might perhaps be overlooked) is the question of the non-triviality of the commutators $$[A^F,B^F]=\omega _F([A^F,B^F])$$ (48) In principle it could happen that all the expectation values on the rhs vanish. In that case the limit algebra would be abelian and the fluctuations classical. In a more general context (cf. e.g. ) this problem is more complicated. In our situation this question can however be answered in a rather straightforward way. We have $$\underset{V}{lim}\omega ([A_V^F,B_V^F])=\underset{V}{lim}\omega ([A_V,V^1B_V])$$ (49) For $`A,B𝒜_0`$, i.e. local, the rhs equals $$\underset{V}{lim}\omega ([A_V,B])$$ (50) We know candidates which lead to a vanishing of the limit for all $`B𝒜_0`$. For $`A`$ chosen s.a. these are the generators of conserved symmetries, written $$Q:=A(x)d^nx$$ (51) Usually they are assumed to commute with the time evolution, expressed as $`Q(t)=Q`$, hence the above limit would also be zero on the full quasi-local algebra. This situation, more specifically the case of spontaneous symmetry breaking (ssb) and Goldstone phenomenon, will be dealt with in more detail in section 5. In any case, as conserved symmetries are usually not so numerous, we may presume that, in the generic case, not all of these commutators will be zero. For $`A,B`$ not necessarily strictly local our above more general formalism is useful. With $$\omega (A(x)B)=F_{AB}(x),\omega (BA(x))=G_{AB}(x)$$ (52) the vanishing of the commutator would imply: $$\begin{array}{c}0=[A^F,B^F]=\underset{R}{lim}R^n\|\widehat{f}(Rp)\|^2(\widehat{F}_{AB}(p)\widehat{G}_{AB}(p))d^np\hfill \\ \hfill =\underset{R}{lim}\|\widehat{f}(p)\|^2(\widehat{F}_{AB}(p/R)\widehat{G}_{AB}(p/R))d^np\\ \hfill =(\widehat{F}_{AB}(0)\widehat{G}_{AB}(0))\|\widehat{f}(p)\|^2d^np\end{array}$$ (53) by the theorem of dominated convergence (note that we are in the $`L^1`$-situation). Hence we have the result ###### Proposition 3.7 $$[A^F,B^F]=0\widehat{F}_{AB}(0)=\widehat{G}_{AB}(0)$$ (54) that is $$F_{AB}(x)d^nx=G_{AB}(x)d^nx$$ (55) or $$(\mathrm{\Omega },[A(x),B]\mathrm{\Omega })d^nx=0$$ (56) which is the same result as in the strictly local case. ### 3.3 The $`(k0)`$-Modes Up to now only the $`(k=0)`$-modes of fluctuation operators, i.e. $`lim_VV^{n/2}_VA(x)d^nx`$, have been studied. For various reasons it is useful to have corresponding formulas at hand for fluctuation observables containing a certain net-momentum. This problem was studied by Verbeure et al in e.g. and the results were applied in e.g. in the analysis of Goldstone modes. In the original (real-space) approach the necessary calculations turned out to be quite involved and far from being simple. This is another case in point to demonstrate the merits of our Fourier space scaling methods. Instead of the original scaling operators, $`A_V^F`$ or $`A_R^F`$, we now study their $`k0`$-variants, $`A_R^F(k)`$. We begin with a technical lemma. ###### Lemma 3.8 $$\widehat{A}(k):=(2\pi )^{n/2}e^{ikx}A(x)d^nx$$ (57) is an operator-valued distribution (We use the convention $`\widehat{f}(k)=(2\pi )^{n/2}e^{ikx}f(x)d^nx`$) Remark: For a systematic use and proofs of such energy-momentum techniques in quantum statistical mechanics we refer to e.g. where also some more mathematical background is provided. Integrating now over $`e^{iqx}f_R(x)`$, we get the $`q`$-mode fluctuation operators. $$\begin{array}{c}A_R^F(q):=R^{n/2}A(x)e^{iqx}f_R(x)d^nx=R^{n/2}\widehat{A}(k+q)\widehat{f}(Rk)d^nk\hfill \\ \hfill =R^{n/2}\widehat{A}(k)\widehat{f}(R(kq))d^nk\end{array}$$ (58) We can now proceed in exactly the same way as above in the case of the zero-mode analysis and calculate the truncated $`l`$-point functions $`A_R^F(1,q_1)\mathrm{}A_R^F(l,q_l)^T`$ (where the indices $`1`$ to $`l`$ label different observables). The only thing that changes are the test functions, i.e. $`f_R(x)e^{iq_kx}f_R(x)`$. We arrive at the conclusion: ###### Theorem 3.9 ($`q`$-Mode Fluctuation Operators) In the case of $`L^1`$-clustering all truncated correlation functions vanish for $`l3`$ and the $`l`$-point functions are again sums of products of $`2`$-point functions. The concrete form of the limit-$`2`$-point functions is given in formula (62). If we calculate the limt-$`2`$-point functions explicitly we get: $$\begin{array}{c}A_R^F(q_1)B_R^F(q_2)^T=\hfill \\ \hfill R^n\widehat{A}(k_1+q_1)\widehat{B}(k_2+q_2)^T\delta (k_1+q_1+k_2+q_2)\widehat{f}(Rk_1)\widehat{f}(Rk_2)𝑑k_1𝑑k_2\\ \hfill =R^n\widehat{A}(k_1+q_1)\widehat{B}((k_1+q_1))^T\widehat{f}(Rk_1)\widehat{f}(R(k_1+q_1+q_2))𝑑k_1\\ \hfill =R^n\widehat{A}(k)\widehat{B}(k)^T\widehat{f}(R(kq_1))\widehat{f}(R(k+q_2))𝑑k\end{array}$$ (59) With $`k^{}:=R(kq_1)`$ we arrive at $$\widehat{W}^T(k^{}/R+q_1)\widehat{f}(k^{})\widehat{f}(k^{}R(q_1+q_2))𝑑k^{}$$ (60) By assumption $`\widehat{W}^T`$ is in $`L^1`$, $`\widehat{f}`$ is of rapid decrease, so the limit can again be carried out under the integral and we have ###### Observation 3.10 For $`q_1+q_20`$ it holds $$\underset{R}{lim}A_R^F(q_1)B_R^F(q_2)^T=0$$ (61) For $`q=q_1=q_2`$ we get on the other side $$\underset{R}{lim}A_R^F(q)B_R^F(q)^T=\widehat{W}^T(q)\widehat{f}(k)\widehat{f}(k)𝑑k$$ (62) In other words, the limit tests the spectral momentum of the two-point function. ## 4 The Case of $`L^2`$-Clustering Before we embark on an investigation of the situation in the regime where phase transitions, vacuum degeneracy and/or spontaneous symmetry breaking (ssb) prevail, we briefly address the case where the clustering is weaker than $`L^1`$ but still $`L^2`$, say. Our above Fourier-space approach can easily handle also this more singular situation. We hence assume now that the truncated $`l`$-point functions cluster only in the $`L^2`$-sense in the difference variables. Now we cannot conclude that the Fourier transform. is bounded and continuous, but we know it is again an $`L^2`$-function. We repeat the first steps of the above calculation with, however, another scaling exponent, $`\alpha `$, which we leave open for the moment. ###### Definition 4.1 In the general case we define fluctuation operators by $$A_R^F:=R^\alpha A(x)f_R(x)d^nx$$ (63) We get $$\begin{array}{c}A_R^F(1)\mathrm{}A_R^F(l)^T=\hfill \\ \hfill constR^{l(n\alpha )}\widehat{f}(Rp_1)\mathrm{}\widehat{f}(Rq_{l1})\widehat{W}^T(q_1,\mathrm{},q_{l1})dq_i\end{array}$$ (64) where the $`\{p_i\}`$ are linear functions of the $`\{q_i\}`$ as described above. We now apply the Cauchy-Schwartz inequality $$\begin{array}{c}\|lhs\|constR^{l(n\alpha )}\left[(\widehat{f}(Rp_1)\mathrm{}\widehat{f}(Rq_{l1}))^2dq_i\right]^{1/2}\hfill \\ \hfill \left[(\widehat{W}^T(q_1,\mathrm{},q_{l1}))^2dq_i\right]^{1/2}\end{array}$$ (65) In the first integral on the rhs we make again a variable transformation from $`q_i`$ to $`q_i^{}:=Rq_i`$, yielding an overall scaling factor $$R^{l(n\alpha )}R^{(l1)n/2}$$ (66) We again want the limits of the $`2`$-point functions to be both finite and non-trivial, i.e. different from zero in general. ###### Proposition 4.2 To make the rhs of (65) finite in the limit for $`l=2`$ the maximal $`\alpha `$ to choose is $$3n4\alpha =0\text{i.e.}\alpha =(3/4)n$$ (67) For a general $`l`$ this leads to the scaling exponent $`(n(1/2)ln)/2`$, which is negative for $`l3`$. Hence, all higher truncated $`l`$-point functions vanish in the limit. However, to guarantee that the result is really non-trivial, we have to analyze the situation in more detail as the above estimate is only an inequality. In the case of $`L^1`$-clustering $`\alpha =n/2`$ was appropriate. The largest value which can occur in the $`L^2`$-case is the above maximal $`\alpha =(3/4)n`$. If we want to avoid that the $`2`$-point functions vanish in the limit we have to choose in the $`L^2`$-case $$(1/2)n<\alpha (3/4)n$$ (68) depending on the concrete decay of the $`2`$-point functions in configuration space. We see that, evidently, the situation is now less canonical as compared to the $`L^1`$-case. Remark: A related situation (on a lattice) was analyzed by Verbeure et al in , where a clustering weaker than $`L^1`$ was considered with, however, the additional input that the local algebras, sitting at the points of the lattice, form a finite-dimensional Lie-algebra. In that case, suitable scaling exponents are chosen to render the auto-correlation functions finite and non-vanishing, while, on the other side, the finiteness of the limit $`3`$-point functions has to be imposed as an extra assumption. Under this proviso one gets the existence of a limit Lie-algebra, but nevertheless results are only partial while perhaps, on the other side, being also more interesting. We do not want to dwell too much on this point at the moment, as progress seems to be to a certain extent model-dependent. Furthermore, we develop a different approach in the last section which is able to cope with any kind of poor cluster behavior. If we want to guarantee the apriori existence or vanishing of the truncated $`3`$-point functions with the help of our above $`L^2`$-estimate (65), we have to restrict the chosen $`\alpha `$ in the following way. ###### Corollary 4.3 If the appropriate $`\alpha `$ fulfills $`\alpha >(2/3)n`$, we get a negative scaling exponent for $`l3`$ as $$n(1/3)ln0\text{for}l3$$ (69) For $`\alpha =2/3`$ the $`3`$-point functions are finite. ###### Remark 4.4 One would get corresponding relations for smaller $`\alpha `$ but higher correlation functions, beginning from a certain order, $`l_0(\alpha )`$ say. On the other hand, one cannot guarantee the apriori existence of the $`l`$-point functions for $`2<l<l_0(\alpha )`$ as the general scaling relation reads for $`ll_0(\alpha )`$: $$l(2\alpha n)>n\text{and}\alpha >(1/2)n$$ (70) and $`\alpha `$ being so chosen that the $`2`$-point functions are non-trivial. ## 5 Spontaneous Symmetry Breaking (SSB) and the Goldstone Phenomenon ### 5.1 General Remarks Before we study fluctuation operators in the regime of vacuum–, ground–, equilibrium–state degeneracy, we want to briefly comment, in order to set the stage, on the (rigorous) implementation of $`ssb`$ in the various areas with particular emphasis on (quantum) statistical mechanics, i.e. condensed matter physics. As this topic has however been much discussed in the past from various points of views, we do not intend to give an exaustive commentary. We only mention some earlier work being of relevance for our argumentation and sketch the general framework. We assume that our state, $`\omega `$ or $`\mathrm{\Omega }`$, is (non-)invariant under some automorphism group of $`𝒜_0`$ or $`𝒜`$. Furthermore, and this is important (while frequently not clearly stated), we assume the time evolution, $`\alpha _t`$, to commute with the automorphism group, $`\alpha _g`$. ###### Definition 5.1 $`\alpha _g`$ is called a symmetry group if $$\alpha _g\alpha _t=\alpha _t\alpha _g$$ (71) ###### Definition 5.2 If $$(\mathrm{\Omega },\alpha _g(A)\mathrm{\Omega })=(\mathrm{\Omega },A\mathrm{\Omega })$$ (72) for all $`A𝒜`$, the symmetry is called conserved and can be implemented by a unitary group of operators in the representation space $$\alpha _g(A)U(g)AU(g^1)$$ (73) On the other side, if $$(\mathrm{\Omega },\alpha _g(A)\mathrm{\Omega })(\mathrm{\Omega },A\mathrm{\Omega })$$ (74) for some $`A`$, $`A`$ the symmetry-breaking observable, the symmetry is called spontaneously broken since it still commutes with the time evolution (i.e. formally: with the Hamiltonian, modulo boundary terms due to long-range correlations). In most cases the (continuous) symmetry group derives from a clearly identifiable generator (we restrict ourselves, for convenience, to one-parameter groups) which is built from a local operator density, i.e. $$U(s)=e^{isQ},Q(t)=q(x,t)d^nx,Q(t)=Q(0):=Q$$ (75) Note that there are a lot of technical subtleties lurking behind these operator identities, all of which we cannot mention in the following (for more details and references see e.g. . A nice review is where many of the widely scattered results have been compiled ). ###### Remark 5.3 In many situations the generator density is the zero-component of a conserved current. Formally the conservation law encodes the time-independence of the global charge, $`Q`$. Furthermore, for convenience, we assume the symmetry to commute with the space translations, i.e. $`U(x)QU(x)=Q`$. This is in fact frequently the case and simplifies certain calculations. The most crucial consequence is that in case the symmetry is spontaneously broken some of the above relations do only hold in a formal or algebraic sense. More specifically: ###### Theorem 5.4 If $`\alpha _g`$ is spontaneously broken the global generator $`Q`$ does only exist in a formal sense as a limit $$Q=\underset{V}{lim}Q_V,Q_V:=_Vq(x)d^nx$$ (76) We have $$ssb\underset{V}{lim}(\mathrm{\Omega },[Q_V,A]\mathrm{\Omega })0$$ (77) for some $`A𝒜`$ and $`Q`$ is in that case only definable as a nasty operator (see below). In the following we will take (77) as the defining relation of $`ssb`$ (the technical details of the various statements can be found in the literature, mentioned above). The notion of $`ssb`$ is closely connected with another phenomenon, the so-called Goldstone-phenomenon. While there exists a clear picture in, say, relativistic quantum field theory, the corresponding picture is a little bit blurred in the non-relativistic regime. In the relativistic context we have sharp zero-mass Goldstone-modes, i.e. true particles due to relativistic covariance. On the other hand, in e.g. condensed matter physics or statistical mechanics the situation is less generic. In general we do no longer have sharp excitation modes; we have rather to expect excitation modes having a finite lifetime for momentum different from zero but becoming infinitely sharply peaked for momentum $`k0`$. The proper view is it to analyze these excitation branches in the full Fourier-space of energy-momentum as has e.g. been done in ref. four of and earlier in the author’s doctoral thesis, the principal object being the spectral-resolution of the $`2`$-point correlation functions (in a neighborhood of $`(E,k)=(0,0)`$). $`SSB`$ or the Goldstone phenomenon manifests itself in this quantity by a singular contribution in the spectral measure. One should mention at this place the work of Bros and Buchholz (see e.g. ) about quantum field theory in temperature (i.e. KMS-) states. In this particlar context the residual causality and locality properties of the underlying relativistic theory lead to a, in some respects, more generic behavior as compared to the ordinary non-relativistic condensed matter regime. In the non-relativistic regime it turns out that the concrete structure of the Goldstone mode depends usually on the details of the microscopic interactions (that means both the so-called energy-momentum dispersion-law which can be, to give an example, quadratic or linear near $`k=0`$ in the case of magnons or phonons, say, and the $`k`$-dependent width of the branch). This led to the desire to characterize the presence of a Goldstone phenomenon by a simple (if qualitative) property. Sometimes one finds in the literature the saying that the Goldstone phenomenon consists of the vanishing of a mass-gap above the ground state. But this statement is in some sense frequently empty. From we know e.g. that a short-ranged Galilei-covariant theory, with a non-vanishing particle density, cannot have a mass-gap due to phonon-excitations which signal the trivial breaking of the Galilei-boosts. Furthermore, in most cases KMS-Hamiltonians have as spectrum the whole real line. ###### Remark 5.5 Models like the famous BCS-model (having a gap) are no case in point as they are implicitly breaking Galilei-invariance as do all such mean-field-models. This becomes apparent when analyzing the interaction part of the corresponding Hamiltonian. The complete fermion- or boson-liquid is, on the other side, again Galilei-invariant, hence has no mass-gap, but may, of course, still display e.g. superfluidity. In the next subsection we will provide a, as we think, more satisfying and completely general characterization of the Goldstone phenomenon which is independent of the details of the model under discussion. ### 5.2 Some Rigorous Results for the Symmetry Generator in the Presence of SSB After the above introductory remarks we want to prove a couple of rigorous results which characterize to some extent the presence of $`ssb`$ in the (non-)relativistic regime. The main observation is that the symmetry generator is no longer defined as a nice operator in the representation (Hilbert- or $`GNS`$-) space when $`ssb`$ is present and that this, at first glance, mathematical result encodes some interesting physics. Let us work, for simplicity, in the context of temperature states. This has the advantage that $`\mathrm{\Omega }`$ is separating, i.e. $$A\mathrm{\Omega }=B\mathrm{\Omega }A=B$$ (78) The first task is to give $`Q:=lim_VQ_V`$ a rigorous meaning. The standard procedure (see the above mentioned literature) is to define $`Q`$ via: $$QA\mathrm{\Omega }:=\underset{V}{lim}[Q_V,A]\mathrm{\Omega },Q\mathrm{\Omega }:=0$$ (79) for e.g. $`A𝒜_0`$. For $`V`$ sufficiently large, the commutator on the rhs becomes independent of $`V`$, hence there is a chance to get a well-defined $`Q`$ (at least on a dense set of vectors) as on the lhs we have by separability $$A\mathrm{\Omega }=B\mathrm{\Omega }A=B[Q_V,AB]=0$$ (80) For $`A𝒜`$ one has to employ cluster properties. ###### Observation 5.6 We have already seen above that, while such a $`Q`$ may exist, the corresponding $`Q_V\mathrm{\Omega }`$ will nevertheless diverge for $`V^n`$! This shows that the connection between the global generator and its local approximations is not that simple. The best one can usually expect, even in the case of symmetry conservation, is a weak convergence on a dense set $$(B\mathrm{\Omega },QA\mathrm{\Omega })=\underset{V}{lim}(B\mathrm{\Omega },Q_V\mathrm{\Omega })$$ (81) but, due to the above divergence of $`Q_V\mathrm{\Omega }`$, we cannot even have weak convergence on the full Hilbert-space. (For more details see the above cited literature; in particular , third ref., where the various possibilities in the respective fields have been compared) We see from the above that $`Q`$ can be defined as a densely defined operator but usually we want to have more. A conserved continuous symmetry is given by a s.a. generator. Let us see under what conditions the above $`Q`$ is at least symmetric provided that the $`Q_V`$ are symmetric. We assume the symmetry to be conserved, i.e. $$\underset{V}{lim}(\mathrm{\Omega },[Q_V,A]\mathrm{\Omega })=0\text{for all}A𝒜$$ (82) We then have $$\begin{array}{c}(B\mathrm{\Omega },QA\mathrm{\Omega })=\underset{V}{lim}(B\mathrm{\Omega },[Q_V,A]\mathrm{\Omega })\hfill \\ \hfill =\underset{V}{lim}\left(([Q_V,B]\mathrm{\Omega },A\mathrm{\Omega })+(Q_v\mathrm{\Omega },B^{}A\mathrm{\Omega })(A^{}B\mathrm{\Omega },Q_V\mathrm{\Omega })\right)\end{array}$$ (83) ###### Conclusion 5.7 $`Q`$ is symmetric if $`lim_V(A\mathrm{\Omega },Q_V\mathrm{\Omega })=0`$ for all $`A𝒜_0`$. Under the same proviso it follows $$(B\mathrm{\Omega },QA\mathrm{\Omega })=\underset{V}{lim}(B\mathrm{\Omega },Q_V\mathrm{\Omega })$$ (84) What is the situation if the symmetry is spontaneously broken? For convenience we replace again the sharp volume-integration by our smooth one, i.e. $$Q_VQ_R:=q(x)f_R(x)d^nx$$ (85) We know that there exists a symmetry-breaking observable $`A`$ s.t. $$\underset{R}{lim}(\mathrm{\Omega },[Q_R,A]\mathrm{\Omega })0QA\mathrm{\Omega }=\underset{R}{lim}[Q_R,A]\mathrm{\Omega }0$$ (86) Due to the assumed translation invariance, i.e. $$U(a)QU(a)=Q\text{or, what is the same,}U(a)q(x)U(a)=q(x+a)$$ (87) we have $$(\mathrm{\Omega },QA\mathrm{\Omega })=(\mathrm{\Omega },QV^1A_V\mathrm{\Omega })$$ (88) and $$QV^1A_V\mathrm{\Omega }=V^1_VU(x)d^nxQA\mathrm{\Omega }$$ (89) $`U(x)`$ the unitary representation of the translations. Remark: As a result of a discussion with Detlev Buchholz, following a seminar talk about the paper, we will give a technically more detailed proof of the above statement in the appendix at the end of the paper. This seems to be advisable since, as we are showing below, the global operator, $`Q`$, turns out to be non-closable, which will make certain limit-manipulations more cumbersome. ###### Lemma 5.8 $$s\underset{V}{lim}V^1_VU(x)d^nx=P_\mathrm{\Omega }$$ (90) $`P_\mathrm{\Omega }`$ the projector on the (in our case) unique vacuum-,ground-, equilibrium-state. Proof: The result is well-known (see e.g. ). We give however a very short and slightly different proof using our smooth volume integration. With $`V_R:=f_R(x)d^nx`$, a spectral resolution yields $$V_R^1U(x)f_R(x)d^nx=const\left(f(x)d^nx\right)^1\widehat{f}(Rp)𝑑E_p$$ (91) Applied to a vector $`\psi `$ we can now employ Lebesgue’s theorem of dominated convergence and get $$\underset{V}{lim}V^1U(x)d^nx\psi =(\widehat{f}(0))^1\widehat{f}(0)P_\mathrm{\Omega }\psi =P_\mathrm{\Omega }\psi \mathrm{}$$ (92) This yields $$0P_\mathrm{\Omega }QA\mathrm{\Omega }=\underset{V}{lim}QV^1A_V\mathrm{\Omega }$$ (93) On the other hand $$\underset{V}{lim}V^1A_V\mathrm{\Omega }=P_\mathrm{\Omega }A\mathrm{\Omega }=0$$ (94) by an analogous reasoning (note that we assumed $`(\mathrm{\Omega },A\mathrm{\Omega })=0`$). We have now a sequence of vectors, $`V^1A_V\mathrm{\Omega }`$, converging to zero in norm while $`QV^1A_V\mathrm{\Omega }`$ converges to $`P_\mathrm{\Omega }QA\mathrm{\Omega }0`$. Summing up what we have shown we arrive at the following conclusion: ###### Conclusion 5.9 (Goldstone Theorem) If we have $`ssb`$ and a separating vector, $`\mathrm{\Omega }`$, (representing the ground or temperature state), $`Q`$ can still be defined as an operator which is however not closable, hence, a fortiori, not symmetric (note that symmetric operators are closable). This abstract result has as a practical consequence the physical property exhibited in the preceding formulas. They express the content of the Goldstone phenomenon in the most general and model independent way. We infer that $`Q`$ induces transitions from a singular part of the continuous spectrum, passing through $`(E,p)=(0,0)`$, to the extremal invariant state $`\mathrm{\Omega }`$. On the other side, a conserved symmetry implies $$Q\mathrm{\Omega }=0,P_\mathrm{\Omega }[Q,A]\mathrm{\Omega }=0P_\mathrm{\Omega }QA\mathrm{\Omega }=0$$ (95) We show now that the above result really contains the original Goldstone phenomenon. Let us e.g. assume that we have the above result and, on the other side, a gap in the energy spectrum above the state $`\mathrm{\Omega }`$. We emphasized above that an important ingredient of the notion of $`ssb`$ is the time independence of, say, the above expression. We employ again the spectral resolution of operators with respect to energy-momentum. We hence have $$0c=P_\mathrm{\Omega }Q\widehat{A}(k,E)e^{itE}𝑑k𝑑E\mathrm{\Omega }$$ (96) with $`c`$ being independent of $`t`$. We choose a real testfunction $`g(t)`$ with $`g(t)𝑑t=1`$. This yields $$0c=P_\mathrm{\Omega }QA(t)g(t)𝑑t\mathrm{\Omega }=P_\mathrm{\Omega }Q\widehat{A}(E)\widehat{g}(E)𝑑E\mathrm{\Omega }$$ (97) If there is a gap above zero we may choose the support of $`\widehat{g}`$ so that $$supp(\widehat{g})supp(spec(H))=0$$ (98) Since, by assumption, $`P_\mathrm{\Omega }`$ has been extracted in the energy-support of $`A`$, we get the result $`c=0`$, that is, no symmetry breaking. But we can infer more about the nature of the energy-momentum spectrum near $`(0,0)`$. We see that $`P_\mathrm{\Omega }QA(g(t))\mathrm{\Omega }`$ depends only on the value of $`\widehat{g}(E)`$ in $`E=0`$, which is one in our case, but not on the shape of $`g`$. Inspecting equation (93) we can infer the following: The Fourier transform of the rhs contracts around $`k=0`$ in the limit $`V\mathrm{}`$. On the other side we learned that in the limit both sides have their energy support concentrated in $`E=0`$. The lhs shows that the limit vector is parallel to $`\mathrm{\Omega }`$. Whereas we do not want to go into the partly intricate details of the limiting processes of non-closable operators (note that it is e.g. dangerous to use the adjoint, $`Q^{}`$, in the reasoning as it is not densely defined), the latter part of the above theorem should now be obvious. This sharp excitation around $`(E,k)=(0,0)`$ extends into the full energy-momentum plane in form of a (usually) smeared excitation branch (having a finite $`k`$-dependent life-time). For the regime of temperature states the situation was analyzed in some detail in the fourth reference of and already in the authors doctoral thesis. We see from the above that a similar situation prevails in the more general case of a separable $`\mathrm{\Omega }`$ and, analogously, for ground-state models where $`Q`$ can be defined in the above way. Even if the above $`Q`$ is not definable as a non-closable limit operator we arrive at a similar result by exploiting the limit-expectation values instead of the strong vector- or operator limits, but we do not want to dwell more into the corresponding details in this paper which deals with a different topic. ## 6 The Canonical (Goldstone) Pair in the Presence of $`SSB`$ As far as we can see, the notion of a canonical Goldstone pair was introduced by Verbeure et al. in . In the following section we want to prove only a few general (model-independent) results, whereas much more could be shown by combining the framework, developed above, with the techniques mentioned in the preceding section. We remarked above that $`ssb`$ is characterized by the non-vanishing (but time-independence) of the following commutator limit $$0c=\underset{V}{lim}(\mathrm{\Omega },[Q_V,A(t)]\mathrm{\Omega })$$ (99) To fix the notation: usually a pure phase is characterized by the non-vanishing of a so-called order parameter in the presence of $`ssb`$. This is an observable, $`B`$ say, with $$(\mathrm{\Omega },B\mathrm{\Omega })=\{\begin{array}{cc}c0\hfill & \text{in the broken phase}\hfill \\ 0\hfill & \text{in the conserved phase (above }T_c\text{, say)}\hfill \end{array}$$ (100) From (99) we see that as order parameter we have to choose $$B:=\underset{V}{lim}[Q_V,A]$$ (101) while $`A`$ is the symmetry breaking observable. ###### Example 6.1 In the Heisenberg-ferromagnet with spontaneous magnetization in, say, the $`z`$-direction the order parameter is $`S_z`$ or $`S_z`$. As generator of the broken symmetry one may take $`S_x`$ and as symmetry breaking obsrvable e.g. $`S_y`$. We have seen that we can write $$0c=\underset{V}{lim}(\mathrm{\Omega },[Q_V,A]\mathrm{\Omega })=\underset{V}{lim}(\mathrm{\Omega },[Q_V,V^1A_V]\mathrm{\Omega })=\underset{R}{lim}(\mathrm{\Omega },[Q_R,V_R^1A_R]\mathrm{\Omega })$$ (102) where $$Q_R:=q(x)f_R(x)d^nx,A_R:=_{S_R}A(x)d^nx$$ (103) with $`V_R`$ the volume of the sphere, $`S_R`$, with radius $`R`$. We can now split the scaling exponent among the two observables (the volume of the unit sphere being absorbed in the constant). $$0const=\underset{R}{lim}(\mathrm{\Omega },[R^\alpha Q_R,R^{(n\alpha )}A_R]\mathrm{\Omega })$$ (104) This form of scaling may yield something reasonable if the scaling exponents can be so adjusted that also $$(\mathrm{\Omega },R^\alpha Q_RR^\alpha Q\mathrm{\Omega })\text{and}(\mathrm{\Omega },R^{(n\alpha )}AR^{(n\alpha )}A\mathrm{\Omega })$$ (105) remain finite in this limit. In general it does not seem to be easy to get both rigorous and general estimates on the scaling behavior of these quantities. Fortunately, in the case of temperature (KMS) states, such estimates are available. In to the following special (real-space-) version of the Bogoliubov-Inequality has been proved and employed for the observables $`Q_R`$ and $`V_R^1A_R`$: $$\|[Q_R,V_R^1A_R]\|^2V_R^1A_RV_R^1A_R[Q_R,[Q_R,H]]$$ (106) The delicate term is the double commutator on the rhs. If $`Q`$ is spontaneously broken, boundary terms will survive in the commutator of $`Q_R`$ and the Hamiltonian, $`H`$, when taking the limit $`R\mathrm{}`$, while in a formal sense they commute. The double commutator saves us two powers of $`R`$, so to say. That is we arrive after some cumbersome manipulations at $$[Q_R,[Q_R,H]]R^{(n2)}\text{for}R\mathrm{}$$ (107) hence $$V_R^1A_RV_R^1A_RR^{(2n)}\text{for}R\mathrm{}$$ (108) as the limit on the lhs is a constant different from zero in the case of $`ssb`$. ###### Theorem 6.2 For temperature states we have for the symmetry breaking observable $$A_RA_RR^{(n+2)}$$ (109) That is, compared with the ordinary, normal scaling behavior ($`R^n`$), the divergence is worse. From this one infers the following decay of the two-point correlation function itself: $$\|A(x)A\|R^{(n2)}$$ (110) Putting all the pieces together we now have to make the following identification: $$n\alpha (n+2)/2\alpha (n2)/2$$ (111) in order that the limit commutator is non-trivial, i.e. non-classical. On the other hand, the divergence behavior of $`Q_RQ_R`$ can frequently be inferred either from covariance properties (as in relativistic quantum field theory; see e.g. the third reference in ) or from an analysis of the spectral behavior in concrete (non-relativistic) models. Summing up we have: ###### Conclusion 6.3 (Canonical Pair) For a covariant four-current in relativistic quantum field theory the two-point function in Fourier space contains a prefactor $`p^2`$ which yields (after some calculations) an $`\alpha =1/2`$ (for space dimension, $`n=3`$). On the other side, if we do not have such nice covariance properties the divergence of $`Q_RQ_R`$ is generically much worser than $`R`$ (in three dimensions). This holds, in particular, for the above temperature states. It follows that for temperature states we cannot find a critical exponent $`\alpha `$ so that both the auto-correlations remain finite in the limit and the commutator non-trivial. That is, for temperature states the limit fluctuations are classical (an observation already made by Verbeure et al for special models, see e.g. ). The situation seems to be less generic for ground state models, i.e. the temperature-zero case. For one, we do not automatically have an a priori estimate as in the above conclusion, from which we can infer that it is the autocorrelation of $`A_R`$ which is ill-behaved. For another, in temperature states, as was shown in e.g. the fourth reference of by the author, the spectral weight has to become infinite along the Goldstone excitation branch in a specific way (which is governed by the dispersion law of the Goldstone mode) for energy-momentum approaching zero. This sort of singularity is mainly responsible for the poor decay of the respective auto-correlation function. This phenomenon may be absent in the case of ground states as has also been shown for certain Bose-gas models in where some of these questions have been dealt with in greater detail. Note in particular that a variety of aspects may depend on the precise shape of the Goldstone mode near energy-momentum equal to $`(0,0)`$ as was shown in the above mentioned paper of the author or in the unpublished doctoral thesis. On the other side, there has been some interesting work of Pitaevskii and Stringari (see e.g. ), who showed that variants of the uncertainty principle may lead to non-trivial results in certain cases for ground state systems if one can exploit and control certain additional sum rules. ###### Remark 6.4 Note that the ordinary uncertainty principle (for e.g. hermitean operators and ignoring possible domain questions) reads $$1/4\|[A,B]\|^2AABB$$ (112) One sees that instead of the double commutator of the local symmetry generator and the hamiltonian now a term like $`Q_RQ_R`$ occurs. While we have an a priori estimate of the large-R-behavior of the double commutator, the behavior of $`Q_RQ_R`$ is probably less generic (in particular in the ground state situation) and we need some extra information of the kind mentioned above. ## 7 The Case of SSB or Very Poor Decay of Correlations In the preceding sections we studied the case of $`L^1`$\- or $`L^2`$-clustering. In this last section we want to briefly show how we can proceed in the case of extremely poor clustering. We want however, for the sake of brevity and in order to better illustrate the method, to concentrate on the simpler case of a uniformly poor decay of all the correlation functions we are discussing. This is of course not always the case but the scheme can be easily generalized (we discuss this topic in more detail in , where we treat this question in the context of the renormalisation group analysis). We hence assume that the truncated $`l`$-point functions cluster weaker than $`L^2`$ or $`L^1`$, say, in the difference variables, $`y_i:=x_{i+1}x_i`$, (see section 3.1). The following reasoning works both in the case of non-$`L^1`$ or non-$`L^2`$ clustering. In the latter case one would again use the Cauchy-Schwarz-inequality (as in section 3.2). To illustrate the method we choose the non-$`L^1`$ procedure. So let us assume $$W^T(y_1,\mathrm{},y_{l1})L^1$$ (113) For each $`l`$ we assume the existence of a weight factor with a suitable exponent, $`\alpha _l`$: $$P_l(y):=(1+y_i^2)^{\alpha _l/2}$$ (114) so that $$F(y):=P_l(y)^1W^T(y)L^1\text{for}\alpha _l>\alpha _l^{inf}$$ (115) On the other side, we define the fluctuation operators with the exponent $`\gamma `$, which will be adjusted later $$A_R^F:=R^\gamma A_R$$ (116) It follows $$W^T(y)=P_l(y)F(y)$$ (117) with $`F(y)`$ an (in general, $`l`$-dependent) $`L^1`$-function. For the limit correlation functions we then get $$A_R^F(1)\mathrm{}A_R^F(l)^T=R^{ln}R^{l\gamma }\widehat{F}(q)\widehat{P}_l(q)\left[\widehat{f}(Rp_1)\mathrm{}\widehat{f}(Rq_{l1})\right]dq_i$$ (118) (cf. section 3.1) ###### Remark 7.1 We write the Fourier transform of $`P_l(y)`$ formally as $$\widehat{P}_l(q)=(1+D_{q_i}^2)^{\alpha _l/2}$$ (119) (with $`D_{q_i}`$ the partial derivations). For non-integer $`\alpha _l/2`$ this is a pseudo-differential operator. At the moment, for the sake of brevity, we do not want to say more about the corresponding mathematical framework (see for a complete discussion). What we in fact only need are the scaling properties of the expression. If one wants to be careful one may equally well take the explicit expression for the Fourier transform of the above polynomial in the $`y`$-coordinates applied to the product of the $`f_R`$’s and exploit its scaling properties. In any case, we get (with this proviso) and the usual variable transformation $`p_i^{}:=Rp_i`$: $$\begin{array}{c}A_R^F(1)\mathrm{}A_R^F(l)^T=\hfill \\ \hfill R^{lnl\gamma (l1)n+\alpha _l}\widehat{F}(q^{}/R)(R^2+D_{q_i^{}}^2)^{\alpha _l/2}\left[\widehat{f}(p_1^{})\mathrm{}\widehat{f}(q_{l1}^{})\right]dq_i^{}\end{array}$$ (120) Again only the explicit scaling prefactor matters in the limit $`R\mathrm{}`$. (Note that for non-minimal $`\alpha _l`$ we may have $`\widehat{F}(0)=0`$. Technical intricacies like this one will be discussed at length in ). To get a finite result for all correlation functions we have to adjust the scaling parameter, $`\gamma `$, so that the exponents vanish or are negative. We choose $`\alpha _2`$ for $`l=2`$ so that the limit two-point function is finite and non-vanishing. That is: $$n2\gamma +\alpha _2=0\gamma =(n+\alpha _2)/2$$ (121) Inserting this $`\gamma `$ in the general expression for $`l3`$, we conclude that the scaling prefactor is finite in the limit provided that $$\alpha _ll\gamma n=((l1)n+l\alpha _2)/2$$ (122) with $`\gamma `$ fixed by the two-point function. For $`\alpha _l<l\gamma n`$ we can even conclude that all(!) higher limit correlation functions vanish and that the resulting theory is (quasi-)free. The latter would, for example, be the case if $$\alpha _l(l1)\alpha _2$$ (123) holds, since we then have (with $`\alpha _2<n`$): $$\alpha _l(l1)\alpha _2<(l1/2)\alpha _2=(2l1)\alpha _2/2<((l1)n+l\alpha _2)/2$$ (124) but nothing can be concluded in general for, say, $`\alpha _l=l\alpha _2`$. We see that it is of tantamount importance to better understand the assymptotic behavior of truncated $`l`$-point functions and, in particular, the rate of decay as a function of $`l`$. We address this topic in more detail in . ## Appendix The rigorous implementation of the formula $$U(a)q(x)U(a)=q(x+a)$$ (125) is $$\begin{array}{c}U(a)Q_RU(a)=U(a)q(x)f_R(x)d^nxU(a)=q(x+a)f_R(x)d^nx\hfill \\ \hfill =q(y)f_R(ya)d^ny=:Q_R(a)\end{array}$$ (126) The first question is: how does the global $`Q`$ behave under translations? To answer this question we have to take recourse to the definition of the global $`Q`$ as a limit of local operations. We have $$U(a)QA\mathrm{\Omega }=U(a)\underset{R}{lim}[Q_R,A]\mathrm{\Omega }=\underset{R}{lim}[Q_R(a),A(a)]\mathrm{\Omega }$$ (127) since it holds $$\underset{n}{lim}U(a)\psi _n=U(a)\underset{n}{lim}\psi _n$$ (128) as $`U(a)`$ is bounded. If $`A`$ is local we have for sufficiently large $`R`$ (and hence, in the limit): $$\underset{R}{lim}[(Q_R(a)Q_R(0)),A(a)]=0$$ (129) We hence arive at $$U(a)QA\mathrm{\Omega }=\underset{R}{lim}[Q_R,A(a)]\mathrm{\Omega }=QA(a)\mathrm{\Omega }=QU(a)A\mathrm{\Omega }$$ (130) ###### Lemma 7.2 On the dense set $`𝒜_0\mathrm{\Omega }`$, $`Q`$ commutes with the translations. In a next step we have to analyse the action of $`Q`$ on integrals or averages like $`_VU(x)AU(x)d^nx\mathrm{\Omega }`$. More specifically, we want to show that $`Q`$ commutes, so to speak, with the operation of integration. We have $$Q_VA(x)d^nx\mathrm{\Omega }:=\underset{R}{lim}[Q_R,_VA(x)d^nx]\mathrm{\Omega }$$ (131) We approximate the integral by a sum, that is: $$_VA(x)d^nx\psi :=\underset{i}{lim}\underset{i}{}d^nx_iA(x_i)\psi $$ (132) and get (as the $`Q_R`$ are assumed to be nice, that is, closed operators) $$[Q_R,_VA(x)d^nx]\mathrm{\Omega }=\underset{i}{lim}[Q_R,\underset{i}{}d^nx_iA(x_i)]\mathrm{\Omega }=\underset{i}{lim}\underset{i}{}d^nx_iU(x_i)[Q_R(x_i),A]\mathrm{\Omega }$$ (133) We again choose $`R`$ so large that $$[Q_R(x),A]=[Q_R,A]\text{for all }xV$$ (134) which leads to $$[Q_R,_VA(x)d^nx]\mathrm{\Omega }=\underset{i}{lim}\underset{i}{}d^nx_iU(x_i)[Q_R,A]\mathrm{\Omega }=_VU(x)d^nx[Q_R,A]\mathrm{\Omega }$$ (135) Taking now the limit $`R\mathrm{}`$, we get ###### Lemma 7.3 $$Q_VA(x)d^nx\mathrm{\Omega }=_VU(x)d^nxQA\mathrm{\Omega }$$ (136) This shows, that our manipulations can be justified.
warning/0003/hep-th0003216.html	ar5iv	text	# The Interplay between 𝜃 and 𝑇 ## 1 Introduction Noncommutative field theories have proven to exhibit fascinating properties (-). As we noted in a recent paper , there is a drastic reduction in the number of degrees of freedom that contribute to the free energy in the non-planar sector, $`F_{\mathrm{np}}`$, at high temperatures. By the free energy in the non-planar sector, we mean those contributions to the free energy which arise from non-planar Feynman graphs. This reduction can be read off by looking at the temperature dependence of $`F_{\mathrm{np}}`$. In this paper, we will present further evidence for this reduction, and study two further theories with very different UV sensitivity to test how generic this high temperature behavior is. We will take $`\theta _{0i}=0`$ throughout. In , we also noted the existence of winding states in noncommutative field theories. Though we have no clean way of treating these modes separately from the conventional degrees of freedom of commutative field theory, it appears plausible that $`F_{\mathrm{np}}`$ is indicative of these additional degrees of freedom, especially at high temperature, at which the winding states become light. In studying the reduction of the degrees of freedom that contribute to $`F_{\mathrm{np}}`$, we emphasize that this phenomenon does not depend on the sensitivity of the corresponding commutative theory to the ultraviolet . Indeed, this reduction occurs in theories that are extremely insensitive to the ultraviolet, like the $`N=4`$, $`D=4`$ SYM theory, as well as theories with strong ultraviolet sensitivity, as eg. $`\lambda \varphi ^3`$ in $`D=6`$. In a nutshell, what we think happens is that as the temperature rises above the noncommutativity scale, the modes of the various fields with momenta of order the temperature, $`T`$, do not contribute to $`F_{\mathrm{np}}`$. This occurs because, as the thermal wavelengths drop below the noncommutativity scale $`\theta ^{1/2}`$, as far as their contribution to $`F_{\mathrm{np}}`$ is concerned, there is no way to distinguish these high momenta modes from each other and thus no way to count their contribution individually. Another way of thinking about this is that the contributions to $`F_{\mathrm{np}}`$ are sensitive to the phase space structure of space: the uncertainty principle among the spatial coordinates $`x_i`$ renders the identification of the degrees of freedom with regard to their contribution to $`F_{\mathrm{np}}`$ impossible; this contribution seems to be summarized by a few quantum degrees of freedom contributing per noncommutative cell whose area is $`\theta `$. In the following, we will refer to this cell as the Moyal cell. In the next section, we briefly review the results for the free energy of the Wess-Zumino model. We introduce a tool to test whether a reduction of degrees of freedom contributing to $`F_{\mathrm{np}}`$ in fact takes place, the calculation of the classical statistical mechanics approximation to F, and discuss when we expect this approximation to be valid. In the third section, we discuss the case of $`N=4`$, $`D=4`$ SYM. This theory also shows a reduction of degrees of freedom at high temperatures, although the theory is rather insensitive to the UV. We see that the nonplanar sector again contains winding states. In the fourth section, we show how these results generalize to higher dimensions where there can be more noncommutative spatial directions. We show that the classical approximation remains valid for $`\lambda \varphi ^3`$ in spite of the high sensitivity of the corresponding commutative theory to the UV. Finally, we end with conclusions. ## 2 The thermodynamics of the noncommutative Wess-Zumino model in perturbation theory The Lagrangian density for the Wess-Zumino model is: $$=i_\mu \overline{\psi }\overline{\sigma }^\mu \psi +A^{}A\frac{1}{2}M\psi \psi \frac{1}{2}M\overline{\psi }\overline{\psi }g\psi \psi Ag\overline{\psi }\overline{\psi }A^{}F^{}F,$$ where $`F`$ is given by $$F=MA^{}gA^{}A^{}.$$ We showed in a previous paper that in the noncommutative Wess-Zumino model, one can distinguish two regimes of high temperature, $`\beta M1`$.<sup>1</sup><sup>1</sup>1We will assume that the Compton wavelength of the fields is bigger than the noncommutative scale, $`\theta M^21`$. These two regimes are: $`\theta T^21`$ and $`\theta T^21`$. In the case $`\theta T^21`$, the thermal wavelength is larger than the noncommutativity scale. Here, $`F_{\mathrm{np}}`$ behaves as $$\frac{F}{V}\|_{\mathrm{np}}g^2T^4.$$ (1) In the other limit, the free energy density behaves as $$\frac{F}{V}\|_{\mathrm{np}}g^2\frac{T^2}{\theta }\mathrm{log}T^2\theta .$$ (2) This expression for the free energy shows a dramatic reduction of the degrees of freedom contributing to $`F_{\mathrm{np}}`$. Such a reduction has a natural interpretation in terms of the novel phase space structure due to the noncommutative nature of space. At any temperature, the typical momenta of the fields have magnitudes of order the temperature or smaller. As the temperature rises, the thermal wavelengths become smaller. When the temperature reaches values of $`O(\theta ^{1/2})`$, modes of the field with momenta of $`O(T)`$ no longer contribute to $`F_{\mathrm{np}}`$. More precisely, what we think happens is that at temperatures $`T^21/\theta `$, the modes $`\varphi _{k_1,k_2,k_3}`$ with $`k_{1,2}\theta ^{1/2}`$ cannot be distinguished from each other insofar as their contribution to $`F_{\mathrm{np}}`$ is concerned, and therefore will not be counted individually in the non-planar contribution to the partition function. These modes lose their individual identity in the non-planar sector because the noncommutativity between $`x_1`$ and $`x_2`$, $`[x_1,x_2]0`$, implies an uncertainty relationship which renders impossible their separate identification. This may explain the reduction in the temperature dependence of the free energy from $`T^4`$ to $`T^2/\theta `$. The non-planar sector yields a contribution to the free energy which, could this sector be isolated, would suggest that it has the number of degrees of freedom of a 1+1 dimensional field theory at temperatures $`T1/\theta ^{1/2}`$. This is reminiscent of the result obtained by Atick and Witten in the context of string theory. There is also a subleading, logarithmic dependence on the temperature in equation (2) for the free energy, which is all that remains from the contributions of the high momenta components of the fields. As announced in the introduction, one can compare this computation to a classical statistical mechanics calculation of the free energy. The classical calculation can be done by just keeping the zero frequency term in the sum over frequencies. The result is $$\frac{F}{V}\|_{\mathrm{np}}g^2\frac{T^2}{\theta }\mathrm{log}\mathrm{\Lambda },$$ (3) where $`\mathrm{\Lambda }`$ is an ultraviolet cutoff. The classical calculation captures the correct power law dependence in the temperature but replaces the logarithmic dependence on temperature obtained in quantum mechanics by $`\mathrm{log}\mathrm{\Lambda }`$. We should remind the reader that classical statistical mechanics is a good approximation at high temperatures for the free energy of systems with few degrees of freedom. For example, the high temperature behavior of the free energy for a particle of mass $`m`$ in a potential $`V(x)`$ is well approximated by classical statistical mechanics when the thermal de-Broglie wavelength $`\lambda _{DB}\frac{1}{(mT)^{1/2}}`$ is smaller than the length scale over which the potential varies. In the path integral formulation, this would correspond to actually shrinking the temporal circle down to zero size and performing a dimensional reduction. In Feynman diagrams, this amounts to only keeping the contributions of the zero frequency components in the sum over frequencies. In contrast, if one uses classical statistical mechanics to evaluate the free energy of systems with a field theory number of degrees of freedom at high temperature, one inevitably encounters a UV catastrophe. On the other hand, classical statistical mechanics is a good approximation in field theory when calculating thermal correlation functions. An example is the calculation of the two point correlation function between two operators, $`O_1(x)`$ and $`O_2(y)`$, $$O_1(x)O_2(y)O_1(x)O_2(y).$$ Indeed, the modes that primarily contribute to this correlation function have wavelength commensurate with the distance $`\|xy\|`$ separating the probes. As one heats the system, the population of these modes increases and hence their contribution to the correlation function is well approximated by classical statistical mechanics. So the classical approximation at high temperature becomes applicable in field theory when there is a bound on the wavelengths of the modes that contribute to the quantity we are calculating. The fact that $`F_{\mathrm{np}}`$ is well approximated by classical statistical mechanics thus suggests that only a reduced set of degrees of freedom contributes to $`F_{\mathrm{np}}`$ at high temperatures. This should be contrasted to the cases of field theories with an arbitrary infinite series of higher derivative terms in the Lagrangian. In such generic cases, if one attempts to approximate the free energy by a classical statistical calculation, one encounters enhanced UV catastrophes as compared to the theory without the higher derivative terms. The averting of the UV catastrophe due to the special choice of higher derivative terms in noncommutative field theories is therefore quite remarkable. ## 3 The thermodynamics of $`N=4`$, $`D=4`$ Super Yang-Mills Our interest in examining this case is to see how crucial the ultraviolet sensitivity of a noncommutative theory is to the existence of winding states and the high temperature behavior of $`F_{\mathrm{np}}`$. We will show that the existence of winding states as well as the reduction at high temperatures in the number of degrees of freedom contributing to $`F_{\mathrm{np}}`$ are independent of the UV behaviour of the noncommutative field theory. The calculation of the free energy is straightforward, the result can be found in the literature . We present for completeness some of the steps in the calculation of the free energy in appendix A. The contribution to the free energy coming from the nonplanar sector is: $$\frac{F}{V}\|_{\mathrm{np}}=4g^2N\frac{d^3p}{(2\pi )^3}\frac{d^3k}{(2\pi )^3}\frac{e^{ip\theta k}}{4\omega _p\omega _k}\left(n_B(\omega _p)+n_F(\omega _p)\right)\left(n_B(\omega _k)+n_F(\omega _k)\right).$$ (4) Up to an overall factor, this is exactly the result we found in the Wess-Zumino case . The presence of winding states can again be detected in the non-planar sector. Indeed, by performing the integration over one set of momenta in eq. (4), one finds contributions to the free energy that are weighted by the length of the temporal circle. This can be seen for example by considering the following factor in eq. (4): $$\frac{d^3k}{(2\pi )^3}\frac{e^{ip\theta k}}{\omega _k}n_B(\omega _k)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n^2\beta ^2+(\theta p)^2}.$$ (5) We can again compare the limits $`T^2\theta 1`$ and $`T^2\theta 1`$. In the case where the thermal wavelength is larger than the noncommutativity scale, the free energy behavior is: $$\frac{F}{V}\|_{\mathrm{np}}g^2NT^4.$$ (6) In the limit where the thermal wavelength is within the noncommmutative scale, $$\frac{F}{V}\|_{\mathrm{np}}g^2N\frac{T^2}{\theta }\mathrm{log}(T^2\theta ).$$ (7) This reveals again a reduction in the number of degrees of freedom contributing to $`F_{\mathrm{np}}`$. There seems to be a universal behavior that whenever the inverse momenta fit into a Moyal cell, the associated components of the field cannot be distinguished in their contributions to $`F_{\mathrm{np}}`$. What is then left over seems to be the contribution due to a few quantum degrees of freedom per Moyal cell. This picture will be further tested in the next section, where we will consider a higher dimensional example. ## 4 The thermodynamics of $`g\varphi ^3`$ in $`D=6`$ dimensions Strictly speaking, this theory does not have good thermodynamic behavior since it lacks a ground state. Therefore, in order to discuss the thermodynamics of this system, we will take the coupling constant to be very small and the temperature to be small enough<sup>2</sup><sup>2</sup>2while still keeping it much larger than the mass such that the excursions of the field $`\varphi `$ away from the local minimum of the potential are within the bounded region of the potential. This implies that the temperature satisfies the inequality $`T\frac{M}{g^{1/2}}`$. The Lagrangian density for this system is: $$=d^6x((\varphi )^2M^2\varphi ^2g\varphi \varphi \varphi ).$$ The free energy in this case has no infrared divergences to leading order in an expansion in the coupling constant g. This is because the high dimensionality of the theory softens the IR sensitivity of the free energy. The non-planar contribution to the free energy at $`O(g^2)`$ is: $$\frac{F}{V}\|_{\mathrm{np}}=g^2T^2\underset{n,l}{}\frac{d^5p}{(2\pi )^5}\frac{d^5k}{(2\pi )^5}\frac{e^{ip\theta k}}{(\frac{4\pi ^2n^2}{\beta ^2}+k^2)(\frac{4\pi ^2l^2}{\beta ^2}+p^2)(\frac{4\pi ^2(n+l)^2}{\beta ^2}+(k+p)^2)}.$$ (8) Again, one finds winding states in this theory. This can easily be seen by rewriting eq. (8) in the form<sup>3</sup><sup>3</sup>3 for details, see appendix B: $$\frac{F}{V}\|_{\mathrm{np}}=g^2T\underset{n,l}{}\frac{d^5p}{(2\pi )^5}_0^1𝑑x_0^{\mathrm{}}\frac{d\alpha _1}{\alpha _1^2}𝑑\alpha _2e^{(\alpha _2+\alpha _1x(1x))(p^2+\frac{4\pi ^2n^2}{\beta ^2})}e^{\frac{(\theta p)^2+4\pi ^2l^2\beta ^2}{\alpha _1}}e^{2\pi inlx},$$ (9) where we can think of $`\alpha _2+\alpha _1x(1x)`$ as the proper time for the propagator of momentum states and $`\frac{1}{\alpha _1}`$ as the proper time for the propogator of winding states. Under the restriction that $`\theta _{0i}=0`$, the most general case in $`D=6`$ is, by a convenient choice of coordinate system, given by nonvanishing $`\theta _{12}`$ and $`\theta _{34}`$. If we take these two parameters to be of the same order of magnitude, one can distinguish two cases: 1. $`T^2\frac{1}{\theta _{12}},\frac{1}{\theta _{34}}`$ 2. $`T^2\frac{1}{\theta _{12}},\frac{1}{\theta _{34}}`$. In the first case, the contribution to the free energy from the nonplanar sector scales, as a function of temperature, like $`T^6`$. This is because the exponential involving the Moyal phase does not oscillate much when the momenta are distributed according to the thermal distributions. In the high temperature limit where $`T^2\frac{1}{\theta _{12}},\frac{1}{\theta _{34}}`$, the possiblity of estimating the behavior of the free energy using classical physics arises. Indeed, as discussed above, if the system has many fewer degrees of freedom contributing to the $`\theta `$ dependence of the free energy, as was the case in the four dimensional examples, then this approximation is valid. Whether or not the approximation is valid is decided a posteriori: if the remaining integrals are finite, then indeed classical statistical mechanics is a good approximation and gives the dominant contribution at high temperature. Performing the classical statistical mechanics calculation, i.e. evaluating the following expression, $$\frac{F}{V}\|_{\mathrm{np}}g^2T^2\frac{d^5k}{(2\pi )^5}\frac{d^5p}{(2\pi )^5}\frac{e^{ip\theta k}}{p^2k^2(p+k)^2},$$ (10) we find no UV divergences. When $`\theta \theta _{12}\theta _{34}`$, $$\frac{F}{V}\|_{\mathrm{np}}g^2\frac{T^2}{\theta _{34}^{}{}_{}{}^{2}\theta _{12}^{}{}_{}{}^{2}}\mathrm{log}\frac{\theta _{12}}{\theta _{34}}\frac{g^2T^2}{\theta ^2}.$$ (11) This behavior is again consistent with the picture on how degrees of freedom, as their inverse momenta fall into Moyal cells, do not contribute to $`F_{\mathrm{np}}`$. ## 5 Conclusions At high temperatures, one observes a substantial reduction in the number of degrees of freedom that contribute to $`F_{\mathrm{np}}`$, the free energy due to the non-planar sector of the theory. The picture that emerges from the previous sections is that this phenomenon can be traced to the modes of the fields with momenta larger than the noncommutativity scales. What happens is that once the wavelengths are within the Moyal cells, there is no way to distinguish and count the separate contributions of these modes to $`F_{\mathrm{np}}`$. What is left over seems to be the contribution of a single degree of freedom per Moyal cell. Because of this severe reduction in the number of degrees of freedom contributing to $`F_{\mathrm{np}}`$, this quantity can be calculated at high temperature using classical statistical mechanics. This behavior does not appear to depend on the details of the noncommutative field theory, the crucial element is the existence of a phase space structure for space. ###### Acknowledgments. We thank Nathan Seiberg for very useful discussions. The work of WF, EG, RMcN, AK-P, SP, PP is supported in part by the Robert Welch Foundation and the NSF under grant number PHY-9511632. SP is also supported by NSF grant PHY-9973543. ## Appendix A The free energy of $`N=4`$, $`D=4`$ SYM We are working in Euclidean spacetime. Let’s begin with the commutative case. We choose the decomposition of the 32 dimensional $`\mathrm{\Gamma }`$-matrices given in . This yields the dimensionally reduced Lagrangian $`=`$ $`2\mathrm{Tr}`$ $`({\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }{\displaystyle \frac{i}{2}}\overline{\chi }_K\overline{\sigma }^\mu D_\mu \chi _K{\displaystyle \frac{1}{2}}D_\mu \varphi _iD^\mu \varphi _i{\displaystyle \frac{1}{2}}D_\mu \phi _iD^\mu \phi _i`$ $`{\displaystyle \frac{i}{2}}g\alpha _{KL}^i(\chi _K[\varphi _i,\chi _L]+\overline{\chi }_K[\varphi _i,\overline{\chi }_L]){\displaystyle \frac{i}{2}}g\beta _{KL}^i(\chi _K[\phi _i,\chi _L]+\overline{\chi }_K[\phi _i,\overline{\chi }_L])`$ $`{\displaystyle \frac{1}{4}}g^2([\varphi _i,\varphi _j][\varphi _i,\varphi _j]+[\phi _i,\phi _j][\phi _i,\phi _j]+2[\varphi _i,\phi _j][\varphi _i,\phi _j])),`$ where $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu ig[A_\mu ,A_\nu ]`$, $`D_\mu =_\mu +ig[A_\mu ,]`$ and $`\mu `$, $`\nu `$ range from 1 to 4. Remember that this theory contains one N=1 vector multiplet and three N=1 chiral multiplets. Thus $`i`$ ranges from 1 to 3, $`K`$ and $`L`$ from 1 to 4. $`\gamma _\mu `$, $`\alpha _i`$ and $`\beta _i`$ satisfy $`\{\gamma _\mu ,\gamma _\nu \}=2\delta _{\mu \nu }`$, $`\{\alpha _i,\alpha _j\}=2\delta _{ij}`$, $`\{\beta _i,\beta _j\}=2\delta _{ij}`$ and commute among each other. The diagrams that contribute to the free energy at two loop are depicted in the following figure: At $`T=0`$, their contribution must vanish. It is easy to see that each diagram is proportional to $$g^2f_{abc}f_{abc}\frac{d^4k}{(2\pi )^4}\frac{d^4p}{(2\pi )^4}\frac{1}{k^2p^2}.$$ The coefficients of the various contributions are given in the following table: $$\begin{array}{cccccccc}& & & & & & & \\ 1& 2& 3& 4& 5& 6& 7& 8\\ & & & & & & & \\ 3& \frac{9}{4}& 4& \frac{1}{4}& \frac{9}{2}& \frac{9}{2}& 12& 0\end{array}$$ (13) We pass to finite temperature by replacing the integration over energies by a sum over even or odd Matsubara frequencies, for bosonic, fermionic degrees of freedom respectively. The noncommutative case does not require any additional calculation. We obtain the noncommutative theory by replacing the commutators in the Lagrangian (A) by Moyal brackets. This yields, for arbitrary fields $`A`$, $`B`$, which take values in the fundamental representation, $$[A,B]_{}(x)=\frac{d^4p}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}\stackrel{~}{A}^a(p)\stackrel{~}{B}^b(k)e^{i(p+k)x}\left([t^a,t^b]\mathrm{cos}\frac{k\theta p}{2}+i\{t^a,t^b\}\mathrm{sin}\frac{k\theta p}{2}\right).$$ For the gauge group $`U(N)`$, both the commutator and the anticommutator of generators of the fundamental close within the algebra: $$[t^a,t^b]=if_{abc}t^c,\{t^a,t^b\}=d_{abc}t^c.$$ The noncommutative result is thus obtained from the commutative one by replacing $$f_{abc}f_{abc}\mathrm{cos}\frac{k\theta p}{2}+d_{abc}\mathrm{sin}\frac{k\theta p}{2}.$$ To obtain (4), we need to perform the sums $`f_{abc}f_{abc}`$ and $`d_{abc}d_{abc}`$. The first is of course identical to the $`SU(N)`$ result, $`N(N^21)`$. The second can be obtained from the SU(N) result, $`N(N\frac{4}{N})`$, where $`\{t^a,t^b\}=d_{abc}^{SU(N)}t^c+\frac{1}{N}\delta _{ab}`$, by including the $`U(1)`$ generator $`t^{N^2}=\frac{1}{\sqrt{2N}}1`$. This gives $`N(N^2+1)`$. ## Appendix B The free energy of the $`D=6`$, $`g\varphi ^3`$ theory We will first briefly show how the winding states appear in this case, then we will perform the classical statistical mechanics calculation of the non-planar contribution to $`O(g^2)`$ to the free energy. The nonplanar contribution to the free energy to $`O(g^2)`$ is $$g^2T^2\underset{n,l}{}\frac{d^5p}{(2\pi )^5}\frac{d^5k}{(2\pi )^5}\frac{e^{i(\theta _{12}(p_1k_2p_2k_1)+\theta _{34}(p_3k_4p_4k_3))}}{(p^2+\frac{4\pi ^2n^2}{\beta ^2})(k^2+\frac{4\pi ^2l^2}{\beta ^2})((p+k)^2+\frac{4\pi ^2(n+l)^2}{\beta ^2})}.$$ (14) The appearance of winding contributions in this integral can be shown by introducing a Feynman parameter, $`x`$, and Schwinger parameters $`\alpha _1`$ and $`\alpha _2`$: $`g^2T^2{\displaystyle \underset{n,l}{}}{\displaystyle \frac{d^5p}{(2\pi )^5}\frac{d^5k}{(2\pi )^5}}`$ $`{\displaystyle _0^1}𝑑x{\displaystyle _0^{\mathrm{}}}𝑑\alpha _1𝑑\alpha _2\alpha _1`$ $`e^{\alpha _2(p^2+\frac{4\pi ^2n^2}{\beta ^2})}e^{\alpha _1k^2\alpha _1(p^2x+2pkx+4\pi ^2\frac{n^2x+2nlx+l^2}{\beta ^2})}e^{ip\theta k}.`$ Performing the integral over $`k`$ followed by a Poisson resummation over $`l`$ gives $$\frac{F}{V}\|_{\mathrm{np}}=g^2T\underset{n,l}{}\frac{d^5p}{(2\pi )^5}_0^1𝑑x_0^{\mathrm{}}\frac{d\alpha _1}{\alpha _1^2}𝑑\alpha _2e^{(\alpha _2+\alpha _1x(1x))(p^2+\frac{4\pi ^2n^2}{\beta ^2})}e^{\frac{(\theta p)^2+4\pi ^2l^2\beta ^2}{\alpha _1}}e^{2\pi inlx},$$ (16) which is equation (9). We now turn to evaluating the dominant contribution to this expression at high temperature. As explained in the text, this is given by the $`n=l=0`$ contribution in equation (14). We will for simplicity limit our discussion to this mode. We introduce Feynman and Schwinger parameters: $`g^2T^2{\displaystyle _0^{\mathrm{}}}𝑑\alpha \alpha ^2{\displaystyle _0^1}𝑑x{\displaystyle _0^{1x}}𝑑y{\displaystyle d^5pd^5ke^{\alpha ((1y)p^2+(1x)k^2+2(1xy)p.k)+ip\theta k}}.`$ (17) The Gaussian integrals can now be performed with the result: $`g^2T^2{\displaystyle _0^{\mathrm{}}}\alpha 𝑑\alpha {\displaystyle _0^1}𝑑x{\displaystyle _0^{1x}}𝑑y{\displaystyle \frac{1}{f(x,y)^{1/2}}}{\displaystyle \frac{1}{\alpha ^2f(x,y)+\theta _{12}^2}}{\displaystyle \frac{1}{\alpha ^2f(x,y)+\theta _{34}^2}},`$ (18) where $`f(x,y)=x(1x)+y(1y)xy`$. The Feynman and Schwinger integrals can also be done exactly and they give: $$g^2T^2\frac{1}{\theta _{12}^2\theta _{34}^2}\mathrm{log}\left(\theta _{12}/\theta _{34}\right).$$ (19) This is the general result when $`M/g^{1/2}TM`$. We note in passing that the zero temperature vacuum energy of this theory is finite: $$g^2\frac{d^6p}{(2\pi )^6}\frac{d^6k}{(2\pi )^6}\frac{e^{i(\theta _{12}(p_1k_2p_2k_1)+\theta _{34}(p_3k_4p_4k_3))}}{p^2k^2(p+k)^2}=g^2\frac{1}{\theta _{12}+\theta _{34}}\frac{1}{\theta _{12}\theta _{34}}.$$ (20)
warning/0003/cond-mat0003391.html	ar5iv	text	# Stochastic Aggregation: Scaling Properties ## I Introduction In the previous study, we introduced a stochastic aggregation process involving both active and passive clusters. We generalized Smoluchowski’s rate equations and obtained exact results for several kernels. In this study, we apply stochastic aggregation to reaction-diffusion, coarsening, and ballistic agglomeration problems. Our goal is to examine the range of validity the mean-field results, and to determine whether the overall scaling behavior extends to low dimensional systems. The rate equations approach is mean-field in nature, i.e., it is valid only when spatial correlations are absent. Formally, it is applicable only in infinite spatial dimension, or in the presence of an effective mixing mechanism. This mean-field theory should also be asymptotically exact when the spatial dimension is sufficiently high. In low spatial dimensions, however, significant spatial correlations eventually develop, and the rate equation approach does not apply in the long time limit. We therefore focus on one-dimensional systems where spatial correlations are most pronounced. We performed numerical simulations of stochastic aggregation processes with both diffusive and ballistic particle transport. The simulations show that the scaling behavior suggested by the mean-field theory is indeed generic, as it extends to one-dimensional systems. We find that two nontrivial model-dependent exponents characterize the survival probabilities of the particles and monomers, respectively. Smoluchowski’s theory provides reasonable estimates for these exponents. Additionally, we studied the limiting mass distribution of passive clusters. Surprisingly, over a substantial mass range, this distribution depends only weakly on the underlying transport mechanism. Furthermore, mean-field theory provides an excellent approximation for the limiting mass distribution. The rest of this paper is organized as follows. The general scaling behavior is outlined in Sec. II. Predictions of the mean-field theory are summarized in Sec. III. Numerical simulations of stochastic aggregation processes with diffusive and ballistic transport mechanisms are described in Secs. IV and V, respectively. A discussion of the results is presented in Sec. VI. ## II Scaling Properties Stochastic aggregation involves two types of clusters: active and passive. Initially, the system consists of active monomers only. When two active clusters merge, the newly-born aggregate remains active with probability $`p`$, or becomes passive (i.e., it never aggregates again) with probability $`q=1p`$. Eventually, all active clusters are depleted and the system consists of passive clusters only. This process can be viewed as an aggregation-annihilation process since it interpolates between aggregation ($`p=1`$) and annihilation ($`p=0`$). Quantities of interest include $`A_k(t)`$ and $`P_k(t)`$, the distributions of active and passive clusters at time $`t`$, as well as the final distribution of passive clusters, $`P_k(\mathrm{})`$. As shown in , two conservation laws underly this system. The first is mass conservation, $`k[A_k(t)+P_k(t)]=\mathrm{const}`$. The second conservation law reflects the fact that changes in the overall densities are coupled, $`qA(t)+(1+q)P(t)=\mathrm{const}`$, where $`A(t)=A_k(t)`$ and $`P(t)=P_k(t)`$ are the number densities of active and passive clusters, respectively. Therefore, it is sufficient to study the time evolution of the number density and the mass density of active clusters, $`A(t)`$ and $`M(t)=kA_k(t)`$, respectively. The latter quantity is the survival probability of an active particle, i.e., the probability that it still belongs to an active cluster at time $`t`$. Smoluchowski’s theory suggests that both quantities decay algebraically in the long time limit $$A(t)t^\nu ,M(t)t^{\nu \psi }.$$ (1) As will be shown below, this as well as other scaling properties suggested by this theory hold qualitatively even for low dimensional stochastic aggregation processes. While the decay exponent $`\nu `$ is typically robust in that it depends only on the major characteristics of the process such as the spatial dimension or the transport mechanism, the exponent $`\psi \psi (p)`$ is non-universal as it depends on the details of the model, i.e., on the probability $`p`$. In turn, this implies a non-universal growth law for the average mass of an active cluster $`k=M/At^{\nu (1\psi )}`$. For the system to follow a scaling behavior, the average mass must be the only relevant scale in the long time limit, and conversely, any scale characterizing the initial mass distribution must be “erased” eventually. In other words, the mass distribution is characterized by a single rescaled variable $$A_k(t)t^{\nu (\psi 2)}F\left(kt^{\nu (\psi 1)}\right),$$ (2) with the time dependent prefactor fixed by the decay laws (1). This scaling behavior is similar to that found for deterministic aggregation-annihilation processes and for aggregation-annihilation of domains in coarsening processes. These studies suggest that another independent exponent describes the decay of small clusters. Specifically, the monomer density decays according to $$A_1(t)t^{\nu \delta },$$ (3) with a model-dependent exponent $`\delta \delta (p)`$. The monomer density decay reflects the small argument behavior of the scaling function $`F(\xi )\xi ^\sigma `$ with $`\delta 1=(1\psi )(1+\sigma )`$. One of our main results is that the mass distribution of active clusters is described by a set of nontrivial exponents $`(\psi ,\delta )`$. These exponents can be viewed as persistence exponents as they characterize the survival probability of an active particle, and an active monomer . Several properties of the scaling exponents are general. For instance, the inequalities $`\psi 1\delta `$ hold since $`A_1A_kkA_k`$. The two exponents are equal $`\psi =\delta =1`$ in the annihilation case ($`p=0`$), since $`A_k(t)=A(t)\delta _{k,1}`$. In the aggregation limit ($`p=1`$) the mass density of active clusters is conserved and therefore $`\psi =0`$. We now turn to the mass distribution of passive clusters. The Smoluchowski theory suggests that the same scaling form underlies both mass distributions $$P_k(t)t^{\nu (\psi 2)}G\left(kt^{\nu (\psi 1)}\right).$$ (4) In contrast with the active cluster distribution, the passive cluster distribution approaches a nontrivial final distribution $`P_k(\mathrm{})`$. Such a time independent final distribution is consistent with the above scaling form only when the scaling function diverges, $`F(\xi )\xi ^\gamma `$ in the limit $`\xi 0`$, with $`\gamma =(2\psi )/(1\psi )`$. As a result, the final mass distribution of passive clusters decays algebraically in the large mass limit $$P_k(\mathrm{})k^\gamma \mathrm{with}\gamma =\frac{2\psi }{1\psi }.$$ (5) At a given time $`t`$, this decay is realized for clusters whose mass $`k`$ does not exceed the characteristic mass $`k^{}t^{\nu (1\psi )}`$. Note also that $`0<\psi <1`$ implies $`2<\gamma <\mathrm{}`$. Generally, the mass conservation restricts the large mass decay exponent to $`\gamma >2`$. Since the $`\psi `$ exponent varies between $`0`$ and $`1`$, we see that the entire range of acceptable exponents is realized by tuning the probability $`p`$. ## III Mean-Field Theory It is well established that spatial correlations can be safely neglected only in spatial dimensions larger than some upper critical dimension, $`d>d_c`$. For example, for irreversible aggregation with mass-independent diffusion and reaction rates, one has $`d_c=2`$; for a general aggregation process, however, the upper critical dimension may be arbitrarily large. Below the upper critical dimension, substantial spatial correlations develop, and the most important features including the scaling exponents and the scaling functions are changed. Generally, the lower the spatial dimension, the larger the difference with the mean-field predictions. Although the Smoluchowski rate equations approach does not apply in low spatial dimensions, it can still serve as a useful approximation after an appropriate modification. This can be accomplished by replacing the overall reaction rate with an effective density-dependent reaction rate $`rr(A)`$ $`{\displaystyle \frac{dA_k}{dt}}`$ $`=`$ $`r\left({\displaystyle \frac{p}{2}}{\displaystyle \underset{i+j=k}{}}A_iA_jA_kA\right),`$ (6) $`{\displaystyle \frac{dP_k}{dt}}`$ $`=`$ $`r\left({\displaystyle \frac{q}{2}}{\displaystyle \underset{i+j=k}{}}A_iA_j\right).`$ (7) We are primarily interested in situations where aggregation is independent of the mass, and therefore we use a mass independent rate kernel. The reaction rate $`r(A)`$ is model dependent. In reaction-diffusion processes, the reaction rate decays algebraically with the density (see, e.g., Refs.). Assuming $`r(A)A^\alpha `$ yields $`\frac{dA}{dt}A^{\alpha +2}`$, and consequently, the density decay exponent is found $`\nu ={\displaystyle \frac{1}{1+\alpha }}.`$ (8) In general, a reduction in the reaction rate, i.e., $`\alpha >0`$, leads to a slowing down in the density decay rate, $`\nu <1`$. Apart from the change in $`\nu `$, all other aspects of this approximation are identical to the Smoluchowski theory with a constant rate kernel. Indeed, the above rate equations reduce to the Smoluchowski’s rate equations with a redefined time variable, $`t\tau =_0^t𝑑t^{}r(t^{})`$. In particular, the scaling exponents $`\psi `$ and $`\delta `$ are independent of $`\alpha `$: $$\psi =2\frac{1p}{2p},\delta =\frac{2}{2p}.$$ (9) One can verify the expected limiting behaviors $`\psi (1)=0`$, and $`\psi (0)=\delta (0)=1`$. Furthermore, the scaling functions are as in the constant kernel solution , and for example, $`F(\xi )`$ is purely exponential. The corresponding small argument exponents $`\gamma =2/p`$ and $`\sigma =0`$ follow from $`\psi `$ and $`\delta `$ using the aforementioned scaling relations. The final mass distribution of passive clusters is independent of the reaction rate $`r`$ $$P_k(\mathrm{})=\frac{q}{p}\frac{\mathrm{\Gamma }(1+2/p)\mathrm{\Gamma }(k)}{\mathrm{\Gamma }(k+2/p)}.$$ (10) Below, we compare these mean-field predictions with simulation results for one-dimensional stochastic aggregation where spatial correlations are most pronounced. We also examine the role of the aggregates’ transport mechanism by considering both diffusive and ballistic transport. ## IV Diffusive Transport In diffusive stochastic aggregation, identical particles are placed onto a $`d`$-dimensional lattice. All particles perform independent random walks, i.e., they hop to a randomly chosen nearest-neighbor site with a constant rate. If this site is occupied, the two particles coalesce irreversibly, and with probability $`p`$ the resulting aggregate remains active, while with probability $`q=1p`$ it becomes passive. Effectively, passive particles are removed from the system. In the case of single-species reaction diffusion processes, the effective reaction rate can be obtained from dimensional analysis. Eq. (6) implies $`[r]=[L]^d[T]^1`$, and since the reaction rate can only be a function of the diffusion coefficient $`[D]=[L]^2[T]^1`$ and the density $`[A]=[L]^d`$, one finds $`rDA^{(2d)/d}`$. Hence, $`\alpha =(2d)/d`$ and Eq. (8) yields the correct decay exponents $`\nu =d/2`$ below the upper critical dimension $`d_c=2`$. To examine the above scaling picture we performed numerical simulations of diffusive stochastic aggregation processes in one dimension. Unless noted otherwise, the data was obtained from an average over 10 independent realizations in a system of size $`L=10^7`$ with periodic boundary conditions. Initially, all sites were occupied. First, we verified that the number density, the mass density, and the monomer density indeed decay algebraically in the long time limit, in accord with Eqs. (1) and (3). The case $`p=1/2`$ is shown in Fig. 1, and the corresponding decay exponents were found: $`\nu =0.500(1)`$, $`\psi =0.6193(3)`$, and $`\delta =1.460(2)`$. Mean-field theory correctly predicts $`\nu =1/2`$. Furthermore, the predictions $`\psi =2/3`$ and $`\delta =4/3`$ provide a reasonable approximation. One can compare with the case of disordered (Sinai) diffusion where a real-space decimation procedure was used to determine exact values of these exponents. Remarkably, the disorder case exponent $`\psi =0.61937`$ is in excellent agreement with the simulation value. There is a small discrepancy with the second exponent $`\delta =1.47041`$. Additionally, we verified that the densities of active and passive clusters follow the scaling forms of Eqs. (2) and (4), respectively (see Fig. 2). In agreement with the mean-field theory, the scaling functions decay exponentially for large masses. We also studied how the exponents vary with the probability $`p`$, as shown in Figs. 3 and 4. The exact exponents found for the disordered case by Le Doussal and Monthus provide an excellent approximation (within 0.1%) for $`\psi `$. In the case of $`\delta `$, a different behavior emerges in the aggregation limit, $`p1`$, where the exact value is $`\delta =3`$ , and the disagreement with both mean-field theory and the disordered case are most pronounced. The above scaling arguments suggest that the limiting mass distribution of passive clusters decays algebraically with the exponent $`\gamma =(2\psi )/(1\psi )`$. For $`p=1/2`$, one therefore expects $`\gamma 3.627`$ (compare with $`\gamma =3.62722`$ and $`\gamma =4`$, predicted by the disordered case and the mean-field theory). This corresponds to a very strong suppression of large masses, and therefore, it is much more difficult to confirm this behavior numerically. Nevertheless, our simulations (Fig. 5) are consistent with a power law decay with an exponent $`\gamma 3.6`$. In one dimension, the diffusion-controlled stochastic aggregation is equivalent to the Potts model with zero-temperature Glauber dynamics. For the $`Q`$-state Potts model with spatially uncorrelated initial conditions, aggregation of domain walls occurs with probability $`p=\frac{Q2}{Q1}`$, and annihilation occurs with probability $`q=\frac{1}{Q1}`$. Therefore, the above can be reformulated in terms of domain walls rather than aggregates. In the coarsening context, the domain wall mass (or number) distribution is dual to domain number distribution . ## V Ballistic Transport The situation when particles move ballistically involves several complications. First, while the annihilation limit is uniquely defined , the aggregation limit has various realizations. In traffic flows, the velocity of a newly-born cluster is the smaller of the two velocities , while in application to astrophysics and granular gases the velocity follows from momentum conservation . Second, numerical results for the annihilation case and analytical results for the traffic case show that the initial conditions are remembered forever, in contrast with the diffusive case. Specifically, the small velocity characteristics of the initial velocity distribution influence the long time asymptotic behavior, including the scaling exponents. We consider the momentum conserving case, also known as “ballistic aggregation” or “sticky gas” . The initial velocities are assigned according to the distribution $`P_0(v)`$. The mass (momentum) of a newly-born cluster is equal to the sum of masses (momenta) of the two colliding clusters. After an agglomeration event, the newborn particle remains active with probability $`p`$, or becomes passive with probability $`q=1p`$. To apply the Smoluchowski rate equations approach, we again use dimensional analysis to calculate the decay exponent $`\nu `$. The collision rate is $`rva^{d1}`$, where $`v`$ is the typical velocity and $`a`$ is the typical radius of an aggregate. A particle of radius $`a`$ contains of the order $`a^d`$ monomers whose initial momenta are uncorrelated. Momentum conservation therefore implies $`va^{d/2}`$. Using $`a^dM/AA^{\psi 1}`$ gives the collision rate $`ra^{(d2)/2}A^{(d2)(\psi 1)/2d}`$. From Eq. (8) one finds $$\nu =\frac{2d}{d+2+\psi (d2)},$$ (11) with $`\psi `$ given by Eq. (9). Apart from the exponent $`\nu `$, features such as the exponential mass distribution and the exponents $`\psi `$ and $`\delta `$ are given by the mean-field theory outlined above. In two dimensions, the collision rate does not depend on $`\psi `$ and hence, the asymptotic behavior $`At^1`$ agrees with that found for deterministic ballistic agglomeration. For $`d2`$, stochastic and deterministic asymptotics differ: stochasticity enhances decay of the number density $`A`$ for $`d<2`$ and weakens it for $`d>2`$. A more detailed mean-field theory can be carried. It yields a factorizing joint mass-velocity distribution, with an exponential mass distribution, and a Gaussian velocity distribution . In the aggregation case, $`\psi =0`$ and therefore the correct scaling exponent $`\nu =2d/(d+2)`$ is recovered from Eq. (11). For the annihilation case, however, initial conditions are “remembered” forever and therefore the above dimensional arguments no longer hold. The predicted exponent in the annihilation case is always mean-field $`\nu =1`$, while one-dimensional numerical simulations yield an exponent continuously varying from $`0`$ to $`1`$ depending on the initial velocity distribution $`P_0(v)`$, e.g., $`\nu 0.8`$ for uniform initial distributions. We have simulated the stochastic aggregation process on a one-dimensional ring with $`10^6`$ particles. The initial velocity distribution was uniform in $`[1,1]`$. We measured the scaling exponent $`\psi `$ via the scaling relation $`MA^\psi `$, rather than directly versus time, since the exponent $`\nu (p)`$ is not known analytically. We have found that the mean-field prediction, $`\psi =(22p)/(2p)`$, provides a reasonable approximation for the exponent $`\psi `$, as shown in Fig. 6. Furthermore, this approximation should improve in higher dimensions. We compared the mean-field prediction for the mass distribution of passive clusters, Eq. (10), with the numerically obtained distributions in both ballistic and diffusive cases. Interestingly, the rate equations provide an excellent approximation for small and moderate masses (see Fig. 7). Given that the discrepancy in $`\psi `$ is maximal for the case $`p=1/2`$, one may expect an even better approximation for other values of $`p`$. Noting the strong decay of this distribution, the contribution of very large masses is extremely small; for example, $`P_{100}(\mathrm{})2.4\times 10^7`$ for $`p=1/2`$. Hence, the most pronounced part of the distribution is well approximated by the rate equations theory. Surprisingly, the transport mechanism does not play an important role as far as the final mass distribution of passive clusters is concerned. ## VI Discussion We have investigated diffusion- and ballistic-controlled stochastic aggregation in one dimension. We have seen that the rate equations approach captures the overall scaling behavior and additionally it provides reasonable estimates for the decay exponents. In general, the mass distribution is characterized by two nontrivial model-dependent decay exponents. In the diffusion-controlled case, the exponent $`\psi `$ underlying the survival probability of a particle is in excellent agreement with the exact results from the disordered case. In fact, one cannot dismiss the possibility that the disordered and the pure values are identical, based on numerics alone. However, there is an evident discrepancy in the exponent $`\delta `$ as the disordered case exponent diverges logarithmically in the aggregation limit. Stochastic aggregation is equivalent to domain coarsening in the zero-temperature Potts-Glauber model. The above exponents $`(\psi ,\delta )`$ characterize the domain wall number distribution in analogy with $`(\psi _D,\delta _D)`$ for the domain number distribution . In the latter case as well, exact values calculated for the disordered case provide an excellent approximation for the domain exponents. In general, the particle survival probability exponent $`\psi `$ is robust, while the monomer survival probability exponent $`\delta `$ is very sensitive to the details of the process. In the ballistic-controlled case, we have shown that even in the absence of a consistent mean-field theory, some characteristics such as the exponent $`\psi `$ are well approximated by the rate equations. Understanding of reaction processes with an underlying ballistic transport remains largely incomplete. The asymptotic behavior is highly sensitive to the initial conditions, and the critical dimension is apparently infinite. In fact, exact analytical results are available mostly in the aggregation limit . The most intriguing property of the stochastic aggregation is the profound lack of universality. Indeed, the weak dependence on the transport mechanism is in contrast with the strong dependence on the parameter $`p`$. For example, our numerical results show that the final distribution of passive clusters is very close in diffusion- and ballistic-controlled situations. Another very impressive manifestation of this is the excellent agreement between the values of the exponent $`\psi (p)`$ in the disordered and pure cases. This research was supported by the DOE (W-7405-ENG-36), NSF (DMR9632059), and ARO (DAAH04-96-1-0114).
warning/0003/cond-mat0003395.html	ar5iv	text	# From localization to delocalization in the quantum Coulomb glass ## I Introduction The discovery of a metal-insulator transition (MIT) in the two-dimensional electron gas in Si-MOSFETs has induced renewed attention to the transport properties of disordered electrons. This MIT is in conflict with the theory of localization for non-interacting electrons which predicts that all states are localized in 2D. The electron density in the Si-MOSFETs is very low which makes the electron-electron interaction particularly important. Thus it is generally assumed that some type of interaction effect is responsible for this MIT. One of the most remarkable findings about the MIT in Si-MOSFETs is that an in-plane magnetic field (which does not couple to the orbital motion of the electrons) strongly suppresses the conducting phase . This suggests that the spin degrees of freedom play an important role for the transition. A complete understanding has, however, not yet been obtained. There have been attempts to explain the experiments based on the perturbative renormalization group , non-perturbative effects , or the transition being a superconductor-insulator transition rather than a MIT . In order to attack the problem of disordered interacting electrons numerically we have developed an efficient method, the Hartree-Fock based diagonalization (HFD) which is related to the quantum-chemical configuration interaction approach. We have used this method to study the influence of interactions on the conductance in one , two , and three dimensions. We found a delocalizing tendency of the interactions for strong disorder but a localizing one for weak disorder. Similar results have been obtained by means of the density-matrix renormalization group in one dimension and exact diagonalization in two dimensions . Since in most of the numerical work in the literature spinless electrons were considered, there are not many results about the importance of the spin degrees of freedom. In this work we address this question by generalizing the HFD method to spinful electrons. We then use it study the influence of the spin degrees of freedom on the Kubo-Greenwood conductance. ## II Model and method The generic model for spinless interacting disordered electrons is the quantum Coulomb glass . In this paper we use a straight-forward generalization of the quantum Coulomb glass to spinful electrons. It is defined on a regular hypercubic lattice with $`g=L^d`$ ($`d`$ is the spatial dimensionality) sites occupied by $`N=N_{}+N_{}=2Kg`$ electrons ($`0<K<1`$). To ensure charge neutrality each lattice site carries a compensating positive charge of $`2Ke`$. The Hamiltonian is given by $$H=t\underset{ij,\sigma }{}(c_{i\sigma }^{}c_{j\sigma }+h.c.)+\underset{i,\sigma }{}\phi _in_{i\sigma }+\frac{1}{2}\underset{ij,\sigma ,\sigma ^{}}{}(n_{i\sigma }K)(n_{j\sigma ^{}}K)U_{ij}+U_H\underset{i}{}n_in_i$$ (1) where $`c_{i\sigma }^{}`$ and $`c_{i\sigma }`$ are the creation and annihilation operators at site $`i`$ and spin $`\sigma `$, and $`ij`$ denotes all pairs of nearest-neighbor sites. $`t`$ is the strength of the hopping term, i.e., the kinetic energy, and $`n_{i\sigma }`$ is the occupation number of spin state $`\sigma `$ at site $`i`$. We parametrize the interaction $`U_{ij}=e^2/r_{ij}`$ by its value $`U`$ between nearest-neighbor sites. The Coulomb repulsion between two electrons at the same site is described by the Hubbard interaction $`U_H`$ The random potential values $`\phi _i`$ are chosen independently from a box distribution of width $`2W_0`$ and zero mean. The boundary conditions are periodic and the Coulomb interaction is treated in the minimum image convention (which implies a cut-off at a distance of $`L/2`$). A numerically exact solution of a quantum many-particle system requires the diagonalization of a matrix whose dimension increases exponentially with system size. This severely limits the possible sample sizes. In order to overcome this problem we have developed the HFD method. The basic idea is to work in a truncated Hilbert space consisting of the corresponding Hartree-Fock (Slater) ground state and the low-lying excited Slater states. For each disorder configuration three steps have to be performed: (i) find the Hartree-Fock solution of the problem, (ii) determine the $`B`$ Slater states with the lowest energies, and (iii) calculate and diagonalize the Hamilton matrix in the subspace spanned by these states. The number $`B`$ of new basis states determines the quality of the approximation, reasonable values have to be found empirically. ## III Results The conductance of a quantum many-particle system can be obtained from linear-response theory. It is essentially determined by the current-current correlation function of the ground state. The real (dissipative) part of the conductance (in units of $`e^2/h`$) is given by the Kubo-Greenwood formula , $$\mathrm{Re}G^{xx}(\omega )=\frac{2\pi ^2}{\omega }\underset{\nu }{}\|0\|j^x\|\nu \|^2\delta (\omega +E_0E_\nu )$$ (2) where $`\omega `$ denotes the frequency. $`j^x`$ is the $`x`$ component of the current operator and $`\nu `$ denotes the eigenstates of the Hamiltonian. Eq. (2) describes an isolated system while in a real d.c. transport experiment the sample is connected to contacts and leads. This results in a finite life time $`\tau `$ of the eigenstates leading to an inhomogeneous broadening $`\gamma =\tau ^1`$ of the $`\delta `$ functions in (2) . To suppress the discreteness of the spectrum of a finite system, $`\gamma `$ should be at least of the order of the single-particle level spacing. For our systems this requires a comparatively large $`\gamma 0.05`$. We tested different $`\gamma `$ and found that the conductance values depend on $`\gamma `$ but the qualitative results do not. The main results of this paper are summarized in Fig. 1 which shows the disorder and interaction dependence of the typical conductance for both spinless and spinful electrons on a two-dimensional lattice of $`4\times 4`$ sites. The qualitative behavior in both cases is similar: In the strongly localized regime (small $`t`$) a moderate interaction delocalizes the electrons while a sufficiently strong interaction always strongly suppress the conductance. This is the precursor of a Wigner crystal or Wigner glass. With decreasing disorder (increasing $`t`$) the interaction-induced enhancement of the conductance also decreases and eventually vanishes. The behavior of the conductance can be attributed to the competition of two effects: First, the interactions destroy the phase of the electrons and thus the interference necessary for localization. This is particularly effective if the localization length is small to begin with. Second, the interactions introduce an additional source of randomness which tends to increase the localization. A comparison of the cases of spinless and spinful electrons shows that the interaction induced delocalization is significantly larger for spinful electrons. Moreover, the enhancement seems to vanish at a larger kinetic energy (which we did not reach in the simulations). A systematic investigation of the dependence of the conductance on $`U`$ and $`U_H`$ will be published elsewhere. In summary, we have studied the influence of electron-electron interactions on Anderson localization for spinless and spinful electrons in two dimensions. For strong disorder moderate interactions significantly enhance the transport. This enhancement is much stronger for spinful than for spinless electrons. Identifying a real phase transition and thus establishing a connection between these findings and the experiments on Si-MOSFETs requires a finite-size scaling analysis of the conductance. This remains a task for the future. This work was supported in part by the NSF under grant no. DMR–98–70597 and by the DFG under grant no. SFB 393/C2. T.V. thanks the Aspen Center for Physics and the University of Oregon for hospitality during the completion of this paper.
warning/0003/cond-mat0003255.html	ar5iv	text	# Interaction effects on the conductance peak height statistics in quantum dots ## Abstract A random interaction matrix model is used to study the statistics of conductance peak heights in Coulomb blockade quantum dots. When the single-particle dynamics conserves time-reversal symmetry, the peak height statistics is insensitive to the interaction strength. But when the single-particle dynamics breaks time-reversal symmetry, the peak height statistics exhibits a crossover from unitary to orthogonal symmetry as the interaction strength increases. This crossover is driven by the time-reversal symmetry of the interaction. Our random interaction matrix model describes features of both the measured peak height and peak spacing statistics. Random matrix theory (RMT) provides a useful tool for describing the universal statistical fluctuations of the spectrum and eigenfunctions of a quantum system whose associated classical dynamics is chaotic. RMT has been successfully applied to the study of mesoscopic phenomena in quantum dots – submicron 2D devices where electrons are confined by electrostatic potentials . In dots with irregular shapes the single-electron dynamics is mostly chaotic, and the use of RMT is justified within a single-particle framework. In dots that are strongly coupled to leads, or open dots, the approximation of non-interacting quasi-particles is reasonable, and interactions can be considered indirectly through their effect on the electron coherence time. When a finite dephasing time is included, RMT can describe quantitatively the observed conductance statistics in open dots . As the dot-leads coupling is made weaker, the charge on the dot becomes quantized, and electron-electron interactions cannot be ignored. In such almost closed dots, the conductance displays sharp peaks versus the gate voltage. The measured peak height distributions were found to agree with RMT predictions . However, it is not clear why RMT should describe the peak height statistics in dots with strong electron-electron interactions. For example, the spacings between peaks were observed to have a Gaussian-like distribution and not the Wigner-Dyson distribution that is expected in the constant interaction model plus single-particle RMT. Numerical simulations of an Anderson model of a small disordered dot ($`10`$ electrons) with Coulomb interactions also found Gaussian peak spacing distributions , while the peak height distributions showed only a weak dependence on interactions . RMT is not limited to non-interacting systems. Indeed it was first introduced to explain the neutron resonance data in the compound nucleus. But the neutron resonances are measured at finite excitations, while the linear conductance experiments in quantum dots probe the ground state of the system with a varying number of electrons. To understand the statistics that emerges from the interplay between one-body chaos and interactions in these dots, it is necessary to construct a random matrix model that includes interactions explicitly and reduces to one-body RMT in the absence of interactions . We can break the time-reversal symmetry of the one-body Hamiltonian (e.g., with a magnetic field), but the time-reversal symmetry of the two-body interaction should be preserved. An important question is whether such a random matrix model can reproduce both the measured peak height and peak spacing statistics. A random interaction matrix model (RIMM) was introduced in Ref. to study the peak spacing distribution in chaotic dots. The model combines a random two-body interaction with a random one-body Hamiltonian, and describes a crossover of the peak spacing statistics from a Wigner-Dyson distribution to a Gaussian distribution in terms of a parameter that measures the fluctuations of the interaction matrix elements. Here, we use the RIMM to study the effects of interactions on the conductance peak height statistics. For a GOE one-body statistics, we find that the partial width (to decay into one of the leads) has a GOE Porter-Thomas distribution independently of the interaction strength. However, for a GUE one-body statistics, we find a crossover from a GUE to a GOE Porter-Thomas distribution as the interaction strength increases. It is well known that, in the absence of interactions, a complete crossover from GOE to GUE statistics can be induced by a magnetic field. Our results indicate that this crossover is not complete when strong interactions are turned on because of a competition between an asymptotic symmetry of the one-body Hamiltonian and the GOE symmetry of the time-reversal–conserving two-body interaction. Finally, for the case of GUE one-body statistics, we find that a Gaussian-like peak spacing distribution develops at interaction strengths for which the peak height statistics is still close to GUE. This explains the observed peak spacing statistics and peak height statistics within a single random matrix model. The $`n`$-th conductance peak height $`G_n`$ corresponds to the tunneling of an electron into a dot with $`n1`$ electrons to form a dot with $`n`$ electrons. At low temperatures, $`G_n=(e^2/h)(\pi \overline{\mathrm{\Gamma }}/4kT)g_n`$, where $`g_n={\displaystyle \frac{1}{\overline{\mathrm{\Gamma }}}}{\displaystyle \frac{\mathrm{\Gamma }_n^l\mathrm{\Gamma }_n^r}{\mathrm{\Gamma }_n^l+\mathrm{\Gamma }_n^r}}.`$ (1) $`\mathrm{\Gamma }_n^{l(r)}`$ is the partial width of the ground state of the $`n`$-electron dot to decay into an electron in the left (right) lead and the ground state of the dot with $`n1`$ electrons: $$\mathrm{\Gamma }_n\left\|\mathrm{\Phi }_{\mathrm{g}.\mathrm{s}.}(n)\|\psi ^{}(𝒓)\|\mathrm{\Phi }_{\mathrm{g}.\mathrm{s}.}(n1)\right\|^2.$$ (2) $`\psi ^{}(𝒓)`$ is the creation operator of an electron at the point $`𝒓`$ \[$`𝒓=𝒓_{l(r)}`$ for the left (right) point contact\], and $`\mathrm{\Phi }_{\mathrm{g}.\mathrm{s}.}(n)`$ is the ground state wavefunction of the $`n`$-electron dot. For non-interacting electrons, $`\mathrm{\Phi }_{\mathrm{g}.\mathrm{s}.}(n)`$ is a Slater determinant of the $`n`$ lowest single-particle eigenfunctions in the dot, and Eq. (2) reduces to $`\mathrm{\Gamma }_n\|\varphi _n(𝒓)\|^2`$, where $`\varphi _n`$ is the $`n`$-th single-particle wavefunction. If the single-particle dynamics is chaotic, $`\|\varphi _n(𝒓)\|^2`$ satisfies Porter-Thomas statistics, leading to universal distributions of the conductance peak heights (1) that are sensitive only to the underlying symmetry class . To determine how interactions might modify the Porter-Thomas statistics of the partial widths, we use the RIMM of spinless interacting fermions $`H={\displaystyle \underset{ij}{}}h_{ij}a_i^{}a_j+{\displaystyle \frac{1}{4}}{\displaystyle \underset{ijkl}{}}\overline{u}_{ijkl}a_i^{}a_j^{}a_la_k.`$ (3) The one-body matrix elements $`h_{ij}`$ are chosen from the appropriate Gaussian random matrix ensemble, while the anti-symmetrized two-body matrix elements $`\overline{u}_{ij;kl}=u_{ij;kl}u_{ij;lk}`$ form a GOE in the two-particle space $`P(h)e^{\frac{\beta }{2a^2}\mathrm{Tr}h^2};P(\overline{u})e^{\mathrm{Tr}\overline{u}^2/2U^2}.`$ (4) The states $`\|i=a_i^{}\|0`$ describe a fixed basis of $`m`$ single-particle states. $`h`$ is an $`m\times m`$ GOE (GUE) matrix when the single-particle dynamics conserves (breaks) time-reversal symmetry. The two-body interaction is assumed to preserve time-reversal symmetry and forms a GOE, irrespective of the symmetry of the one-body Hamiltonian. To calculate the statistics of the partial widths $`\mathrm{\Gamma }_n`$, we expand $`\psi ^{}(𝒓)=_i\psi _i(𝒓)a_i^{}`$, where $`\psi _i(𝒓)𝒓\|i`$ is the wavefunction of the fixed state $`\|i`$. It follows from the orthogonal invariance of the ensemble (4) that the statistics of $`\mathrm{\Gamma }_n`$ is identical to the statistics of $$\mathrm{\Gamma }_n^i\left\|\mathrm{\Phi }_{\mathrm{g}.\mathrm{s}.}(n)\|a_i^{}\|\mathrm{\Phi }_{\mathrm{g}.\mathrm{s}.}(n1)\right\|^2$$ (5) for any $`i`$. For each realization $`H`$ of the ensemble (3), we calculate the ground states for $`n1`$ and $`n`$ electrons, and compute $`\mathrm{\Gamma }_n^i`$ using (5). The distributions of the normalized width $`\widehat{\mathrm{\Gamma }}=\mathrm{\Gamma }/\overline{\mathrm{\Gamma }}`$ for a GOE single-particle statistics are shown in Fig. 1 for several values of $`U/\mathrm{\Delta }`$. We show $`P(\mathrm{ln}\widehat{\mathrm{\Gamma }})`$ rather than $`P(\widehat{\mathrm{\Gamma }})`$ in order to display more clearly the small values of $`\widehat{\mathrm{\Gamma }}`$. The distributions are independent of $`U/\mathrm{\Delta }`$ and are well described by the GOE Porter-Thomas distribution $`P_{\mathrm{GOE}}(\mathrm{ln}\widehat{\mathrm{\Gamma }})=(\widehat{\mathrm{\Gamma }}/2)^{1/2}\mathrm{exp}(\widehat{\mathrm{\Gamma }}/2)`$ (solid line). As a reference we also show the GUE Porter-Thomas distribution $`P_{\mathrm{GUE}}(\mathrm{ln}\widehat{\mathrm{\Gamma }})=\widehat{\mathrm{\Gamma }}\mathrm{exp}(\widehat{\mathrm{\Gamma }})`$ (dashed line). The conductance peak heights are calculated from (1) using two uncorrelated “sites” $`i`$ and $`j`$ for the left and right leads. The peak height distributions $`P(g)`$ (for a GOE $`h`$) are shown in the insets of Fig. 1. They are all in good agreement with the GOE distribution $`P_{\mathrm{GOE}}(g)=\sqrt{2/\pi g}e^{2g}`$ (solid lines) irrespective of $`U/\mathrm{\Delta }`$. The dashed lines describe the GUE distribution $`P(g)=4g[K_0(2g)+K_1(2g)]\text{e}^{2g}`$ ($`K_0`$ and $`K_1`$ are Bessel functions). The results for a GUE one-body statistics are shown in Fig. 2 for the same values of $`U/\mathrm{\Delta }`$ as in Fig. 1. When the interaction strength increases, the width and peak height distributions make a crossover from the corresponding GUE distribution (at $`U=0`$) to the GOE distribution (at large $`U`$). Equivalently, the transition from GOE to GUE statistics due to a time-reversal symmetry-breaking one-body field is not complete because of the competing GOE symmetry of the two-body interaction. The crossover distributions are compared with distributions obtained from an RMT ensemble that describes the crossover between the orthogonal and unitary symmetries. This ensemble is $`H=S+i\alpha A`$ , where $`S`$ and $`A`$ are $`N\times N`$ symmetric and antisymmetric uncorrelated Gaussian matrices and $`\alpha `$ is a parameter . Its wavefunction statistics depends on the parameter $`\lambda \alpha \sqrt{N}/\pi `$. In particular, the width distribution is given in closed form by $`P_\lambda (\widehat{\mathrm{\Gamma }})={\displaystyle _0^1}dtP_\lambda (t){\displaystyle \frac{1+t^2}{t}}e^{\left(\frac{1+t^2}{t}\right)^2\widehat{\mathrm{\Gamma }}}I_0\left({\displaystyle \frac{1t^4}{t^2}}\widehat{\mathrm{\Gamma }}\right),`$ (6) where $`I_0`$ is the modified Bessel function of order zero. $`t`$ in (6) is a “shape” parameter that fluctuates in the interval $`[0,1]`$ according to a known distribution $`P_\lambda (t)=\pi ^2`$ $`(1/t^3t)\lambda ^2e^{\frac{\pi ^2}{2}\lambda ^2\left(t1/t\right)^2}\{\varphi _1(\lambda )+`$ (8) $`[(t+t^1)^2/41/(2\pi ^2\lambda ^2)][1\varphi _1(\lambda )]\},`$ where $`\varphi _1(\lambda )=_0^1dye^{2\pi ^2\lambda ^2(1y^2)}`$. Eq. (6) reduces to the GOE and GUE Porter-Thomas distributions for $`\lambda =0`$ and $`\lambda 1`$, respectively. In practice, the crossover from the GOE to the GUE already occurs for $`\lambda 1`$. For each $`U/\mathrm{\Delta }`$, we find $`\lambda `$ by fitting the distribution (6) to the computed width distribution. The distributions (6), shown by the short-dashed lines in Fig. 2, accurately describe the width distributions of the RIMM with a GUE $`h`$. Good agreement with closed RMT expressions is also obtained for the peak height distributions $`P_\lambda (g)`$ shown in the insets of Fig. 2. An important issue is the universality associated with the RIMM (3) and (4). The model depends on three parameters: $`m`$, $`n`$, and $`U/\mathrm{\Delta }`$. The top panel of Fig. 3 shows, for a GUE one-body $`h`$, the crossover parameter $`\lambda `$ as a function of $`U/\mathrm{\Delta }`$ for $`m=10`$ and $`n=4,5,6,7`$ (symbols), and for a “reference” case $`m=12`$ and $`n=4`$ (solid line). The curves depend on both $`m`$ and $`n`$, but can all be scaled on the reference curve after scaling the interaction strength by a constant, $`U_{\mathrm{eff}}f(m,n)U/\mathrm{\Delta }`$ (see bottom panel of Fig. 3). The values of the scaling factor $`f(m,n)`$, shown in the bottom inset of Fig. 3, are essentially the same as those needed to obtain a universal peak spacing statistics . We conclude that the peak spacing statistics as well as the peak height statistics in the RIMM depend only on a single parameter $`U_{\mathrm{eff}}`$. In the range $`U_{\mathrm{eff}}1`$, $`\lambda 1`$ and the statistics is essentially GUE. For $`U_{\mathrm{eff}}11.5`$, the peak height statistics is close to GUE while the peak spacings already follow a Gaussian distribution. This explains the measured RMT-like peak height distributions and the Gaussian-like shape of the peak spacing distributions within a single random matrix model. We remark that the small deviations from GUE statistics observed in the experiment of Ref. and in the calculations of Ref. in the presence of a magentic field are consistent with a crossover from GUE to GOE. The average width $`\overline{\mathrm{\Gamma }}`$ is a monotonically decreasing function of $`U/\mathrm{\Delta }`$ and saturates at large values of $`U`$. The top inset of Fig. 3 shows $`\overline{\mathrm{\Gamma }}(U)/\overline{\mathrm{\Gamma }}(0)`$ as a function of $`U/\mathrm{\Delta }`$ for several $`m`$ and $`n`$. This dependence is well described by $`\overline{\mathrm{\Gamma }}(U)\overline{\mathrm{\Gamma }}(\mathrm{})=(\overline{\mathrm{\Gamma }}(0)\overline{\mathrm{\Gamma }}(\mathrm{}))/(1+bU^2/\mathrm{\Delta }^2)`$. The parameter $`\overline{\mathrm{\Gamma }}(\mathrm{})`$ depends on $`m`$ and $`n`$, while $`b`$ is found to be independent of $`m`$ and to depend only weakly on $`n`$. Next we compare the predictions of the RIMM with those of a model of a quantum dot. We studied a 2D Anderson model with on-site disorder parameter $`W`$ and hopping matrix element $`V=1`$. The electrons are interacting with a Coulomb interaction whose strength over one lattice spacing $`a`$ is $`U_c=e^2/a`$ . We choose periodic boundary conditions in both directions so that the average width and width statistics are independent of the specific sites to which the leads are attached. With these boundary conditions, a background charge (that was included in the calculations of Ref. ) only shifts the energy levels by an $`n`$-dependent constant but does not affect the wavefunctions. The left panel of Fig. 4 shows the width statistics in the absence of magnetic flux for a $`4\times 5`$ lattice with disorder parameter $`W=5`$ and for $`n=4`$ electrons. The distributions are approximately described by the GOE Porter-Thomas distribution (solid line) irrespective of the value of $`U_c`$ and in agreement with the RIMM. The dependence of the average width $`\overline{\mathrm{\Gamma }}`$ on the interaction strength (not shown) is similar to that observed in the RIMM. The right panel shows width distributions in the presence of magnetic flux $`\mathrm{\Phi }=0.14\mathrm{\Phi }_0`$ and for $`W=3`$. The $`U_c=0`$ distribution (circles) agrees with the GUE Porter-Thomas distribution (dashed line), while the $`U_c=12`$ distribution (triangles) is described by (6) with $`\lambda =0.17`$. This is the distribution obtained in the RIMM for $`U_{\mathrm{eff}}4`$. For weaker interactions the calculated width distributions exhibit some deviations from (6). We note however that, for the small lattices used, it is difficult to find a disorder strength for which the model is in the metallic diffusive regime and displays universal RMT statistics. In particular, in the presence of a magnetic flux we could not find values of $`W`$ for which both the spectral and wavefunction statistics are GUE. In conclusion, we have investigated the width and peak height statistics in Coulomb-blockade quantum dots using a random interaction matrix model with interactions that preserve time-reversal symmetry. For a GOE one-body symmetry the statistics is insensitive to the interaction strength. However, at strong interactions, a time-reversal symmetry-breaking field leads only to a partial crossover from GOE to GUE statistics. Our random interaction matrix model can reproduce both the observed Gaussian-like shape of the peak spacing distribution and the RMT statistics of the peak height distributions. We thank Y. Gefen for useful discussions. This work was supported in part by the U.S. DOE grant No. DE-FG-0291-ER-40608. A.W. acknowledges a fellowship from the Studienstiftung des Deutschen Volkes.
warning/0003/hep-ph0003100.html	ar5iv	text	# 1 Introduction ## 1 Introduction Hard processes in massless gauge field theories as QED and QCD exhibit Mellin–structures. The observables $`\sigma `$ can be written as Mellin convolutions $`\sigma (x)=\left[AB\right](x)={\displaystyle _0^1}𝑑x_1{\displaystyle _0^1}𝑑x_2\delta (xx_1x_2)A_1(x_1)A_2(x_2).`$ (1) Here, the functions $`A_i(x_i)`$ denote hard cross sections and splitting functions, respectively. In practice even multiple Mellin convolutions occur in the calculations of higher order Feynman diagrams. Integrals of this type can be diagonalized by the Mellin transform $`\text{M}\left[\sigma (x)\right](N)={\displaystyle _0^1}𝑑xx^N\sigma (x)=\text{M}\left[A_1(x_1)\right](N)\text{M}\left[A_2(x_2)\right](N)`$ (2) and $`\sigma (x)`$ can be found by a single inverse Mellin transform after evaluating the functions $`\text{M}\left[A_i(x_i)\right](N)`$, which are simpler in general. For $`Nϵ\text{N}^+`$ the Mellin transforms of the functions $`A_i(x)`$ can be represented in terms of linear combinations of finite harmonic sums \[1–3\] and their polynomials. $$S_{k_1,\mathrm{},k_m}(N)=\underset{n_1=1}{\overset{N}{}}\frac{\left[\mathrm{sign}(k_1)\right]^{n_1}}{n_1^{\|k_1\|}}\underset{n_2=1}{\overset{n_1}{}}\frac{\left[\mathrm{sign}(k_2)\right]^{n_2}}{n_2^{\|k_2\|}}\mathrm{}\underset{n_m=1}{\overset{n_{m1}}{}}\frac{\left[\mathrm{sign}(k_m)\right]^{n_m}}{n_m^{\|k_m\|}}.$$ (3) Up to the level of two–loop order only harmonic sums up to rank 4 contribute. The set of functions $`A_i(x)`$ is formed by polynomials of Nielsen integrals $$S_{n,p}(x)=\frac{(1)^{n+p1}}{(n1)!p!}_0^1\frac{dz}{z}\mathrm{log}^{n1}(z)\mathrm{log}^p(1zx).$$ (4) As was shown in Ref. the representation of the Mellin transforms in terms of finite harmonic sums may be reduced significantly by applying algebraic relations, cf. also , between these sums which are implied by index permutation and decomposition. Up to the level of two–loop order only 25 basic functions remain out of which the coefficient functions and splitting functions can be assembled as Mellin polynomials. Any of the multiple Mellin convolutions discussed above can thus be traced back to polynomials of the basic functions and a single numerical Mellin inversion which is performed as a complex contour integral. The Mellin transforms, which are firstly obtained at the positive integers, have to be analytically continued to the complex $`N`$–plane. This analytic continuation is unique . For single harmonic sums the continuation is well–known $`S_k(N)`$ $`=`$ $`(1)^{k1}{\displaystyle \frac{1}{(k1)!}}\psi ^{(k1)}(N+1)+c_k^+`$ (5) $`S_k(N)`$ $`=`$ $`(1)^{k1+N}{\displaystyle \frac{1}{(k1)!}}\beta ^{(k1)}(N+1)c_k^{},`$ (6) with $`\psi (z)`$ the logarithmic derivative of the Euler $`\mathrm{\Gamma }`$–function and $`\beta (z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\psi \left({\displaystyle \frac{z+1}{2}}\right)\psi \left({\displaystyle \frac{z}{2}}\right)\right]`$ (7) $`c_1^+`$ $`=`$ $`\gamma _E`$ (8) $`c_k^+`$ $`=`$ $`\zeta (k),k2`$ (9) $`c_1^{}`$ $`=`$ $`\mathrm{log}(2)`$ (10) $`c_k^+`$ $`=`$ $`\left(1{\displaystyle \frac{1}{2^{k1}}}\right)\zeta (k),k2.`$ (11) Here $`\gamma _E`$ denotes the Euler–Mascheroni constant and $`\zeta (k)`$ is the Riemann $`\zeta `$–function. It is the aim of the present paper to calculate the analytic continuations of the Mellin transforms of these 25 basic functions and to provide a code for numerical evaluations. The paper is organized as follows. The basic functions are introduced in section 2 and their structure is discussed. In section 3 the representation of the basic functions is given for positive integers. These representations which are given in terms of multiple harmonic sums do also provide first detailed numerical tests for the representations used for the analytic continuations later. The Mellin transforms for complex argument are given in section 4 for the functions of the type $`f_i(x)/(1+x)`$ and for the functions $`(f_k(x)f_k(1))/(x1)`$ in section 5. Here we aim on high numerical accuracy. Section 6 describes the basic options of the code ANCONT. The results are summarized in section 7. ## 2 Basic Functions Using the algebraic relations given in Ref. one may show that the linear representations up to threefold harmonic sums up to transcendentality four can be expressed by the single harmonic sums $`S_{\pm k}(N)`$ and the Mellin–transforms of the following 25 basic functions : $$\begin{array}{ccccc}g_1(x)\hfill & =& \frac{\mathrm{log}(1+x)}{x+1}\hfill & & S_{1,1}(N)\hfill \\ g_2(x)\hfill & =& \frac{\mathrm{log}^2(1+x)}{x+1}\hfill & & S_{1,1,1}(N)\hfill \\ g_3(x)\hfill & =& \frac{\text{Li}_2(x)}{x+1}\hfill & & S_{2,1}(N)\hfill \\ g_4(x)\hfill & =& \frac{\text{Li}_2(x)}{x+1}\hfill & & S_{2,1}(N)\hfill \\ g_5(x)\hfill & =& \frac{\mathrm{log}(x)\text{Li}_2(x)}{x+1}\hfill & & S_{2,2}(N),S_{3,1}(N)\hfill \\ g_6(x)\hfill & =& \frac{\text{Li}_3(x)}{x+1}\hfill & & S_{3,1}(N)\hfill \\ g_7(x)\hfill & =& \frac{\text{Li}_3(x)}{x+1}\hfill & & S_{3,1}(N)\hfill \\ g_8(x)\hfill & =& \frac{\text{S}_{1,2}(x)}{x+1}\hfill & & S_{2,1,1}(N)\hfill \\ g_9(x)\hfill & =& \frac{\text{S}_{1,2}(x)}{x+1}\hfill & & S_{2,1,1}(N)\hfill \\ g_{10}(x)\hfill & =& \frac{I_1(x)}{1+x}\hfill & & S_{1,2,1}(N),S_{2,1,1}(N)\hfill \\ g_{11}(x)\hfill & =& \frac{\mathrm{log}(1x)\text{Li}_2(x)}{1+x}\hfill & & S_{1,2,1}(N),S_{2,1,1}\hfill \\ g_{12}(x)\hfill & =& \mathrm{log}(1x)\frac{\text{Li}_2(x)}{x+1}\hfill & & S_{1,2,1}(N),S_{2,1,1}(N),S_{2,1,1}(N)\hfill \\ g_{13}(x)\hfill & =& \frac{\mathrm{log}(1+x)}{1+x}\text{Li}_2(x)\hfill & & S_{1,2,1}(N),S_{2,1,1}(N)\hfill \\ g_{14}(x)\hfill & =& \frac{\mathrm{log}^2(1+x)\mathrm{log}^2(2)}{x1}\hfill & & S_{1,1,1}(N)\hfill \end{array}$$ (12) $`\begin{array}{ccccc}g_{15}(x)\hfill & =& {\displaystyle \frac{\mathrm{log}(1+x)\mathrm{log}(2)}{x1}}\text{Li}_2(x)\hfill & & S_{1,2,1}(N),S_{2,1,1}(N),S_{2,1,1}(N)\hfill \\ g_{16}(x)\hfill & =& {\displaystyle \frac{\mathrm{log}(1+x)\mathrm{log}(2)}{x1}}\text{Li}_2(x)\hfill & & S_{1,2,1}(N),S_{2,1,1}(N)\hfill \\ g_{17}(x)\hfill & =& {\displaystyle \frac{\mathrm{log}(x)\mathrm{log}^2(1+x)}{x1}}\hfill & & S_{1,1,2}(N),S_{1,2,1}(N),S_{2,1,1}(N)\hfill \\ g_{18}(x)\hfill & =& {\displaystyle \frac{\text{Li}_2(x)\zeta (2)}{x1}}\hfill & & S_{2,1}(N)\hfill \\ g_{19}(x)\hfill & =& {\displaystyle \frac{\text{Li}_2(x)+\zeta (2)/2}{x1}}\hfill & & S_{2,1}(N)\hfill \\ g_{20}(x)\hfill & =& {\displaystyle \frac{\text{Li}_3(x)\zeta (3)}{x1}}\hfill & & S_{3,1}(N)\hfill \\ g_{21}(x)\hfill & =& {\displaystyle \frac{\text{S}_{1,2}(x)\zeta (3)}{x1}}\hfill & & S_{2,1,1}(N)\hfill \\ g_{22}(x)\hfill & =& {\displaystyle \frac{\mathrm{log}(x)\text{Li}_2(x)}{x1}}\hfill & & S_{3,1}(N)\hfill \\ g_{23}(x)\hfill & =& {\displaystyle \frac{\text{Li}_3(x)+3\zeta (3)/4}{x1}}\hfill & & S_{3,1}(N)\hfill \\ g_{24}(x)\hfill & =& {\displaystyle \frac{I_1(x)+(5/8)\zeta (3)}{x1}}\hfill & & S_{2,1,1}(N),S_{2,1,1}(N)\hfill \\ g_{25}(x)\hfill & =& {\displaystyle \frac{\text{S}_{1,2}(x)\zeta (3)/8}{x1}}\hfill & & S_{2,1,1}(N),\hfill \end{array}`$ (24) with $`I_1(x)={\displaystyle _0^x}{\displaystyle \frac{dz}{z}}\mathrm{log}(1z)\mathrm{log}(1+z).`$ (25) Here we listed also all harmonic sums of rank 4 which contribute to the respective Mellin transforms. The functions $`g_i(x)`$ belong to the class of Nielsen–integrals $`\mathrm{S}_{n,p}(x)`$, Eq. (4). The degree of transcendentality of these functions is defined to be $`\tau =p+n`$. The transcendentality of the product of two of these these functions is the sum of their transcendentalities. The measures $`dx/(x\pm 1)`$ are defined to be of transcendentality 1. The polylogarithms are given by $`\text{Li}_n(x)={\displaystyle \frac{d\text{Li}_{n+1}(x)}{d\mathrm{log}(x)}}\text{S}_{n1,1}(x)={\displaystyle \frac{(1)^{n1}}{(n2)!}}{\displaystyle _0^1}{\displaystyle \frac{dz}{z}}\mathrm{log}^{n2}(z)\mathrm{log}(1zx)\mathrm{for}n2.`$ (26) The logarithms $`\mathrm{log}(1\pm x)`$ are related to the dilogarithm by $`{\displaystyle \frac{d\text{Li}_2(\pm x)}{d\mathrm{log}(x)}}=\text{Li}_1(\pm x)=\mathrm{log}(1x)`$ (27) and $`\text{Li}_0(x)={\displaystyle \frac{x}{1x}}.`$ (28) Similarly, $`{\displaystyle \frac{d\text{S}_{n,p}(x)}{d\mathrm{log}(x)}}=\text{S}_{n1,p}(x)`$ (29) holds. As a generalization of the Nielsen–integrals (4) one may wish to consider $$\mathrm{S}_{n,p,q}(x)=\frac{(1)^{n+p+q1}}{(n1)!p!q!}_0^1\frac{dz}{z}\mathrm{log}^{n1}(z)\mathrm{log}^p(1zx)\mathrm{log}^q(1+zx).$$ (30) Both the classes of functions (4) and (30) are non–minimal and large subsets of them can be represented already by the polylogarithms $`\text{Li}_l(y)`$, where $`y`$ denotes a function of $`x`$. Relations between the different Nielsen–integrals and their generalizations are, however, more easily established using non–minimal representations, since the argument structure turns out to be simple in the latter case. Following this line we have introduced in the above set of functions $`I_1(x)S_{1,1,1}(x)={\displaystyle _0^x}{\displaystyle \frac{dz}{z}}\mathrm{log}(1z)\mathrm{log}(1+z)={\displaystyle \frac{1}{2}}\text{S}_{1,2}(x^2)\text{S}_{1,2}(x)\text{S}_{1,2}(x),`$ (31) instead referring to $`\text{S}_{1,2}(x^2)`$. The Mellin transforms of the functions (12) are related to a series of finite harmonic sums through which all harmonic sums up to level three and transcendentality four may be expressed in terms of polynomials, except of the well–known case of single harmonic sums and polynomials which are made of only single harmonic sums. To establish the connection more closely we mentioned above those sums in which the respective Mellin transforms occur in explicit form, see also Ref. . ## 3 Representations of the Mellin Transforms at Positive Integers To compare the numerical expressions of the analytic continuations of the Mellin–transforms given in the forthcoming sections we summarize the exact representations of the Mellin–transforms for $`Nϵ𝐍^+`$ for the basic functions in terms of harmonic sums. These relations are used as one test for the numerical accuracy of the analytic continuations. $`\text{M}\left[g_1(x)\right](N)`$ $`=`$ $`(1)^N\{S_{1,1}(N)+{\displaystyle \frac{1}{2}}\mathrm{log}^2(2)[S_1(N)S_1(N)]\mathrm{log}(2)`$ (32) $`S_1(N)S_1(N)S_2(N)\}`$ $`\text{M}\left[g_2(x)\right](N)`$ $`=`$ $`2(1)^N\{S_{1,1,1}(N)\mathrm{log}(2)[S_{1,1}(N)S_{1,1}(N)]`$ (33) $`{\displaystyle \frac{1}{2}}\mathrm{log}^2(2)[S_1(N)S_1(N)]+{\displaystyle \frac{1}{6}}\mathrm{log}^3(2)\}`$ $`\text{M}\left[g_3(x)\right](N)`$ $`=`$ $`(1)^{N+1}\left[S_{2,1}(N)\zeta (2)S_1(N)+{\displaystyle \frac{5}{8}}\zeta (3)\zeta (2)\mathrm{log}(2)\right]`$ (34) $`\text{M}\left[g_4(x)\right](N)`$ $`=`$ $`(1)^{N+1}\{S_{2,1}(N)+\mathrm{log}(2)[S_2(N)S_2(N)]+{\displaystyle \frac{1}{2}}\zeta (2)S_1(N)`$ (35) $`{\displaystyle \frac{1}{4}}\zeta (3)+{\displaystyle \frac{1}{2}}\zeta (2)\mathrm{log}(2)\}`$ $`\text{M}\left[g_5(x)\right](N)`$ $`=`$ $`(1)^N\left[S_{2,2}(N)+2S_{3,1}(N)2\zeta (2)S_2(N){\displaystyle \frac{3}{40}}\zeta (2)^2\right]`$ (36) $`\text{M}\left[g_6(x)\right](N)`$ $`=`$ $`(1)^N[S_{3,1}(N)\zeta (2)S_2(N)+\zeta (3)S_1(N)`$ (37) $`+{\displaystyle \frac{3}{5}}\zeta (2)^22\text{Li}_4\left({\displaystyle \frac{1}{2}}\right){\displaystyle \frac{3}{4}}\zeta (3)\mathrm{log}(2)+{\displaystyle \frac{1}{2}}\zeta (2)\mathrm{log}^2(2){\displaystyle \frac{1}{12}}\mathrm{log}^4(2)]`$ $`\text{M}\left[g_7(x)\right](N)`$ $`=`$ $`(1)^N\{S_{3,1}(N)+\mathrm{log}(2)[S_3(N)S_3(N)]+{\displaystyle \frac{1}{2}}\zeta (2)S_2(N)`$ (38) $`{\displaystyle \frac{3}{4}}\zeta (3)S_1(N)+{\displaystyle \frac{1}{8}}\zeta (2)^2{\displaystyle \frac{3}{4}}\zeta (3)\mathrm{log}(2)\}`$ $`\text{M}\left[g_8(x)\right](N)`$ $`=`$ $`(1)^{N+1}[S_{2,1,1}(N)\zeta (3)S_1(N)+\text{Li}_4\left({\displaystyle \frac{1}{2}}\right){\displaystyle \frac{1}{8}}\zeta (2)^2{\displaystyle \frac{1}{8}}\zeta (3)\mathrm{log}(2)`$ (39) $`{\displaystyle \frac{1}{4}}\zeta (2)\mathrm{log}^2(2)+{\displaystyle \frac{1}{24}}\mathrm{log}^4(2)]`$ $`\text{M}\left[g_9(x)\right](N)`$ $`=`$ $`(1)^{N+1}\{S_{2,1,1}(N)+\mathrm{log}(2)[S_{2,1}(N)S_{2,1}(N)]`$ (40) $`{\displaystyle \frac{1}{2}}\mathrm{log}^2(2)\left[S_2(N)S_2(N)\right]{\displaystyle \frac{1}{8}}\zeta (3)S_1(N)`$ $`3\text{Li}_4\left({\displaystyle \frac{1}{2}}\right)+{\displaystyle \frac{6}{5}}\zeta (2)^2{\displaystyle \frac{11}{4}}\zeta (3)\mathrm{log}(2)+{\displaystyle \frac{3}{4}}\zeta (2)\mathrm{log}^2(2){\displaystyle \frac{1}{8}}\mathrm{log}^4(2)\}`$ $`\text{M}\left[g_{10}(x)\right](N)`$ $`=`$ $`(1)^{N+1}\{S_{2,1,1}(N)+S_{2,1,1}(N)\mathrm{log}(2)[S_{2,1}(N)S_{2,1}(N)]`$ (41) $`+{\displaystyle \frac{1}{2}}\left[\zeta (2)\mathrm{log}^2(2)\right]\left[S_2(N)S_2(N)\right]+{\displaystyle \frac{5}{8}}\zeta (3)\left[S_1(N)+\mathrm{log}(2)\right]`$ $`{\displaystyle \frac{3}{20}}\zeta (2)^2\}`$ $`\text{M}\left[g_{11}(x)\right](N)`$ $`=`$ $`(1)^N\{S_{1,2,1}(N)+2S_{2,1,1}(N)2\zeta (3)S_1(N)\zeta (2)S_{1,1}(N)`$ (42) $`{\displaystyle \frac{29}{40}}\zeta (2)^2+3\text{Li}_4\left({\displaystyle \frac{1}{2}}\right){\displaystyle \frac{1}{4}}\zeta (2)\mathrm{log}^2(2)+{\displaystyle \frac{1}{4}}\mathrm{log}^4(2)\}`$ $`\text{M}\left[g_{12}(x)\right](N)`$ $`=`$ $`(1)^N\{S_{2,1,1}(N)+S_{2,1,1}(N)+S_{1,2,1}(N)`$ (43) $`\mathrm{log}(2)\left[S_{2,1}(N)+S_{1,2}(N)\right]+{\displaystyle \frac{1}{2}}\zeta (2)S_{1,1}(N)`$ $`+\left[S_1(N)S_2(N)+S_3(N)\right]\mathrm{log}(2)`$ $`+{\displaystyle \frac{1}{2}}\left[\zeta (2)\mathrm{log}^2(2)\right]\left[S_2(N)S_2(N)\right]+{\displaystyle \frac{5}{8}}\zeta (3)S_1(N)`$ $`4\text{Li}_4\left({\displaystyle \frac{1}{2}}\right)+{\displaystyle \frac{3}{2}}\zeta )2)^2{\displaystyle \frac{21}{8}}\zeta (3)\mathrm{log}(2)+{\displaystyle \frac{3}{4}}\zeta (2)\mathrm{log}^2(2){\displaystyle \frac{1}{6}}\mathrm{log}^4(2)\}`$ $`\text{M}\left[g_{13}(x)\right](N)`$ $`=`$ $`(1)^N\{S_{1,2,1}(N)+2S_{2,1,1}(N)+{\displaystyle \frac{1}{2}}\zeta (2)S_{1,1}(N)`$ (44) $`\mathrm{log}^2(2)\left[S_2(N)S_2(N)\right]`$ $`+\mathrm{log}(2)[S_{2,1}(N)S_{1,2}(N)2S_{2,1}(N)+S_1(N)S_2(N)`$ $`+S_3(N){\displaystyle \frac{1}{2}}\zeta (2)S_1(N)]`$ $`3\text{Li}_4\left({\displaystyle \frac{1}{2}}\right)+{\displaystyle \frac{6}{5}}\zeta (2)^2{\displaystyle \frac{21}{8}}\zeta (3)\mathrm{log}(2)+{\displaystyle \frac{1}{2}}\zeta (2)\mathrm{log}^2(2){\displaystyle \frac{1}{8}}\mathrm{log}^4(2)\}`$ $`\text{M}\left[g_{14}(x)\right](N)`$ $`=`$ $`2S_{1,1,1}(N)2\mathrm{log}(2)\left[S_{1,1}(N)S_{1,1}(N)\right]\mathrm{log}^2(2)S_1(N)`$ (45) $`{\displaystyle \frac{1}{4}}\zeta (3)+\zeta (2)\mathrm{log}(2){\displaystyle \frac{2}{3}}\mathrm{log}^3(2)`$ $`\text{M}\left[g_{15}(x)\right](N)`$ $`=`$ $`S_{1,2,1}(N)+S_{2,1,1}(N)+S_{2,1,1}(N)\zeta (2)S_{1,1}(N)+{\displaystyle \frac{5}{18}}\zeta (3)S_1(N)`$ (46) $`\left[{\displaystyle \frac{5}{8}}\zeta (3)\zeta (2)\mathrm{log}(2)\right]\left[S_1(N)S_1(N)\right]+\mathrm{log}(2)\left[S_{2,1}(N)2\zeta (3)\right]`$ $`{\displaystyle \frac{1}{2}}\left(\zeta (2)\mathrm{log}^2(2)\right)\left[S_2(N)S_2(N)\right]\mathrm{log}(2)\zeta (2)S_1(N)`$ $`+{\displaystyle \frac{19}{4}}\zeta (2)^2\text{Li}_4\left({\displaystyle \frac{1}{2}}\right)+{\displaystyle \frac{7}{4}}\zeta (3)\mathrm{log}(2){\displaystyle \frac{1}{4}}\zeta (2)\mathrm{log}^2(2){\displaystyle \frac{1}{24}}\mathrm{log}^4(2)`$ $`\text{M}\left[g_{16}(x)\right](N)`$ $`=`$ $`2S_{2,1,1}(N)+S_{1,2,1}(N)`$ (47) $`+\mathrm{log}(2)\left[2\left(S_{2,1}(N)S_{2,1}(N)\right)+S_{1,2}(N)S_{1,2}(N)\right]`$ $`+{\displaystyle \frac{1}{2}}\zeta (2)S_{1,1}(N)\mathrm{log}^2(2)\left[S_2(N)S_2(N)\right]{\displaystyle \frac{1}{4}}\zeta (3)S_1(N)`$ $`\left({\displaystyle \frac{1}{4}}\zeta (3){\displaystyle \frac{1}{2}}\zeta (2)\mathrm{log}(2)\right)\left[S_1(N)S_1(N)\right]`$ $`+\mathrm{log}(2)\left[S_{2,1}(N)\mathrm{log}(2)\left(S_2(N)S_2(N)\right)+{\displaystyle \frac{1}{2}}\zeta (2)S_1(N)\right]`$ $`{\displaystyle \frac{33}{20}}\zeta (2)^2+4\text{Li}_4\left({\displaystyle \frac{1}{2}}\right)+{\displaystyle \frac{13}{4}}\zeta (3)\mathrm{log}(2){\displaystyle \frac{3}{4}}\zeta (2)\mathrm{log}^2(2)+{\displaystyle \frac{1}{6}}\mathrm{log}^4(2)`$ $`\text{M}\left[g_{17}(x)\right](N)`$ $`=`$ $`2\left[S_{1,2,1}(N)+S_{1,1,2}(N)+S_{2,1,1}(N)\right]`$ (48) $`+\mathrm{log}(2)\left[S_{1,2}(N)S_{1,2}(N)S_{2,1}(N)+S_{2,1}(N)\right]`$ $`+{\displaystyle \frac{1}{8}}\zeta (3)S_1(N)+{\displaystyle \frac{1}{2}}\mathrm{log}^2(2)\left[S_2(N)S_2(N)\right]{\displaystyle \frac{1}{4}}\zeta (2)S_{1,1}(N)`$ $`+{\displaystyle \frac{7}{4}}\zeta (2)^24\text{Li}_4\left({\displaystyle \frac{1}{2}}\right){\displaystyle \frac{21}{4}}\zeta (3)\mathrm{log}(2)+{\displaystyle \frac{5}{2}}\zeta (2)\mathrm{log}^2(2){\displaystyle \frac{1}{6}}\mathrm{log}^4(2)`$ $`\text{M}\left[g_{18}(x)\right](N)`$ $`=`$ $`S_{2,1}(N)+2\zeta (3)`$ (49) $`\text{M}\left[g_{19}(x)\right](N)`$ $`=`$ $`S_{2,1}(N)+\mathrm{log}(2)\left[S_2(N)S_2(N)\right]{\displaystyle \frac{5}{8}}\zeta (3)`$ (50) $`\text{M}\left[g_{20}(x)\right](N)`$ $`=`$ $`S_{3,1}(N)\zeta (2)S_2(N)+{\displaystyle \frac{1}{2}}\zeta (2)^2`$ (51) $`\text{M}\left[g_{21}(x)\right](N)`$ $`=`$ $`S_{2,1,1}(N)+{\displaystyle \frac{6}{5}}\zeta (2)^2`$ (52) $`\text{M}\left[g_{22}(x)\right](N)`$ $`=`$ $`2S_{3,1}(N)+{\displaystyle \frac{1}{2}}S_4(N)+{\displaystyle \frac{1}{2}}S_2^2(N)2\zeta (2)S_2(N)+{\displaystyle \frac{3}{10}}\zeta (2)^2`$ (53) $`\text{M}\left[g_{23}(x)\right](N)`$ $`=`$ $`S_{3,1}(N)\mathrm{log}(2)\left[S_3(N)S_3(N)\right]+{\displaystyle \frac{1}{2}}\zeta (2)S_2(N)`$ (54) $`+2\text{Li}_4\left({\displaystyle \frac{1}{2}}\right){\displaystyle \frac{11}{10}}\zeta (2)^2+{\displaystyle \frac{7}{4}}\zeta (3)\mathrm{log}(2){\displaystyle \frac{1}{2}}\zeta (2)\mathrm{log}^2(2)+{\displaystyle \frac{1}{12}}\mathrm{log}^4(2)`$ $`\text{M}\left[g_{24}(x)\right](N)`$ $`=`$ $`S_{2,1,1}(N)S_{2,1,1}(N)+\mathrm{log}(2)\left[S_{2,1}(N)S_{2,1}(N)\right]`$ (55) $`+{\displaystyle \frac{1}{2}}\left[\zeta (2)\mathrm{log}^2(2)\right]\left[S_2(N)S_2(N)\right]{\displaystyle \frac{5}{8}}\zeta (3)S_1(N)`$ $`+{\displaystyle \frac{1}{4}}\zeta (2)^22\text{Li}_4\left({\displaystyle \frac{1}{2}}\right){\displaystyle \frac{7}{4}}\zeta (3)\mathrm{log}(2)+{\displaystyle \frac{1}{2}}\zeta (2)\mathrm{log}^2(2){\displaystyle \frac{1}{12}}\mathrm{log}^4(2)`$ $`\text{M}\left[g_{25}(x)\right](N)`$ $`=`$ $`S_{2,1,1}(N)\mathrm{log}(2)\left[S_{2,1}(N)S_{2,1}(N)\right]`$ (56) $`+{\displaystyle \frac{1}{2}}\mathrm{log}^2(2)\left[S_2(N)S_2(N)\right]+{\displaystyle \frac{3}{40}}\zeta (2)^2`$ The above relations can be checked by numerical integration directly. In the code ANCONT we use the routines DAIND, Ref. , and verified these relations numerically for the integer moments $`Nϵ[1,20]`$ at an accuracy of better than at least $`410^9`$, and in many cases even of $`O(10^{12})`$. ## 4 Mellin Transforms for the Functions $`𝐟(𝐱)/(𝐱+\mathrm{𝟏})`$ For functions of the type $$g_i(x)=\frac{f_i(x)}{x+1}$$ (57) and $`f_i(x)ϵ𝒞^{\mathrm{}}[0,1[`$, cf. , one may expand $`f_i(x)`$ into a Taylor series around $`x=0`$, $$f_i(x)=\underset{k=0}{\overset{\mathrm{}}{}}a_kx^k.$$ (58) The Mellin–transform is then given by $$\text{M}\left[\frac{f_i(x)}{x+1}\right](N)=\underset{k=0}{\overset{\mathrm{}}{}}a_k\beta (N+k+1).$$ (59) Since the Mellin–transforms of the type $`\text{M}[\mathrm{log}^m(x)f(x)](N)`$ can be calculated using the relation $`\text{M}\left[\mathrm{log}^m(x)f(x)\right](N)={\displaystyle \frac{^N}{N^k}}\text{M}\left[f(x)\right](N),`$ (60) one obtains $$\text{M}\left[\mathrm{log}^m(x)\frac{f_i(x)}{1+x}\right](N)=\underset{k=0}{\overset{\mathrm{}}{}}a_k\beta ^{(m)}(N+k+1).$$ (61) The function $`\beta ^{(m)}(z)`$ obeys the recursion relation $$\beta ^{(m)}(z+1)=(1)^mm!\left[\frac{1}{(1+z)^{m+1}}\frac{1}{z^{m+1}}\right]+\beta ^{(m)}(z).$$ (62) For the polylogarithms $`\text{Li}_l(\pm x)`$ one obtains $$\text{M}\left[\mathrm{log}^m(x)\frac{\text{Li}_l(\pm x)}{x+1}\right](N)=\underset{k=1}{\overset{\mathrm{}}{}}\frac{(\pm 1)^k}{k^l}\beta ^{(m)}(N+k+1),$$ (63) as an example. Although these representations are quite general and may be applied to other Nielsen–integrals as well, the corresponding series may not converge fast enough. Alternatively, the analytic continuation of the Mellin transforms containing the factor $`(x+1)^1`$ can be performed using the transformation $`\text{M}\left[{\displaystyle \frac{f(x)}{1+x}}\right](N){\displaystyle _0^1}𝑑x{\displaystyle \frac{x^N}{x+1}}f(x)`$ $`=`$ $`\mathrm{log}(2)f(1){\displaystyle _0^1}𝑑xx^{N1}\mathrm{log}(1+x)\left[Nf(x)+xf^{}(x)\right]`$ (64) $`=`$ $`\mathrm{log}(2)f(1){\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}\{N\text{M}[f(x)](N+k1)`$ $`+\text{M}[f^{}(x)](N+k)\},`$ if $`f(x)`$ is differentiable in $`]0,1[`$. The function $`\mathrm{log}(1+x)`$ can be approximated by $$\mathrm{ln}(1+x)\underset{k=1}{\overset{9}{}}a_k^{(1)}x^k,$$ (65) with an accuracy better than $`3\times 10^8`$. The polynomial is determined using the approximation given by the MINIMAX–method with an adaptive choice of arguments $`x`$. The coefficients $`a_k^{(1)}`$ are given by $`\begin{array}{cccccc}a_1^{\left(1\right)}\hfill & =& \hfill 0.999999974532238\mathrm{E}+0& a_2^{\left(1\right)}\hfill & =& \hfill 0.499995525889840\mathrm{E}+0\\ a_3^{\left(1\right)}\hfill & =& \hfill 0.333203435557262\mathrm{E}+0& a_4^{\left(1\right)}\hfill & =& \hfill 0.248529457782640\mathrm{E}+0\\ a_5^{\left(1\right)}\hfill & =& \hfill 0.191451164719161\mathrm{E}+0& a_6^{\left(1\right)}\hfill & =& \hfill 0.137466222728331\mathrm{E}+0\\ a_7^{\left(1\right)}\hfill & =& \hfill 0.792107412244877\mathrm{E}1& a_8^{\left(1\right)}\hfill & =& \hfill 0.301109656912626\mathrm{E}1\\ a_9^{\left(1\right)}\hfill & =& \hfill 0.538406208663153\mathrm{E}2& & & \end{array}`$ (67) Similar representations are found in the literature cf. \[11–13\] and have been used in Ref. before. Up to two–loop order also the representations of $`\mathrm{log}^k(1+x)`$ are needed with $`k=2,3`$, $`\mathrm{ln}^2(1+x)`$ $``$ $`{\displaystyle \underset{k=2}{\overset{11}{}}}a_k^{(2)}x^k,`$ (68) $`\mathrm{ln}^3(1+x)`$ $``$ $`{\displaystyle \underset{k=3}{\overset{13}{}}}a_k^{(3)}x^k.`$ (69) The coefficients are $`\begin{array}{cccccc}a_2^{\left(2\right)}\hfill & =& \hfill 0.999999980543793\mathrm{E}+0& a_3^{\left(2\right)}\hfill & =& \hfill 0.999995797779624\mathrm{E}+0\\ a_4^{\left(2\right)}\hfill & =& \hfill 0.916516447393493\mathrm{E}+0& a_5^{\left(2\right)}\hfill & =& \hfill 0.831229921350708\mathrm{E}+0\\ a_6^{\left(2\right)}\hfill & =& \hfill 0.745873737923571\mathrm{E}+0& a_7^{\left(2\right)}\hfill & =& \hfill 0.634523908078600\mathrm{E}+0\\ a_8^{\left(2\right)}\hfill & =& \hfill 0.467104011423750\mathrm{E}+0& a_9^{\left(2\right)}\hfill & =& \hfill 0.261348046799178\mathrm{E}+0\\ a_{10}^{\left(2\right)}\hfill & =& \hfill 0.936814286867420\mathrm{E}1& a_{11}^{\left(2\right)}\hfill & =& \hfill 0.156249375012462\mathrm{E}1\end{array}`$ (71) and $`\begin{array}{cccccc}a_3^{\left(3\right)}\hfill & =& \hfill 0.999999989322696\mathrm{E}+0& a_4^{\left(3\right)}\hfill & =& \hfill 0.149999722020708\mathrm{E}+1\\ a_5^{\left(3\right)}\hfill & =& \hfill 0.174988008499745\mathrm{E}+1& a_6^{\left(3\right)}\hfill & =& \hfill 0.187296689068405\mathrm{E}+1\\ a_7^{\left(3\right)}\hfill & =& \hfill 0.191539974617231\mathrm{E}+1& a_8^{\left(3\right)}\hfill & =& \hfill 0.185963744001295\mathrm{E}+1\\ a_9^{\left(3\right)}\hfill & =& \hfill 0.162987195424434\mathrm{E}+1& a_{10}^{\left(3\right)}\hfill & =& \hfill 0.117982353224299\mathrm{E}+1\\ a_{11}^{\left(3\right)}\hfill & =& \hfill 0.628710122994999\mathrm{E}+0& a_{12}^{\left(3\right)}\hfill & =& \hfill 0.211307487211713\mathrm{E}+0\\ a_{13}^{\left(3\right)}\hfill & =& \hfill 0.328953352932140\mathrm{E}1& & & \end{array}`$ (73) Eqs. (68,69) hold at an accuracy better than $`2\times 10^8`$ for $`xϵ[0,1]`$. The Mellin–transforms of the functions $`\mathrm{log}(1+x)/(1+x)`$ and $`\mathrm{log}^2(1+x)/(1+x)`$ are obtained by $`\text{M}\left[{\displaystyle \frac{\mathrm{log}^k(1+x)}{1+x}}\right](N)`$ $`=`$ $`{\displaystyle \frac{1}{k+1}}\left[\mathrm{log}^{k+1}(2)N\text{M}\left[\mathrm{log}^{k+1}(1+x)\right](N1)\right]`$ (74) $`=`$ $`{\displaystyle \frac{1}{k+1}}\left[\mathrm{log}^{k+1}(2)N{\displaystyle \underset{l=k+1}{\overset{L_1(k+1)}{}}}{\displaystyle \frac{a_l^{(k+1)}}{N+l}}\right].`$ with $`L_1(2)=11`$ and $`L_1(3)=13`$. Similarly the Mellin–transform of the function $`g_{17}(x)`$ is obtained by $$\text{M}\left[\frac{\mathrm{log}(x)}{x1}\mathrm{log}^2(1+x)\right](N)=\underset{k=2}{\overset{11}{}}a_k^{(2)}\psi ^{}(N+k+1).$$ (75) The remaining Mellin–transforms are obtained from the representation Eq. (64) and are expressed in terms of the known Mellin–transforms for the functions $`f(x)`$ and $`f^{}(x)`$, cf. : $`\text{M}\left[\text{Li}_2(x)\right](N)`$ $`=`$ $`{\displaystyle \frac{1}{N+1}}\left[\zeta (2){\displaystyle \frac{S_1(N+1)}{N+1}}\right]`$ (76) $`\text{M}\left[\text{Li}_2^{}(x)\right](N)`$ $`=`$ $`\text{M}\left[\mathrm{log}(1x)\right](N1)={\displaystyle \frac{S_1(N)}{N}}`$ (77) etc. The single harmonic sums $`S_{\pm k}(N)`$ are given in Eqs. (5,6). One obtains the representations : $`\text{M}\left[{\displaystyle \frac{\text{Li}_2(x)}{1+x}}\right](N)`$ $`=`$ $`\zeta (2)\mathrm{log}(2){\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}\left[{\displaystyle \frac{N}{N+k}}\zeta (2)+{\displaystyle \frac{k}{(N+k)^2}}S_1(N+k)\right]`$ (78) $`\text{M}\left[{\displaystyle \frac{\text{Li}_2(x)}{1+x}}\right](N)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\zeta (2)\mathrm{log}(2)+{\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}\{{\displaystyle \frac{N}{N+k}}{\displaystyle \frac{\zeta (2)}{2}}+{\displaystyle \frac{k}{(N+k)^2}}[\mathrm{log}(2)`$ (79) $`\beta (N+k+1)]\}`$ $`\text{M}\left[\mathrm{log}(x){\displaystyle \frac{\text{Li}_2(x)}{1+x}}\right](N)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{9}{}}}{\displaystyle \frac{k}{(N+k)^2}}\left[\zeta (2)+\psi ^{}(k+1)2{\displaystyle \frac{S_1(N+k)}{N+k}}\right]`$ (80) $`\text{M}\left[{\displaystyle \frac{\text{Li}_3(x)}{1+x}}\right](N)`$ $`=`$ $`\zeta (3)\mathrm{log}(2)`$ (81) $`{\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}\left\{{\displaystyle \frac{N}{N+k}}\zeta (3)+{\displaystyle \frac{k}{(N+k)^2}}\left[\zeta (2){\displaystyle \frac{S_1(N+k)}{N+k}}\right]\right\}`$ $`\text{M}\left[{\displaystyle \frac{\text{Li}_3(x)}{1+x}}\right](N)`$ $`=`$ $`{\displaystyle \frac{3}{4}}\zeta (3)\mathrm{log}(2)+{\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}\{{\displaystyle \frac{N}{N+k}}{\displaystyle \frac{3}{4}}\zeta (3)+{\displaystyle \frac{k}{(N+k)^2}}{\displaystyle \frac{1}{2}}\zeta (2)`$ (82) $`{\displaystyle \frac{k}{(N+k)^3}}[\mathrm{log}(2)\beta (N+k+1)]\}`$ $`\text{M}\left[{\displaystyle \frac{\text{S}_{1,2}(x)}{1+x}}\right](N)`$ $`=`$ $`\mathrm{log}(2)\zeta (3){\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}\{{\displaystyle \frac{N}{N+k}}\zeta (3)`$ (83) $`+{\displaystyle \frac{k}{(N+k)^2}}{\displaystyle \frac{1}{2}}[S_1^2(N+k)+S_2(N+k)]\}`$ $`\text{M}\left[{\displaystyle \frac{\text{S}_{1,2}(x)}{1+x}}\right](N)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\zeta (3)\mathrm{log}(2){\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}{\displaystyle \frac{N}{N+k}}\left[{\displaystyle \frac{\zeta (3)}{8}}{\displaystyle \frac{1}{2}}{\displaystyle \underset{l=2}{\overset{11}{}}}{\displaystyle \frac{a_l^{(2)}}{N+k+l}}\right]`$ (84) $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=3}{\overset{13}{}}}{\displaystyle \frac{a_k^{(3)}}{N+k}}`$ The function $`\text{S}_{1,2}(x^2)`$ does not occur singly but rather in the combination $`\text{S}_{1,2}(x^2)/2\text{S}_{1,2}(x)\text{S}_{1,2}(x)`$ in the Mellin–transforms which are related to the harmonic sums up to the level being investigated in the present paper. It emerges via the integral $$I_1(x)=_0^x\frac{dz}{z}\mathrm{log}(1z)\mathrm{log}(1+z)=\frac{1}{2}\text{S}_{1,2}(x^2)\text{S}_{1,2}(x)\text{S}_{1,2}(x).$$ (85) We firstly calculate the Mellin–transform for this integral. The Mellin–transform for $`\text{S}_{1,2}(x^2)`$ is then easily obtained by a linear combination using the relations given above. $`\text{M}\left[{\displaystyle \frac{I_1(x)}{1+x}}\right](N)`$ $`=`$ $`{\displaystyle \frac{5}{8}}\zeta (3)\mathrm{log}(2)+{\displaystyle \underset{k=2}{\overset{11}{}}}a_k^{(2)}{\displaystyle \frac{S_1(N+k)}{N+k}}`$ (86) $`+{\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}{\displaystyle \frac{N}{N+k}}\left[{\displaystyle \frac{5}{8}}\zeta (3){\displaystyle \underset{l=1}{\overset{9}{}}}a_l^{(1)}{\displaystyle \frac{S_1(N+k+l)}{N+k+l}}\right].`$ \| $`N`$ \| Eq. (87) \| Eq. (88) \| Eq. (78) \| \| --- \| --- \| --- \| --- \| \| 0 \| -0.717E-03 \| 0.969E-04 \| 0.212E-08 \| \| 1 \| -0.117E-02 \| 0.205E-03 \| 0.442E-08 \| \| 2 \| -0.160E-02 \| 0.292E-03 \| 0.761E-08 \| \| 3 \| -0.207E-02 \| 0.378E-03 \| 0.117E-07 \| \| 4 \| -0.257E-02 \| 0.452E-03 \| 0.169E-07 \| \| 5 \| -0.308E-02 \| 0.505E-03 \| 0.232E-07 \| \| 6 \| -0.359E-02 \| 0.531E-03 \| 0.307E-07 \| \| 7 \| -0.410E-02 \| 0.529E-03 \| 0.397E-07 \| \| 8 \| -0.461E-02 \| 0.501E-03 \| 0.501E-07 \| \| 9 \| -0.511E-02 \| 0.449E-03 \| 0.622E-07 \| \| 10 \| -0.559E-02 \| 0.376E-03 \| 0.758E-07 \| \| 11 \| -0.607E-02 \| 0.286E-03 \| 0.909E-07 \| \| 12 \| -0.652E-02 \| 0.181E-03 \| 0.107E-06 \| \| 13 \| -0.697E-02 \| 0.637E-04 \| 0.125E-06 \| \| 14 \| -0.740E-02 \| -0.630E-04 \| 0.144E-06 \| \| 15 \| -0.782E-02 \| -0.197E-03 \| 0.164E-06 \| \| 16 \| -0.823E-02 \| -0.338E-03 \| 0.184E-06 \| \| 17 \| -0.862E-02 \| -0.483E-03 \| 0.205E-06 \| \| 18 \| -0.900E-02 \| -0.632E-03 \| 0.226E-06 \| \| 19 \| -0.936E-02 \| -0.784E-03 \| 0.247E-06 \| \| 20 \| -0.972E-02 \| -0.936E-03 \| 0.268E-06 \| Table 1: Comparison of the relative accuracy of different approximations for $`\text{M}\left[\text{Li}_2(x)/(1+x)\right](N)`$. Since the harmonic sum $`S_{2,1}(N)`$ contributes already to the next–to–leading order anomalous dimensions, approximations for $`\text{M}\left[\text{Li}_2(x)/(1+x)\right](N)`$ have been used in the literature for some time . Two approximations used before are $`\text{M}\left[{\displaystyle \frac{\text{Li}_2(x)}{1+x}}\right](N)`$ $``$ $`{\displaystyle \frac{1.01}{N+2}}{\displaystyle \frac{0.846}{N+3}}+{\displaystyle \frac{1.155}{N+4}}{\displaystyle \frac{1.074}{N+5}}+{\displaystyle \frac{0.55}{N+6}}`$ (87) or $`\text{M}\left[{\displaystyle \frac{\text{Li}_2(x)}{1+x}}\right](N)`$ $``$ $`{\displaystyle \frac{1.004}{N+2}}{\displaystyle \frac{0.846}{N+3}}+{\displaystyle \frac{1.342}{N+4}}{\displaystyle \frac{1.532}{N+5}}+{\displaystyle \frac{0.839}{N+6}},`$ (88) cf. , which can be compared to our representation Eq. (78). In Table 1 the integer moments of these representations are given for $`Nϵ[0,20]`$. Eq. (78) is more precise than Eqs. (87,88) by three to four orders of magnitude. Other quantities were approximated by similar representations as (87,88) in recently. A Mellin–transform of a function carrying the factor $`\mathrm{log}(1x)/(1+x)`$ can be performed if the remainder factor obeys a polynomial representation, since the former factor may be transformed into the Mellin–transform of $`\mathrm{log}(1+x)/(1+x)`$, cf. Ref. . Eqs. (44,45), by $$\text{M}\left[\frac{\mathrm{log}(1x)}{1+x}\right](N)=\text{M}\left[\frac{\mathrm{log}(1+x)}{1+x}\right](N)\beta (N+1)\left[\psi (N+1)+\gamma _E\mathrm{log}(2)\right]+\beta ^{}(N+1).$$ (89) Using Eq. (74) one obtains $`\text{M}\left[{\displaystyle \frac{\mathrm{log}(1x)\text{Li}_2(x)}{1+x}}\right](N)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=1}{\overset{9}{}}}{\displaystyle \frac{a_k^{(1)}}{k}}\{[\mathrm{log}^2(2){\displaystyle \underset{l=2}{\overset{11}{}}}a_l^{(2)}{\displaystyle \frac{N+k}{N+k+l}}]\beta ^{}(N+k+1).`$ (90) $`+\beta (N+k+1)[S_1(N+k)\mathrm{log}(2)]\}.`$ The Mellin–transform of the function $`\mathrm{log}(1x)\text{Li}_2(x)/(1+x)`$ can be calculated similarly to the case which was discussed before. $`\text{Li}_2(x)`$ may be represented by $$\text{Li}_2(x)P_0^{(1)}(x)+\left[P_1^{(1)}(x)+\mathrm{log}(1x)P_2^{(1)}(x)\right],$$ (91) where $`P_{0,1,2}^{(1)}(x)`$ are the polynomials $`P_0^{(1)}(x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{11}{}}}c_k^{(1)}x^k`$ (92) $`P_1^{(1)}(x)`$ $`=`$ $`{\displaystyle \frac{49}{36}}+\zeta (2)+{\displaystyle \frac{11}{6}}x{\displaystyle \frac{7}{12}}x^2+{\displaystyle \frac{1}{9}}x^3`$ (93) $`P_2^{(1)}(x)`$ $`=`$ $`{\displaystyle \frac{11}{6}}3x+{\displaystyle \frac{3}{2}}x^2{\displaystyle \frac{1}{3}}x^3.`$ (94) with $`\begin{array}{cccccc}c_0^{\left(1\right)}\hfill & =& \hfill 0.283822933724932\mathrm{E}+0& c_1^{\left(1\right)}\hfill & =& \hfill 0.999994319023731\mathrm{E}+0\\ c_2^{\left(1\right)}\hfill & =& \hfill 0.124975762907682\mathrm{E}+1& c_3^{\left(1\right)}\hfill & =& \hfill 0.607076808008983\mathrm{E}+0\\ c_4^{\left(1\right)}\hfill & =& \hfill 0.280403220046588\mathrm{E}1& c_5^{\left(1\right)}\hfill & =& \hfill 0.181869786537805\mathrm{E}+0\\ c_6^{\left(1\right)}\hfill & =& \hfill 0.532318519269331\mathrm{E}+0& c_7^{\left(1\right)}\hfill & =& \hfill 0.107281686995035\mathrm{E}+1\\ c_8^{\left(1\right)}\hfill & =& \hfill 0.138194913357518\mathrm{E}+1& c_9^{\left(1\right)}\hfill & =& \hfill 0.111100841298484\mathrm{E}+1\\ c_{10}^{\left(1\right)}\hfill & =& \hfill 0.506649587198046\mathrm{E}+0& c_{11}^{\left(1\right)}\hfill & =& \hfill 0.100672390783659\mathrm{E}+0\end{array}`$ The representation Eq. (91) holds at an accuracy of better than $`3\times 10^8`$ for $`xϵ[0,1]`$. The Mellin transform of $`\mathrm{log}(1x)\text{Li}_2(x)/(1+x)`$ is thus expressed as a sum over Mellin transforms $`\text{M}[\mathrm{log}(1x)/(1+x)]N)`$, Eq. (89) and $`\text{M}[\mathrm{log}^2(1x)/(1+x)](N)`$. The latter expression is related to the harmonic sum $`S_{1,1,1}(N)`$, Eq. (60). This sum can be expressed by the sum $`S_{1,1,1}(N)`$ and lower order harmonic sums through $`S_{1,1,1}(N)`$ $`=`$ $`S_{1,1,1}(N)+S_1(N)S_{1,1}(N)+S_{2,1}(N)+S_{1,2}(N)`$ (96) $`+{\displaystyle \frac{1}{2}}\left[S_1^2(N)S_1(N)S_1(N)S_2(N)\right]S_3(N).`$ The harmonic sum $`S_{1,1,1}(N)`$ is related to $`\text{M}\left[\mathrm{log}^2(1+x)/(1+x)\right](N)`$, which can be calculated using the representation (69) and lower order Mellin transforms. A simpler representation of $`\text{M}\left[\mathrm{log}^2(1x)/(1+x)\right](N)`$ is given by $`\text{M}\left[{\displaystyle \frac{\mathrm{log}^2(1x)}{1+x}}\right](N)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}\left\{{\displaystyle \frac{k}{N+k}}\left[S_1^2(N+k)+S_2(N+k)\right]\left[S_1^2(k)+S_2(k)\right]\right\}`$ (97) $`+{\displaystyle \frac{7}{4}}\zeta (3)\zeta (2)\mathrm{log}(2)+{\displaystyle \frac{1}{3}}\mathrm{log}^3(2)`$ referring to (65). Similarly $`\text{M}\left[{\displaystyle \frac{\mathrm{log}(1x)}{1+x}}\right](N)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}\left\{{\displaystyle \frac{k}{N+k}}S_1(N+K)S_1(k)\right\}+{\displaystyle \frac{1}{2}}\left[\mathrm{log}^2(2)\zeta (2)\right].`$ (98) holds. The above relations yield $`\text{M}\left[{\displaystyle \frac{\mathrm{log}(1x)\text{Li}_2(x)}{1+x}}\right](N)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{11}{}}}c_k^{(1)}\{{\displaystyle \frac{1}{2}}[\mathrm{log}^2(2)\zeta (2)]`$ (99) $`{\displaystyle \underset{l=1}{\overset{9}{}}}a_l^{(1)}[{\displaystyle \frac{l}{N+k+l}}S_1(N+k+l)S_1(l)]\}`$ $`+{\displaystyle \underset{k=0}{\overset{3}{}}}P_{2,k}^{(1)}\{{\displaystyle \frac{7}{4}}\zeta (3)\zeta (2)\mathrm{log}(2)+{\displaystyle \frac{1}{3}}\mathrm{log}^3(2)`$ $`+{\displaystyle \underset{l=1}{\overset{9}{}}}a_l^{(1)}[{\displaystyle \frac{l}{N+k+l}}[S_1^2(N+k+l)+S_2(N+k+l)]`$ $`S_1^2(l)S_2(l)]\}.`$ Finally, the Mellin–transform of $`\mathrm{log}(1+x)\text{Li}_2(x)/(1+x)`$ is given by $`\text{M}\left[{\displaystyle \frac{\mathrm{log}(1+x)\text{Li}_2(x)}{1+x}}\right](N)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\zeta (2)\mathrm{log}^2(2)+{\displaystyle \frac{1}{2}}\{{\displaystyle \underset{k=3}{\overset{13}{}}}{\displaystyle \frac{a_k^{(3)}}{N+k}}`$ (100) $`+{\displaystyle \underset{k=2}{\overset{11}{}}}a_k^{(2)}{\displaystyle \frac{N}{N+k}}[{\displaystyle \frac{1}{2}}\zeta (2){\displaystyle \frac{\mathrm{log}(2)\beta (N+k+1)}{N+k}}]\}`$ ## 5 Mellin Transforms for the Functions $`𝐟(𝐱)/(𝐱\mathrm{𝟏})`$ For this class of Mellin–transforms the role of $`\mathrm{log}(1+x)`$ in the previous section is taken by $`\mathrm{log}(1x)`$, provided the functions $`f(x)`$ in the numerator do vanish at a sufficient degree as $`x1`$. Unlike $`\mathrm{log}(1+x)`$, the function $`\mathrm{log}(1x)`$ has no simple polynomial representation in the range $`xϵ[0,1[`$ at a comparable accuracy to Eq. (65). One may be tempted to express the numerator functions in a series of $`\mathrm{log}(1x)`$ instead. This is indeed possible for a wide class of Nielsen integrals as $`\mathrm{S}_{1,p}(x)`$ in a simple manner, cf. also , $$\mathrm{S}_{1,p}(x)=\frac{1}{p!}\underset{k=0}{\overset{\mathrm{}}{}}\frac{B_k}{k!}\frac{(1)^k}{p+k}\mathrm{log}^{p+k}(1x),$$ (101) with $`B_k`$ the Bernoulli numbers.<sup>1</sup><sup>1</sup>1For other Nielsen integrals, as already for $`\text{Li}_l(x),l3`$, the serial expansion in $`\mathrm{log}(1x)`$ leads no longer to optimal representations due to the emergence of nested multiple series, see for an example. The Taylor expansion of $`\mathrm{log}^k(1x)`$ is given by $`\mathrm{log}^k(1x)=(1)^kk!{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\left[\begin{array}{c}n\\ m\end{array}\right]{\displaystyle \frac{x^l}{l!}}.`$ (104) Here $`\left[\begin{array}{c}n\\ m\end{array}\right]`$ denote the Stirling–numbers of the first kind . The Mellin–transform is given by $`\text{M}\left[\mathrm{log}^k(1x)\right](N)=(1)^kk!{\displaystyle \frac{1}{N}}S_{\underset{k}{\underset{}{\text{1, …, 1}}}}(N)={\displaystyle \frac{^k}{b^k}}B(N,b+1)\|_{b=0},`$ (105) where $`B(a,b)`$ denote Euler’s Beta function. The finite multiple harmonic sums $`S_{\underset{k}{\underset{}{\text{1, …, 1}}}}(N)`$ are given in , Eq. (158), and obey a determinant–representation. The elements of the corresponding matrix are numbers and the single harmonic sums $`S_l(N)={\displaystyle \underset{k=1}{\overset{N}{}}}{\displaystyle \frac{1}{k^l}}={\displaystyle \frac{(1)^{l1}}{(l1)!}}\psi ^{(l1)}(N+1)+c_l`$ (106) with $`c_1=\gamma _E`$ and $`c_l=\zeta (l)`$ for $`l>1`$. In this way one obtains the analytic continuations. With growing values of $`k`$ the explicit expressions become rather lengthly. For numerical computations one may use the recursion relation, Eq. (164), in Ref. , $`S_{\underset{k}{\underset{}{\text{1, …, 1}}}}(N)={\displaystyle \frac{1}{k}}{\displaystyle \underset{l=1}{\overset{k}{}}}S_l(N)S_{\underset{kl}{\underset{}{\text{1, …, 1}}}}(N)`$ (107) for complex values of $`N`$. Alternatively to these exact expressions one might also use the Taylor series Eq. (104) for which the Mellin–transform is easily calculated. The Stirling–numbers of the first kind can be represented using a recursion relation, cf. Eqs. (166,167) in Ref. . The Mellin transform of $`\mathrm{S}_{1,p}(x)`$ is thus given by $$\text{M}\left[\mathrm{S}_{1,p}(x)\right](N)=\frac{(1)^p}{p!}\underset{k=0}{\overset{\mathrm{}}{}}\frac{B_k}{k!}\frac{1}{p+k}\frac{1}{N}S_{\underset{p+k}{\underset{}{\text{1, …, 1}}}}(N).$$ (108) On the other hand the representation $$\text{M}\left[\mathrm{S}_{1,p}(x)\right](N)=\frac{\zeta (n+1)}{N}\frac{1}{N^2}S_{\underset{p}{\underset{}{\text{1, …, 1}}}}(N)$$ (109) is obtained by partial integration. This relation is more compact than the former. Eqs. (108,109) imply an interesting relation between the multiple finite harmonic sums of a single index and the Bernoulli numbers and the $`\zeta `$–function. Since not all Nielsen integrals which we are considering in the present paper can be expressed easily as a series in $`\mathrm{log}(1x)`$ and the corresponding series are found to converge not fast enough in a series of cases we are going to apply a somewhat modified representation. To guarantee a fast convergence also in the range of large values of $`x\stackrel{<}{}1`$ we subtract from some of the numerator functions a polynomial in $`\mathrm{log}(1x)`$ and $`x`$ of low degree. The remainder function is expanded into a polynomial by the MINIMAX–method. In both cases the Mellin transforms can be calculated analytically afterwards. In some cases further partial integrations have to be performed. Let us now discuss the individual cases in detail. Some of the Mellin–transforms are of the type $$\text{M}\left[\frac{\mathrm{log}^k(1+x)\mathrm{log}^k(2)}{x1}f(x)\right](N),$$ (110) where $`k=1,2`$. The first factor in Eq. (110) can be represented by a polynomial. Here the denominator is divided out of a polynomial representation of the numerator. The corresponding approximations are found applying the MINIMAX–method, $$\frac{\mathrm{log}^l(1+x)\mathrm{log}^l(2)}{x1}\underset{k=0}{\overset{L(k)}{}}b_k^{(l)}x^k,$$ (111) with $`L(1)=8`$ and $`L(2)=9`$. The expansion coefficients are given by $`\begin{array}{cccccc}b_0^{\left(1\right)}\hfill & =& \hfill 0.693147166991375\mathrm{E}+0& b_1^{\left(1\right)}\hfill & =& \hfill 0.306850436868254\mathrm{E}+0\\ b_2^{\left(1\right)}\hfill & =& \hfill 0.193078041088284\mathrm{E}+0& b_3^{\left(1\right)}\hfill & =& \hfill 0.139403892894644\mathrm{E}+0\\ b_4^{\left(1\right)}\hfill & =& \hfill 0.105269615988049\mathrm{E}+0& b_5^{\left(1\right)}\hfill & =& \hfill 0.746801353858524\mathrm{E}1\\ b_6^{\left(1\right)}\hfill & =& \hfill 0.427339135378207\mathrm{E}1& b_7^{\left(1\right)}\hfill & =& \hfill 0.161809049989783\mathrm{E}1\\ b_8^{\left(1\right)}\hfill & =& \hfill 0.288664611077007\mathrm{E}2& & & \end{array}`$ (113) and $`\begin{array}{cccccc}b_0^{\left(2\right)}\hfill & =& \hfill 0.480453024731510\mathrm{E}+0& b_1^{\left(2\right)}\hfill & =& \hfill 0.480450679641120\mathrm{E}+0\\ b_2^{\left(2\right)}\hfill & =& \hfill 0.519463586324817\mathrm{E}+0& b_3^{\left(2\right)}\hfill & =& \hfill 0.479285947990175\mathrm{E}+0\\ b_4^{\left(2\right)}\hfill & =& \hfill 0.427765744446172\mathrm{E}+0& b_5^{\left(2\right)}\hfill & =& \hfill 0.360855321373065\mathrm{E}+0\\ b_6^{\left(2\right)}\hfill & =& \hfill 0.263827078164263\mathrm{E}+0& b_7^{\left(2\right)}\hfill & =& \hfill 0.146927719341510\mathrm{E}+0\\ b_8^{\left(2\right)}\hfill & =& \hfill 0.525105367350968\mathrm{E}1& b_9^{\left(2\right)}\hfill & =& \hfill 0.874144396622167\mathrm{E}2\end{array}`$ (115) The related Mellin–transforms have the representation $`\text{M}\left[{\displaystyle \frac{\mathrm{log}(1+x)\mathrm{log}(2)}{x1}}\right](N)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{8}{}}}{\displaystyle \frac{b_k^{(1)}}{N+k+1}},`$ (116) $`\text{M}\left[{\displaystyle \frac{\mathrm{log}^2(1+x)\mathrm{log}^2(2)}{x1}}\right](N)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{9}{}}}{\displaystyle \frac{b_k^{(2)}}{N+k+1}},`$ (117) $`\text{M}\left[{\displaystyle \frac{\mathrm{log}(1+x)\mathrm{log}(2)}{x1}}\text{Li}_2(x)\right](N)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{8}{}}}{\displaystyle \frac{b_k^{(1)}}{N+k+1}}\left[\zeta (2){\displaystyle \frac{S_1(N+k+1)}{N+k+1}}\right]`$ (118) $`\text{M}\left[{\displaystyle \frac{\mathrm{log}(1+x)\mathrm{log}(2)}{x1}}\text{Li}_2(x)\right](N)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{8}{}}}{\displaystyle \frac{b_k^{(1)}}{N+k+1}}\left[{\displaystyle \frac{1}{2}}\zeta (2)+{\displaystyle \frac{\mathrm{log}(2)\beta (N+k+2)}{N+k+1}}\right]`$ In the case of the remaining Mellin–transforms of the type $$\text{M}\left[\frac{f(x)f(1)}{x1}\right](N)$$ (120) firstly the integral $$F(x)=_0^x𝑑z\frac{f(z)f(1)}{z1}$$ (121) is evaluated leading to $$\text{M}\left[\frac{f(x)f(1)}{x1}\right](N)=F(1)N\text{M}\left[F(x)\right](N1).$$ (122) We now seek for an appropriate representation of the functions $`F(x)`$. The respective functions are : $`{\displaystyle _0^x}𝑑z{\displaystyle \frac{\text{Li}_2(z)\zeta (2)}{z1}}`$ $`=`$ $`\left[\text{Li}_2(x)\zeta (2)\right]\mathrm{log}(1x)+2\text{S}_{1,2}(x)`$ (123) $`{\displaystyle _0^x}𝑑z{\displaystyle \frac{\text{Li}_2(z)+\zeta (2)/2}{z1}}`$ $`=`$ $`\left[\text{Li}_2(x)+{\displaystyle \frac{1}{2}}\zeta (2)\right]\mathrm{log}(1x)+I_1(x)`$ (124) $`{\displaystyle _0^x}𝑑z{\displaystyle \frac{\text{Li}_3(z)\zeta (3)}{z1}}`$ $`=`$ $`\left[\text{Li}_3(x)\zeta (3)\right]\mathrm{log}(1x)+{\displaystyle \frac{1}{2}}\text{Li}_2^2(x)`$ (125) $`{\displaystyle _0^x}𝑑z{\displaystyle \frac{\text{S}_{1,2}(z)\zeta (3)}{z1}}`$ $`=`$ $`\left[\text{S}_{1,2}(x)\zeta (3)\right]\mathrm{log}(1x)+3\mathrm{S}_{1,3}(x)`$ (126) Here one may use the representation $`\text{Li}_2(x)={\displaystyle _0^x}{\displaystyle \frac{dy}{y}}{\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}y^k={\displaystyle \underset{k=1}{\overset{9}{}}}{\displaystyle \frac{a_k^{(1)}}{k}}x^k.`$ (127) We derive approximations for $`\text{Li}_3(x),\text{S}_{1,2}(x),\text{Li}_2^2(x)`$ and $`I_1(x)`$ to be able to perform the Mellin–transform using these semi–analytic expressions. Here it is of particular importance to account for the $`\mathrm{log}^k(1x)`$ behavior as $`x1`$. One obtains $`\text{Li}_3(x)`$ $``$ $`P_0^{(2)}(x)+\left[P_1^{(2)}(x)+\mathrm{log}(1x)P_2^{(2)}(x)\right],`$ (128) $`\text{S}_{1,2}(x)`$ $``$ $`P_0^{(3)}(x)+\left[P_1^{(3)}(x)+\mathrm{log}(1x)P_2^{(3)}(x)+\mathrm{log}^2(1x)P_3^{(3)}(x)\right],`$ (129) $`\text{Li}_2^2(x)`$ $``$ $`P_0^{(4)}(x)+\left[P_1^{(4)}(x)+\mathrm{log}(1x)P_2^{(4)}(x)+\mathrm{log}^2(1x)P_3^{(4)}(x)\right],`$ (130) where $`P_{0,1,2}^{(k)}(x)`$ are the polynomials $`P_0^{(2)}(x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{12}{}}}c_k^{(2)}x^k`$ (131) $`P_1^{(2)}(x)`$ $`=`$ $`\zeta (3){\displaystyle \frac{11}{6}}\zeta (2)+{\displaystyle \frac{4}{3}}+\left(3\zeta (2){\displaystyle \frac{13}{4}}\right)x\left({\displaystyle \frac{3}{2}}\zeta (2){\displaystyle \frac{5}{2}}\right)x^2+\left({\displaystyle \frac{1}{3}}\zeta (2){\displaystyle \frac{7}{12}}\right)x^3`$ $`P_2^{(2)}(x)`$ $`=`$ $`1+{\displaystyle \frac{5}{2}}x2x^2+{\displaystyle \frac{1}{2}}x^3,`$ (133) with $`\begin{array}{cccccc}c_0^{\left(2\right)}\hfill & =& \hfill 0.480322239287459\mathrm{E}+0& c_1^{\left(2\right)}\hfill & =& \hfill 0.168480825099837\mathrm{E}+1\\ c_2^{\left(2\right)}\hfill & =& \hfill 0.209270571633447\mathrm{E}+1& c_3^{\left(2\right)}\hfill & =& \hfill 0.101728150522737\mathrm{E}+1\\ c_4^{\left(2\right)}\hfill & =& \hfill 0.160180000661971\mathrm{E}+0& c_5^{\left(2\right)}\hfill & =& \hfill 0.351982379713689\mathrm{E}+0\\ c_6^{\left(2\right)}\hfill & =& \hfill 0.141033369447519\mathrm{E}+1& c_7^{\left(2\right)}\hfill & =& \hfill 0.353344124579927\mathrm{E}+1\\ c_8^{\left(2\right)}\hfill & =& \hfill 0.593934899678262\mathrm{E}+1& c_9^{\left(2\right)}\hfill & =& \hfill 0.660019998525006\mathrm{E}+1\\ c_{10}^{\left(2\right)}\hfill & =& \hfill 0.466330491799074\mathrm{E}+1& c_{11}^{\left(2\right)}\hfill & =& \hfill 0.189825521858848\mathrm{E}+1\\ c_{12}^{\left(2\right)}\hfill & =& \hfill 0.339773000152805\mathrm{E}+0& & & \end{array}`$ $`P_0^{(3)}(x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{9}{}}}c_k^{(3)}x^k`$ (135) $`P_1^{(3)}(x)`$ $`=`$ $`\zeta (3){\displaystyle \frac{2035}{1728}}+{\displaystyle \frac{205}{144}}x{\displaystyle \frac{95}{288}}x^2+{\displaystyle \frac{43}{432}}x^3{\displaystyle \frac{1}{64}}x^4`$ (136) $`P_2^{(3)}(x)`$ $`=`$ $`{\displaystyle \frac{205}{144}}{\displaystyle \frac{25}{12}}x+{\displaystyle \frac{23}{24}}x^2{\displaystyle \frac{13}{36}}x^3+{\displaystyle \frac{1}{16}}x^4`$ (137) $`P_3^{(3)}(x)`$ $`=`$ $`{\displaystyle \frac{25}{24}}+2x{\displaystyle \frac{3}{2}}x^2+{\displaystyle \frac{2}{3}}x^3{\displaystyle \frac{1}{8}}x^4,`$ (138) with $`\begin{array}{cccccc}c_0^{\left(3\right)}\hfill & =& \hfill 0.243948949064443\mathrm{E}1& c_1^{\left(3\right)}\hfill & =& \hfill 0.000005136294145\mathrm{E}+0\\ c_2^{\left(3\right)}\hfill & =& \hfill 0.249849075518710\mathrm{E}+0& c_3^{\left(3\right)}\hfill & =& \hfill 0.498290708990997\mathrm{E}+0\\ c_4^{\left(3\right)}\hfill & =& \hfill 0.354866791547134\mathrm{E}+0& c_5^{\left(3\right)}\hfill & =& \hfill 0.522116678353452\mathrm{E}1\\ c_6^{\left(3\right)}\hfill & =& \hfill 0.648354706049337\mathrm{E}1& c_7^{\left(3\right)}\hfill & =& \hfill 0.644165053822532\mathrm{E}1\\ c_8^{\left(3\right)}\hfill & =& \hfill 0.394927322542075\mathrm{E}1& c_9^{\left(3\right)}\hfill & =& \hfill 0.100879370657869\mathrm{E}1\end{array}`$ $`P_0^{(4)}(x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{12}{}}}c_k^{(4)}x^k`$ (140) $`P_1^{(4)}(x)`$ $`=`$ $`{\displaystyle \frac{257}{144}}{\displaystyle \frac{205}{72}}\zeta (2)+\zeta (2)^2\left({\displaystyle \frac{167}{36}}{\displaystyle \frac{25}{6}}\zeta (2)\right)x+\left({\displaystyle \frac{101}{24}}{\displaystyle \frac{23}{12}}\zeta (2)\right)x^2`$ (141) $`\left({\displaystyle \frac{59}{36}}{\displaystyle \frac{13}{18}}\zeta (2)\right)x^3+\left({\displaystyle \frac{41}{144}}{\displaystyle \frac{1}{8}}\zeta (2)\right)x^4`$ $`P_2^{(4)}(x)`$ $`=`$ $`\left({\displaystyle \frac{167}{36}}{\displaystyle \frac{25}{6}}\zeta (2)\right)+\left({\displaystyle \frac{235}{18}}8\zeta (2)\right)x\left({\displaystyle \frac{40}{3}}6\zeta (2)\right)x^2+\left({\displaystyle \frac{109}{18}}{\displaystyle \frac{8}{3}}\zeta (2)\right)x^3`$ (142) $`\left({\displaystyle \frac{41}{36}}{\displaystyle \frac{1}{2}}\zeta (2)\right)x^4`$ $`P_3^{(4)}(x)`$ $`=`$ $`{\displaystyle \frac{35}{12}}{\displaystyle \frac{26}{3}}x+{\displaystyle \frac{19}{2}}x^2{\displaystyle \frac{14}{3}}x^3+{\displaystyle \frac{11}{12}}x^4,`$ (143) with $`\begin{array}{cccccc}c_0^{\left(4\right)}\hfill & =& \hfill 0.192962504274437\mathrm{E}+0& c_1^{\left(4\right)}\hfill & =& \hfill 0.000005641557253\mathrm{E}+0\\ c_2^{\left(4\right)}\hfill & =& \hfill 0.196891075399448\mathrm{E}+1& c_3^{\left(4\right)}\hfill & =& \hfill 0.392919138747074\mathrm{E}+1\\ c_4^{\left(4\right)}\hfill & =& \hfill 0.290306105685546\mathrm{E}+1& c_5^{\left(4\right)}\hfill & =& \hfill 0.992890266001707\mathrm{E}+0\\ c_6^{\left(4\right)}\hfill & =& \hfill 0.130026190226546\mathrm{E}+1& c_7^{\left(4\right)}\hfill & =& \hfill 0.341870577921103\mathrm{E}+1\\ c_8^{\left(4\right)}\hfill & =& \hfill 0.576763902370864\mathrm{E}+1& c_9^{\left(4\right)}\hfill & =& \hfill 0.645554138192407\mathrm{E}+1\\ c_{10}^{\left(4\right)}\hfill & =& \hfill 0.459405622046138\mathrm{E}+1& c_{11}^{\left(4\right)}\hfill & =& \hfill 0.188510809558304\mathrm{E}+1\\ c_{12}^{\left(4\right)}\hfill & =& \hfill 0.340476080290674\mathrm{E}+0& & & \end{array}`$ (145) These representations hold at an accuracy of better than $`2\times 10^8,3\times 10^8`$, and $`2\times 10^8`$, respectively, for $`xϵ[0,1]`$. In the case of the function $`I_1(x)`$, Eq. (25), the function $`\mathrm{log}(1+x)`$ in the integrand can be approximated by Eq. (65). The integral can then be evaluated analytically and takes the form $$I_1(x)=P_0^{(5)}(x)+P_2^{(5)}(x)\mathrm{log}(1x),$$ (146) with $`P_0^{(5)}(x)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{9}{}}}c_k^{(5)}x^k`$ (147) $`P_2^{(5)}(x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{9}{}}}d_k^{(5)}x^k.`$ (148) The coefficients $`c_k^{(5)}`$ and $`d_k^{(5)}`$ are algebraically related to the coefficients $`a_k^{(1)}`$, Eq. (65). Their numerical values are $`\begin{array}{cccccc}c_1^{\left(5\right)}\hfill & =& \hfill 0.822467033400776\mathrm{E}+0& c_2^{\left(5\right)}\hfill & =& \hfill 0.887664705657325\mathrm{E}1\\ c_3^{\left(5\right)}\hfill & =& \hfill 0.241549406045162\mathrm{E}1& c_4^{\left(5\right)}\hfill & =& \hfill 0.965074750946139\mathrm{E}2\\ c_5^{\left(5\right)}\hfill & =& \hfill 0.470587487919749\mathrm{E}2& c_6^{\left(5\right)}\hfill & =& \hfill 0.246014308378549\mathrm{E}2\\ c_7^{\left(5\right)}\hfill & =& \hfill 0.116431121874067\mathrm{E}2& c_8^{\left(5\right)}\hfill & =& \hfill 0.395705193848026\mathrm{E}3\\ c_9^{\left(5\right)}\hfill & =& \hfill 0.664699010014505\mathrm{E}4& & & \end{array}`$ (150) $`\begin{array}{cccccc}d_0^{\left(5\right)}\hfill & =& \hfill 0.822467033400776\mathrm{E}+0& d_1^{\left(5\right)}\hfill & =& \hfill 0.999999974532241\mathrm{E}+0\\ d_2^{\left(5\right)}\hfill & =& \hfill 0.249997762945014\mathrm{E}+0& d_3^{\left(5\right)}\hfill & =& \hfill 0.111067811851394\mathrm{E}+0\\ d_4^{\left(5\right)}\hfill & =& \hfill 0.621323644338330\mathrm{E}1& d_5^{\left(5\right)}\hfill & =& \hfill 0.382902328987004\mathrm{E}1\\ d_6^{\left(5\right)}\hfill & =& \hfill 0.229110370338977\mathrm{E}1& d_7^{\left(5\right)}\hfill & =& \hfill 0.113158200819689\mathrm{E}1\\ d_8^{\left(5\right)}\hfill & =& \hfill 0.376387065979726\mathrm{E}2& d_9^{\left(5\right)}\hfill & =& \hfill 0.598229109013054\mathrm{E}3\end{array}`$ (152) This representation holds at an accuracy of better than $`3\times 10^8`$ for $`xϵ[0,1]`$. The Mellin–transforms of the above polynomials are : $`\text{M}\left[{\displaystyle \underset{k=0}{\overset{m}{}}}A_kx^k\right](N)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{m+1}{}}}{\displaystyle \frac{A_{k1}}{N+k}}`$ (153) $`\text{M}\left[{\displaystyle \underset{k=0}{\overset{m}{}}}A_kx^k\mathrm{log}^l(1x)\right](N)`$ $`=`$ $`(1)^ll!{\displaystyle \underset{k=1}{\overset{m}{}}}{\displaystyle \frac{A_{k1}}{N+k}}S_{\underset{l}{\underset{}{\text{1, …, 1}}}}(N+k).`$ (154) Here, $`S_{1,1}(N)`$ $`=`$ $`{\displaystyle \frac{1}{2!}}\left[S_1^2(N)+S_2(N)\right]`$ (155) $`S_{1,1,1}(N)`$ $`=`$ $`{\displaystyle \frac{1}{3!}}\left[S_1^3(N)+3S_1(N)S_2(N)+2S_3(N)\right].`$ (156) The Mellin transforms of the functions $`g_{18}(x)g_{22}(x)`$ are then given represented by $`\text{M}\left[{\displaystyle \frac{\text{Li}_2(x)\zeta (2)}{x1}}\right](N)`$ $`=`$ $`{\displaystyle \frac{1}{N}}\left[S_1^2(N)+S_2(N)\right]\zeta (2)S_1(N)+{\displaystyle \underset{k=0}{\overset{11}{}}}c_k^{(1)}{\displaystyle \frac{N}{N+k}}S_1(N+k)`$ (157) $`{\displaystyle \underset{k=0}{\overset{3}{}}}P_{1,k}^{(2)}{\displaystyle \frac{N}{N+k}}\left[S_1^2(N+k)+S_2(N+k)\right]`$ $`\text{M}\left[{\displaystyle \frac{\text{Li}_2(x)+\zeta (2)/2}{x1}}\right](N)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\zeta (2)S_1(N){\displaystyle \underset{k=1}{\overset{9}{}}}{\displaystyle \frac{a_k^{(1)}}{k}}S_1(N+k)`$ (158) $`\text{M}\left[{\displaystyle \frac{\text{Li}_3(x)\zeta (3)}{x1}}\right](N)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\zeta (2)^2\zeta (3)S_1(N)+{\displaystyle \underset{k=0}{\overset{12}{}}}c_k^{(2)}{\displaystyle \frac{N}{N+k}}S_1(N+k)`$ (159) $`{\displaystyle \underset{k=0}{\overset{3}{}}}P_{2,k}^{(2)}{\displaystyle \frac{N}{N+k}}\left[S_1^2(N+k)+S_2(N+k)\right]{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=0}{\overset{12}{}}}c_k^{(4)}{\displaystyle \frac{N}{N+k}}`$ $`+{\displaystyle \frac{N}{2}}{\displaystyle \underset{k=0}{\overset{4}{}}}\left[P_{2,k}^{(4)}{\displaystyle \frac{S_1(N+k)}{N+k}}P_{3,k}^{(4)}{\displaystyle \frac{S_1^2(N+k)+S_2(N+k)}{N+k}}\right]`$ $`\text{M}\left[{\displaystyle \frac{\text{S}_{1,2}(x)\zeta (3)}{x1}}\right](N)`$ $`=`$ $`\zeta (3)S_1(N)+{\displaystyle \frac{1}{2N}}\left[S_1^3(N)+2S_1(N)S_2(N)+2S_3(N)\right]`$ (160) $`+{\displaystyle \underset{k=0}{\overset{9}{}}}c_k^{(3)}{\displaystyle \frac{N}{N+k}}S_1(N+k)+{\displaystyle \underset{k=0}{\overset{4}{}}}{\displaystyle \frac{N}{N+k}}\{P_{3,k}^{(3)}`$ $`\left[S_1^3(N+k)+3S_1(N+k)S_2(N+k)+2S_3(N+k)\right]`$ $`P_{2,k}^{(3)}[S_1^2(N+k)+S_2(N+k)]\}`$ $`\text{M}\left[{\displaystyle \frac{\mathrm{log}(x)\text{Li}_2(x)}{x1}}\right](N)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{11}{}}}c_k^{(1)}\psi ^{}(N+k+1)`$ (161) $`{\displaystyle \underset{k=0}{\overset{3}{}}}P_{2,k}^{(1)}\left[S_1(N+k)\psi ^{}(N+k+1){\displaystyle \frac{1}{2}}\psi ^{\prime \prime }(N+k+1)\right]`$ A set of other integrals we are going to re–write in terms of integrals of the type $`_0^x𝑑zf_2(z)/(1+z)`$, the Mellin–transforms of which are then evaluated using Eq. (64). The following integral relations are obtained : $`{\displaystyle _0^x}𝑑z{\displaystyle \frac{\text{Li}_3(z)+(3/4)\zeta (3)}{z1}}`$ $`=`$ $`\left[\text{Li}_3(x)+{\displaystyle \frac{3}{4}}\zeta (3)\right]\mathrm{log}(1x)+\text{Li}_2(x)\text{Li}_2(x)`$ (162) $`+\text{Li}_3(x)\mathrm{log}(1+x){\displaystyle _0^x}𝑑z{\displaystyle \frac{\text{Li}_3(z)}{1+z}}`$ $`{\displaystyle _0^x}𝑑z{\displaystyle \frac{I_1(z)+(5/8)\zeta (3)}{z1}}`$ $`=`$ $`\left[I_1(x)+{\displaystyle \frac{5}{8}}\zeta (3)\right]\mathrm{log}(1x)2\text{S}_{1,2}(x)\mathrm{log}(1+x)`$ (163) $`+2{\displaystyle _0^x}𝑑z{\displaystyle \frac{\text{S}_{1,2}(z)}{1+z}}`$ $`{\displaystyle _0^x}𝑑z{\displaystyle \frac{\text{S}_{1,2}(z)\zeta (3)/8}{z1}}`$ $`=`$ $`\left[\text{S}_{1,2}(x){\displaystyle \frac{1}{8}}\zeta (3)\right]\mathrm{log}(1x){\displaystyle \frac{1}{2}}I_1(x)\mathrm{log}(1+x)`$ (164) $`+{\displaystyle \frac{1}{2}}{\displaystyle _0^x}𝑑z{\displaystyle \frac{I_1(z)}{1+z}}`$ Here one may represent $`\text{Li}_3(x)`$ and $`\text{S}_{1,2}(x)`$ by $`\text{Li}_3(x)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{9}{}}}{\displaystyle \frac{a_k^{(1)}}{k^2}}x^k`$ (165) $`\text{S}_{1,2}(x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=2}{\overset{11}{}}}{\displaystyle \frac{a_k^{(2)}}{k}}x^k.`$ (166) The Mellin transforms pf $`g_{22}(x)g_{25}(x)`$ are represented by $`\text{M}\left[{\displaystyle \frac{\text{Li}_3(x)+3\zeta (3)/3}{x1}}\right](N)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\zeta (2)^2+\zeta (3)\mathrm{log}(2)+{\displaystyle \frac{3}{4}}\zeta (3)S_1(N)\text{M}\left[{\displaystyle \frac{\text{Li}_3(x)}{1+x}}\right](N)`$ (167) $`{\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}{\displaystyle \frac{N}{N+k}}\{\zeta (3){\displaystyle \frac{\zeta (2)}{N+k}}`$ $`{\displaystyle \frac{\zeta (2)}{k}}+S_1(N+k)[{\displaystyle \frac{1}{(N+k)^2}}+{\displaystyle \frac{1}{k^2}}+{\displaystyle \frac{1}{k(N+k)}}]\}`$ $`\text{M}\left[{\displaystyle \frac{I_1(x)+5\zeta (3)/8}{x1}}\right](N)=2\zeta (3)\mathrm{log}(2)+2\text{M}\left[{\displaystyle \frac{S_{1,2}(x)}{1+x}}\right](N)`$ (168) $`+{\displaystyle \frac{5}{8}}\zeta (3)S_1(N)+{\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}{\displaystyle \frac{N}{N+k}}\left[2\zeta (3){\displaystyle \frac{S_1^2(N+k)+S_2(N+k)}{N+k}}\right]`$ $`+{\displaystyle \underset{k=1}{\overset{9}{}}}c_k^{(5)}{\displaystyle \frac{N}{N+k}}S_1(N+k){\displaystyle \underset{k=0}{\overset{9}{}}}d_k^{(5)}{\displaystyle \frac{N}{N+k}}\left[S_1^2(N+k)+S_2(N+k)\right]`$ $`\text{M}\left[{\displaystyle \frac{\text{S}_{1,2}(x)\zeta (3)/8}{x1}}\right](N)`$ $`=`$ $`{\displaystyle \frac{5}{16}}\zeta (3)\mathrm{log}(2)+{\displaystyle \frac{1}{2}}\text{M}\left[{\displaystyle \frac{I_1(x)}{1+x}}\right](N){\displaystyle \frac{1}{8}}\zeta (3)S_1(N)`$ (169) $`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=2}{\overset{11}{}}}{\displaystyle \frac{a_k^{(2)}}{k}}{\displaystyle \frac{N}{N+k}}S_1(N+k)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}{\displaystyle \frac{N}{N+k}}[{\displaystyle \frac{5}{8}}\zeta (3)`$ $`+{\displaystyle \underset{l=1}{\overset{9}{}}}a_l^{(1)}{\displaystyle \frac{S_1(N+k+l)}{N+k+l}}]`$ We finally would like to add a remark on a recent analysis of the two–loop coefficient functions of deep inelastic scattering . The representation given there also contains the harmonic sum $`S_{1,1,1,1}(N)`$ which was not needed to express the individual Mellin transforms in Refs. , cf. Ref. , which may be caused due to the algebraic relations being applied in . This sum has the representation $`S_{1,1,1,1}(N)=(1)^{N+1}{\displaystyle \frac{1}{6}}\text{M}\left[{\displaystyle \frac{\mathrm{log}^3(1x)}{1+x}}\right](N)\text{Li}_4\left({\displaystyle \frac{1}{2}}\right).`$ (170) The Mellin transform of the function $`\mathrm{log}^3(1x)/(1+x)`$ for complex argument is easily obtained, $`\text{M}\left[{\displaystyle \frac{\mathrm{log}^3(1x)}{1+x}}\right](N)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{9}{}}}a_k^{(1)}\{{\displaystyle \frac{k}{N+k}}[S_1^3(N+k)+3S_1(N+k)S_2(N+k)+2S_3(N+k)]`$ (171) $`[S_1^3(k)+3S_1(k)S_2(k)+2S_3(k)]\}6\text{Li}_4({\displaystyle \frac{1}{2}}).`$ The above relations allow to express the Mellin transforms of all functions emerging in the coefficient functions and anomalous dimensions of massless gauge theories up to two–loop order and to evaluate their analytic continuation from integer to complex arguments. ## 6 The Code ANCONT ### 6.1 General Structure The code ANCONT calculates the Mellin transforms of the basic functions $`g_i(x)`$ both for integer and complex values of the Mellin index N. The parameters of the code are initialized in SUBROUTINE ACINI. The calculations are performed in SUBROUTINE ACRUN. The code ends with SUBROUTINE ACEND. ### 6.2 USER routines The user may access parts of the code via the user routines UINIT, URUN and UOUT. These routines are called in ACINI, ACRUN and ACEND, respectively. SUBROUTINE UINIT may be used to actualize the running parameters. Via URUN the user may built her/his own Mellin convolutions or other functions of the basic functions and their Mellin transforms for positive integer or complex Mellin index using the functions ACG$`i`$, FCT$`i`$ or XCG$`i`$, which are described below, as a library. SUBROUTINE UOUT may be used as an output-interface at the end of the code. To transfer data between different user routines the user may define COMMON-blocks named COMMON/USxxxx/ … ### 6.3 Initialization The main parameters and constants of the code are defined in SUBROUTINE ACINI. This routine calls the subroutines INVINI, DEFAUL and the user routine UINIT. SUBROUTINE INVINI sets parameters and constants related to the Mellin inversion. The default running parameters of the code are set in SUBROUTINE DEFAUL. These are : IRUN = 0 ITEST1 = 1 ITEST2 = 1 ITEST3 = 1 IMIN = 1 IMAX = 26 IAPP = 1 NMIN = 1 NMAX = 20 EPS = 1.0D-9 The running parameters are printed by SUBROUTINE WROUT. The parameters ITEST$`i`$ initialize tests of the code, which are inactive for ITEST$`i`$.NE.1. For ITEST1.EQ.1 the representation of the moments of the basic functions $`g_i(x)`$ for positive integer index by harmonic sums are compared to those obtained by numerical integration. ITEST2.EQ.1 induces a comparison of the representations for the analytic continuation of the Mellin transform of the basic functions with those by the harmonic sums for positive integer index. For ITEST3.EQ.1 a comparison is made for the Mellin inversion using the analytic continuations of the Mellin transforms of the basic functions and the numerical representations of the basic functions in the range $`xϵ[10^7,0.99]`$. IMIN and IMAX mark the index range of basic functions to be used. Likewise NMIN and NMAX set the minimum and maximum positive integer moment for the tests. IAPP selects the representation of the analytic continuation for $`g_3(x)=\text{Li}_2(x)/(1+x)`$ to compare the representations Eq. (78) IAPP = 1 and (87) IAPP=2, (88) IAPP=3, respectively. EPS (REAL\8) is a pilot parameter for the numerical integration of the program DAIND and denotes the relative numerical accuracy to be obtained. The above parameters are available through the COMMON-blocks COMMON/TEST/ ITEST1, ITEST2,ITEST3 COMMON/RUN / IRUN COMMON/MOMPA/ NMIN,NMAX COMMON/FUNPA/ IMIN,IMAX COMMON/IAPP/ IAPP COMMON/EP/ EPS. ### 6.4 Running The code provides three main lines, which are initialized setting the parameters ITEST1, ITEST2, ITEST3 = 1, respectively. For ITEST1= 1 the representations of the integer moments of the basic functions $`g_i(x)`$ in terms of harmonic sums are tested comparing them with the corresponding numerical integrals. The value of the Mellin moment is calculated. The relative accuracy comparing both calculations RAT = VAL1/VAL2 -1 and the value of the moment are provided. The following test-output is obtained for $`k=9`$, $`g_9(x)=\text{S}_{1,2}(x)/(1+x)`$, $`n_1=2,n_2=20,val=1.0`$D-9. Here N denotes the Mellin index, RATk the relative accuracy comparing the numerical integration and the representation by harmonic sums, and VAL the value of the Mellin moment. ``` ************************************************************* N,RAT9,VAL= 2 -1.983302411190380E-12 1.784970126521851E-02 N,RAT9,VAL= 3 -1.992739306899693E-12 1.447619078834311E-02 N,RAT9,VAL= 4 -1.742939126359033E-12 1.216425178651454E-02 N,RAT9,VAL= 5 -1.722177955798543E-12 1.048399327332255E-02 N,RAT9,VAL= 6 -1.469158128486470E-12 9.208975593568651E-03 N,RAT9,VAL= 7 -1.485589429250922E-12 8.209008584718454E-03 N,RAT9,VAL= 8 -1.187272502534142E-12 7.404084086538535E-03 N,RAT9,VAL= 9 -1.256883486178140E-12 6.742389451876160E-03 N,RAT9,VAL= 10 -1.036948304999896E-12 6.188923213501677E-03 N,RAT9,VAL= 11 -8.926193117986259E-13 5.719203265306224E-03 N,RAT9,VAL= 12 -8.733014311701481E-13 5.315598831318883E-03 N,RAT9,VAL= 13 -7.922551503725117E-13 4.965094655720756E-03 N,RAT9,VAL= 14 -5.925260282424460E-13 4.657875778975199E-03 N,RAT9,VAL= 15 -6.212808045802376E-13 4.386402245697693E-03 N,RAT9,VAL= 16 -3.986810881428937E-13 4.144786898528877E-03 N,RAT9,VAL= 17 -4.581890422628021E-13 3.928366556207149E-03 N,RAT9,VAL= 18 -1.886268918838141E-13 3.733399977739974E-03 N,RAT9,VAL= 19 -3.899103262483550E-13 3.556850973662370E-03 N,RAT9,VAL= 20 -1.719735465144367E-13 3.396229940556720E-03 ************************************************************ ``` For ITEST2= 1 the representations of the Mellin transforms of the basic functions $`g_i(x)`$ valid for complex arguments are compared to the representation in terms of harmonic sums at positive integer argument. The value of the Mellin moment is calculated. The relative accuracy comparing both calculations RAT = VAL1/VAL2 -1 and the value of the moment are provided. The relative accuracy of these representations are given in Table 2 for $`N=2`$ to 20. \| $`N`$ \| 1 \| 2 \| 3 \| 4 \| 5 \| 6 \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| 2 \| -1.06E-09 \| 5.62E-10 \| 7.61E-09 \| 3.52E-09 \| -6.94E-09 \| 4.67E-09 \| \| 3 \| -2.09E-09 \| 1.06E-09 \| 1.18E-08 \| 6.27E-09 \| -1.15E-08 \| 7.85E-09 \| \| 4 \| -3.48E-09 \| 1.72E-09 \| 1.69E-08 \| 9.92E-09 \| -1.79E-08 \| 1.20E-08 \| \| 5 \| -5.28E-09 \| 2.54E-09 \| 2.32E-08 \| 1.46E-08 \| -2.58E-08 \| 1.72E-08 \| \| 6 \| -7.50E-09 \| 3.55E-09 \| 3.07E-08 \| 2.04E-08 \| -3.62E-08 \| 2.36E-08 \| \| 7 \| -1.02E-08 \| 3.97E-08 \| 4.74E-09 \| 2.75E-08 \| -4.92E-08 \| 3.14E-08 \| \| 8 \| -1.34E-08 \| 6.15E-09 \| 5.01E-08 \| 3.61E-08 \| -6.47E-08 \| 4.07E-08 \| \| 9 \| -1.72E-08 \| 7.78E-09 \| 6.22E-08 \| 4.63E-08 \| -8.21E-08 \| 5.17E-08 \| \| 10 \| -2.17E-08 \| 9.67E-09 \| 7.58E-08 \| 5.82E-08 \| -1.01E-07 \| 6.43E-08 \| \| 11 \| -2.69E-08 \| 1.18E-08 \| 9.09E-08 \| 7.18E-08 \| -1.12E-07 \| 7.86E-08 \| \| 12 \| -3.28E-08 \| 1.43E-08 \| 1.07E-07 \| 8.70E-08 \| -1.38E-07 \| 9.45E-08 \| \| 13 \| -3.94E-08 \| 1.70E-08 \| 1.25E-07 \| 1.04E-07 \| -1.56E-07 \| 1.12E-07 \| \| 14 \| -4.67E-08 \| 2.01E-08 \| 1.44E-07 \| 1.22E-07 \| -1.72E-07 \| 1.30E-07 \| \| 15 \| -5.47E-08 \| 2.34E-08 \| 1.64E-07 \| 1.41E-07 \| -1.85E-07 \| 1.50E-07 \| \| 16 \| -6.34E-08 \| 2.71E-08 \| 1.84E-07 \| 1.61E-07 \| -1.96E-07 \| 1.70E-07 \| \| 17 \| -7.27E-08 \| 3.11E-08 \| 2.05E-07 \| 1.82E-07 \| -2.03E-07 \| 1.91E-07 \| \| 18 \| -8.26E-08 \| 3.54E-08 \| 2.26E-07 \| 2.03E-07 \| -2.08E-07 \| 2.12E-07 \| \| 19 \| -9.29E-08 \| 3.99E-08 \| 2.47E-07 \| 2.24E-07 \| -2.09E-07 \| 2.34E-07 \| \| 20 \| -1.04E-07 \| 4.47E-08 \| 2.68E-07 \| 2.46E-07 \| -2.07E-07 \| 2.56E-07 \| \| $`N`$ \| 7 \| 8 \| 9 \| 10 \| 11 \| 12 \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| 2 \| 3.67E-09 \| 2.47E-08 \| 6.04E-09 \| -9.99E-09 \| 9.43E-09 \| -3.16E-10 \| \| 3 \| 6.48E-09 \| 3.25E-08 \| 9.77E-09 \| -1.03E-08 \| 1.40E-08 \| -6.46E-10 \| \| 4 \| 1.02E-08 \| 4.15E-08 \| 1.45E-08 \| -9.44E-09 \| 1.93E-08 \| -1.05E-09 \| \| 5 \| 1.49E-08 \| 5.16E-08 \| 2.02E-08 \| -7.45E-09 \| 2.54E-08 \| -1.52E-09 \| \| 6 \| 2.08E-08 \| 6.29E-08 \| 2.72E-08 \| -4.20E-09 \| 3.22E-08 \| -2.07E-09 \| \| 7 \| 2.80E-08 \| 7.56E-08 \| 3.57E-08 \| 4.18E-10 \| 3.98E-08 \| -2.72E-09 \| \| 8 \| 3.67E-08 \| 8.94E-08 \| 4.57E-08 \| 6.48E-09 \| 4.80E-08 \| -3.46E-09 \| \| 9 \| 4.71E-08 \| 1.04E-07 \| 5.74E-08 \| 1.40E-08 \| 5.68E-08 \| -4.31E-09 \| \| 10 \| 5.91E-08 \| 1.20E-07 \| 7.09E-08 \| 2.31E-08 \| 6.59E-08 \| -5.28E-09 \| \| 11 \| 7.28E-08 \| 1.37E-07 \| 8.60E-08 \| 3.36E-08 \| 7.53E-08 \| -6.37E-09 \| \| 12 \| 8.81E-08 \| 1.54E-07 \| 1.03E-07 \| 4.55E-08 \| 8.48E-08 \| -7.57E-09 \| \| 13 \| 1.05E-07 \| 1.72E-07 \| 1.21E-07 \| 5.90E-08 \| 9.42E-08 \| -8.89E-09 \| \| 14 \| 1.23E-07 \| 1.91E-07 \| 1.40E-07 \| 7.28E-08 \| 1.03E-07 \| -1.03E-08 \| \| 15 \| 1.42E-07 \| 2.09E-07 \| 1.60E-07 \| 8.80E-08 \| 1.12E-07 \| -1.18E-08 \| \| 16 \| 1.62E-07 \| 2.27E-07 \| 1.81E-07 \| 1.04E-07 \| 1.21E-07 \| -1.34E-08 \| \| 17 \| 1.83E-07 \| 2.46E-07 \| 2.03E-07 \| 1.21E-07 \| 1.29E-07 \| -1.51E-08 \| \| 18 \| 2.04E-07 \| 2.64E-07 \| 2.26E-07 \| 1.38E-07 \| 1.36E-07 \| -1.68E-08 \| \| 19 \| 2.26E-07 \| 2.81E-07 \| 2.48E-07 \| 1.55E-07 \| 1.43E-07 \| -1.86E-08 \| \| 20 \| 2.48E-07 \| 2.98E-07 \| 2.71E-07 \| 1.73E-07 \| 1.49E-07 \| -2.04E-08 \| Table 2: Relative accuracy of the Mellin transforms comparing the approximative representations valid for complex arguments and the representation in terms of harmonic sums for $`Nϵ[2,20]`$ \| $`N`$ \| 13 \| 14 \| 15 \| 16 \| 17 \| 18 \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| 2 \| -9.79E-10 \| -2.70E-10 \| -9.02E-10 \| -7.02E-10 \| 9.31E-10 \| -6.00E-10 \| \| 3 \| -2.08E-09 \| -3.63E-10 \| -1.14E-09 \| -9.13E-10 \| 1.18E-09 \| -1.13E-09 \| \| 4 \| -3.56E-09 \| -4.59E-10 \| -1.40E-09 \| -1.14E-09 \| 1.44E-09 \| -1.80E-09 \| \| 5 \| -5.45E-09 \| -5.60E-10 \| -1.68E-09 \| -1.34E-09 \| 1.72E-09 \| -2.60E-09 \| \| 6 \| -7.80E-09 \| -6.65E-10 \| -1.96E-09 \| -1.64E-09 \| 2.01E-09 \| -3.54E-09 \| \| 7 \| -1.06E-08 \| -7.75E-10 \| -2.27E-09 \| -1.92E-09 \| 2.32E-09 \| -4.60E-09 \| \| 8 \| -1.40E-08 \| -8.93E-10 \| -2.60E-09 \| -2.23E-09 \| 2.65E-09 \| -5.80E-09 \| \| 9 \| -1.81E-08 \| -1.02E-09 \| -2.95E-09 \| -2.57E-09 \| 3.00E-09 \| -7.14E-09 \| \| 10 \| -2.28E-08 \| -1.15E-09 \| -3.32E-09 \| -2.92E-09 \| 3.39E-09 \| -8.60E-09 \| \| 11 \| -2.82E-08 \| -1.30E-09 \| -3.69E-09 \| -3.29E-09 \| 3.79E-09 \| -1.02E-08 \| \| 12 \| -3.43E-08 \| -1.45E-09 \| -4.06E-09 \| -3.68E-09 \| 4.22E-09 \| -1.19E-08 \| \| 13 \| -4.11E-08 \| -1.62E-09 \| -4.43E-09 \| -4.06E-09 \| 4.66E-09 \| -1.38E-08 \| \| 14 \| -4.86E-08 \| -1.78E-09 \| -4.79E-09 \| -4.44E-09 \| 5.11E-09 \| -1.58E-08 \| \| 15 \| -5.68E-08 \| -1.96E-09 \| -5.14E-09 \| -4.80E-09 \| 5.56E-09 \| -1.80E-08 \| \| 16 \| -6.55E-08 \| -2.14E-09 \| -5.47E-09 \| -5.16E-09 \| 6.02E-09 \| -2.04E-08 \| \| 17 \| -7.49E-08 \| -2.31E-09 \| -5.77E-09 \| -5.49E-09 \| 6.47E-09 \| -2.29E-08 \| \| 18 \| -8.47E-08 \| -2.49E-09 \| -6.06E-09 \| -5.80E-09 \| 6.91E-09 \| -2.56E-08 \| \| 19 \| -9.49E-08 \| -2.66E-09 \| -6.32E-09 \| -6.09E-09 \| 7.34E-09 \| -2.85E-08 \| \| 20 \| -1.06E-07 \| -2.83E-09 \| -6.55E-09 \| -6.36E-09 \| 7.75E-09 \| -3.15E-08 \| \| $`N`$ \| 19 \| 20 \| 21 \| 22 \| 23 \| 24 \| 25 \| 26 \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| 2 \| -0.248E-08 \| 0.108E-08 \| 0.646E-09 \| -0.131E-09 \| 0.431E-08 \| -0.408E-08 \| 0.353E-08 \| -0.489E-08 \| \| 3 \| -0.303E-08 \| 0.212E-08 \| 0.114E-08 \| -0.166E-09 \| 0.686E-08 \| -0.500E-08 \| 0.461E-08 \| -0.795E-08 \| \| 4 \| -0.353E-08 \| 0.351E-08 \| 0.174E-08 \| -0.200E-09 \| 0.963E-08 \| -0.586E-08 \| 0.568E-08 \| -0.115E-07 \| \| 5 \| -0.397E-08 \| 0.525E-08 \| 0.243E-08 \| -0.235E-09 \| 0.126E-07 \| -0.668E-08 \| 0.674E-08 \| -0.154E-07 \| \| 6 \| -0.435E-08 \| 0.733E-08 \| 0.321E-08 \| -0.270E-09 \| 0.157E-07 \| -0.746E-08 \| 0.779E-08 \| -0.198E-07 \| \| 7 \| -0.470E-08 \| 0.976E-08 \| 0.409E-08 \| -0.305E-09 \| 0.189E-07 \| -0.821E-08 \| 0.882E-08 \| -0.248E-07 \| \| 8 \| -0.499E-08 \| 0.126E-07 \| 0.506E-08 \| -0.341E-09 \| 0.223E-07 \| -0.893E-08 \| 0.985E-08 \| -0.303E-07 \| \| 9 \| -0.525E-08 \| 0.157E-07 \| -0.377E-09 \| 0.612E-08 \| 0.257E-07 \| -0.963E-08 \| 0.109E-07 \| -0.365E-07 \| \| 10 \| -0.546E-08 \| 0.192E-07 \| 0.728E-08 \| -0.413E-09 \| 0.293E-07 \| -0.103E-07 \| 0.118E-07 \| -0.432E-07 \| \| 11 \| -0.562E-08 \| 0.231E-07 \| 0.855E-08 \| -0.449E-09 \| 0.329E-07 \| -0.110E-07 \| 0.128E-07 \| -0.506E-07 \| \| 12 \| -0.575E-08 \| 0.274E-07 \| 0.993E-08 \| -0.486E-09 \| 0.366E-07 \| -0.116E-07 \| 0.138E-07 \| -0.587E-07 \| \| 13 \| -0.584E-08 \| 0.321E-07 \| 0.114E-07 \| -0.523E-09 \| 0.404E-07 \| -0.123E-07 \| 0.147E-07 \| -0.675E-07 \| \| 14 \| -0.590E-08 \| 0.371E-07 \| 0.131E-07 \| -0.561E-09 \| 0.442E-07 \| -0.129E-07 \| 0.157E-07 \| -0.770E-07 \| \| 15 \| -0.592E-08 \| 0.426E-07 \| 0.148E-07 \| -0.600E-09 \| 0.481E-07 \| -0.135E-07 \| 0.166E-07 \| -0.872E-07 \| \| 16 \| -0.592E-08 \| 0.485E-07 \| 0.166E-07 \| -0.641E-09 \| 0.520E-07 \| -0.141E-07 \| 0.175E-07 \| -0.981E-07 \| \| 17 \| -0.589E-08 \| 0.549E-07 \| 0.186E-07 \| -0.683E-09 \| 0.560E-07 \| -0.147E-07 \| 0.184E-07 \| -0.110E-06 \| \| 18 \| -0.583E-08 \| 0.617E-07 \| 0.207E-07 \| -0.729E-09 \| 0.601E-07 \| -0.153E-07 \| 0.192E-07 \| -0.122E-06 \| \| 19 \| -0.576E-08 \| 0.690E-07 \| 0.229E-07 \| -0.776E-09 \| 0.642E-07 \| -0.159E-07 \| 0.201E-07 \| -0.135E-06 \| \| 20 \| -0.566E-08 \| 0.769E-07 \| 0.251E-07 \| -0.827E-09 \| 0.683E-07 \| -0.165E-07 \| 0.210E-07 \| -0.149E-06 \| Table 2 : continued For ITEST3= 1 the Mellin inversion for the basic functions $`g_i(x)`$ based on the representations of the Mellin transforms for complex argument are compared to the numerical expressions for the functions $`g_i(x)`$ in the range $`xϵ[10^7,0.99]`$. The value of $`g_i`$ as a function of $`x`$ is calculated. The relative accuracy comparing both calculations RAT = VAL1/VAL2 -1 is given. The inverse Mellin transform to $`x`$-space is performed by a contour integral numerically. The singularities of the Mellin transforms for $`\mathrm{N}ϵ𝐂`$ are situated at the non–positive integers for the coefficient and splitting functions in massless fixed–order perturbation theory. <sup>2</sup><sup>2</sup>2Some all order resummations, as e.g. the small–$`x`$ resummation, lead also to singularities outside the real axis, cf. . The non–perturbative input densities are usually expressed in terms of polynomials of the type $`h(x)={\displaystyle \underset{i}{}}A_ix^{\alpha _i}(1x)^{\beta _i},`$ (172) the Mellin transform of which is a sum of Euler Beta–functions. Their singularities are as well situated on the real axis left of an upper bound. The inverse Mellin transform is given by $`h(x)={\displaystyle _0^{\mathrm{}}}𝑑z\mathrm{𝖨𝗆}\left[e^{i\mathrm{\Phi }}x^{c(z)}f[c(z)]\right]`$ (173) where $`c(z)=c_0+ze^{i\mathrm{\Phi }}`$, cf. also . The parameter integral over $`z`$ can be performed by standard algorithms. The code segments the integration path logarithmically into N=20 or N=50 pieces on each of which the integral is calculated by the 8-point or 32-point Gauss formula. The starting point C is varied according to the rightmost singularity of the function to be inverted and is chosen close to it on the right. The angle of the linear path w.r.t. the real axis is chosen by $`\mathrm{\Phi }=(3/4)\pi `$ for $`x<0.8`$, $`\mathrm{\Phi }=(7/8)\pi `$ for $`x>0.8`$. and $`\mathrm{\Phi }=(19/20)\pi `$ for $`x>0.98`$. As an example the relative accuracy of the representation for the function $`g_{16}=\left[\mathrm{log}(1+x)\mathrm{log}(2)\right]/(x1)\text{Li}_2(x)`$ between $`x=10^7`$ and $`x=0.99`$ is ``` ************************************************** X,RAT16,VAL= .10000D-06 -.44220D-08 -.69315D-07 X,RAT16,VAL= .10000D-05 -.10120D-07 -.69315D-06 X,RAT16,VAL= .10000D-04 -.13908D-07 -.69314D-05 X,RAT16,VAL= .10000D-03 -.16350D-07 -.69310D-04 X,RAT16,VAL= .10000D-02 -.15094D-07 -.69267D-03 X,RAT16,VAL= .10000D-01 .61713D-08 -.68838D-02 X,RAT16,VAL= .50000D-01 .94884D-08 -.33499D-01 X,RAT16,VAL= .10000D+00 -.19742D-07 -.64836D-01 X,RAT16,VAL= .15000D+00 -.92017D-08 -.94219D-01 X,RAT16,VAL= .20000D+00 .14400D-07 -.12183D+00 X,RAT16,VAL= .25000D+00 .20299D-07 -.14783D+00 X,RAT16,VAL= .30000D+00 .50955D-08 -.17236D+00 X,RAT16,VAL= .35000D+00 -.15289D-07 -.19554D+00 X,RAT16,VAL= .40000D+00 -.23563D-07 -.21747D+00 X,RAT16,VAL= .45000D+00 -.13795D-07 -.23827D+00 X,RAT16,VAL= .50000D+00 .61787D-08 -.25800D+00 X,RAT16,VAL= .55000D+00 .21289D-07 -.27676D+00 X,RAT16,VAL= .60000D+00 .19772D-07 -.29461D+00 X,RAT16,VAL= .65000D+00 .16678D-08 -.31161D+00 X,RAT16,VAL= .70000D+00 -.19671D-07 -.32783D+00 X,RAT16,VAL= .75000D+00 .22255D-05 -.34332D+00 X,RAT16,VAL= .80000D+00 .22936D-05 -.35811D+00 X,RAT16,VAL= .85000D+00 .30694D-05 -.37226D+00 X,RAT16,VAL= .90000D+00 .38631D-05 -.38581D+00 X,RAT16,VAL= .95000D+00 -.25871D-02 -.39879D+00 X,RAT16,VAL= .99000D+00 -.20484D-01 -.40879D+00 **************************************************** ``` Similar accuracies are obtained also in the other cases. We would like to mention that for $`i=3`$ and IAPP=2,3, respectively the accuracy is lower by three orders of magnitude up to $`x0.6`$ relative to the results for IAPP=1 and remains less above. The central functions of the code are : FCT$`i`$(N), XCG$`i`$, ACG$`i`$ and FKN$`i`$. The REAL\8 FUNCTIONS FCT$`i`$(N), i=1..26 provide the positive integer moments of the basic functions $`g_i(x)`$, N $``$ 1 obtained by numerical integration using the program DAIND . The integrands are defined in the REAL\8 FUNCTIONS FKT$`i`$(N), i=1..26. A second representation of the positive integer moments of the basic functions is given using the finite harmonic sums in the REAL\8 FUNCTIONS XCG$`i`$(N), i=1..26. The analytic continuations of the Mellin transforms are given in the COMPLEX\16 FUNCTIONS ACG$`i`$(N), i=1..26. The basic functions $`g_i(x)`$ are given by REAL\8 FUNCTION FKN$`i`$(X), i=1..26. One may use the above options without performing the aforementioned comparisons. The corresponding pilot parameter is IRUN, which takes values between 0 and 5. IRUN = 0: implies a test run. IRUN = 1: The Mellin moments of the basic functions are calculated by numerical integration for positive Mellin indices. IRUN = 2: The Mellin moments of the basic functions are calculated by their representation in terms of harmonic sums for positive Mellin indices. IRUN = 3: The Mellin moments of the basic functions are calculated by their representation in terms of the analytic continuations ACG$`i`$ for complex argument at positive integer moments. IRUN = 4: The basic functions $`g_i(x)`$ are calculated for $`xϵ[10^7,0.99]`$ using numerical representations for the corresponding functions including the Nielsen integrals. IRUN = 5: The basic functions $`g_i(x)`$ are calculated for $`xϵ[10^7,0.99]`$ using the Mellin moment inversion by a complex contour integral. ### 6.5 Subsidiary Routines A series of subsidiary routines representing mathematical functions and procedures are contained in the code. #### 6.5.1 $`\psi ^{(𝐤)}(𝐳)`$ SUBROUTINE PSI$`k`$(ZZ,RES), k = 0..3 provides the value of the functions RES = $`\psi ^{(k)}(`$ZZ$`)`$. Both ZZ and RES are COMPLEX\16. #### 6.5.2 $`\beta ^{(𝐤)}(𝐳)`$ SUBROUTINE BET$`k`$(ZZ,RES), k = 0..3 provides the value of the functions RES = $`\beta ^{(k)}(`$ZZ$`)`$, cf. Eq. (7). Both ZZ and RES are COMPLEX\16. #### 6.5.3 $`𝚪(𝐳)`$ SUBROUTINE GAMMA(Z,RES) is the Euler Gamma function RES = $`\mathrm{\Gamma }(`$Z$`)`$ for complex argument, with $$\mathrm{\Gamma }(z+1)=z\mathrm{\Gamma }(z).$$ (174) Z and RES are COMPLEX\16. For $`\|z\|>10,\|\mathrm{arg}\|(z)<\pi `$ the function is represented by $`\mathrm{log}[\mathrm{\Gamma }(z)]`$ $``$ $`z\left[\mathrm{log}(z)1\right]+{\displaystyle \frac{1}{2}}\mathrm{log}(2\pi /z)`$ $`+{\displaystyle \frac{1}{12}}{\displaystyle \frac{1}{z}}{\displaystyle \frac{1}{360}}{\displaystyle \frac{1}{z^3}}+{\displaystyle \frac{1}{1260}}{\displaystyle \frac{1}{z^5}}{\displaystyle \frac{1}{1680}}{\displaystyle \frac{1}{z^7}}{\displaystyle \frac{691}{360360}}{\displaystyle \frac{1}{z^{11}}}+{\displaystyle \frac{1}{156}}{\displaystyle \frac{1}{z^{13}}}{\displaystyle \frac{3617}{122400}}{\displaystyle \frac{1}{z^{15}}}.`$ The SUBROUTINE GAMMAL(Z,RES) delivers RES = $`\mathrm{log}[\mathrm{\Gamma }(`$Z$`)]`$. #### 6.5.4 $`𝐁(𝐚,𝐛)`$ SUBROUTINE BETA(A,B,RES) is the Euler Beta function RES = $`B(`$A,B$`)`$ for complex argument, with $`B(a,b)={\displaystyle \frac{\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)}{\mathrm{\Gamma }(a+b)}}.`$ (176) A,B and RES are COMPLEX\16. This function may be used to represent the Mellin transforms of the partonic input densities for complex argument. #### 6.5.5 Harmonic Sums The code contains routines for the (alternating) harmonic sums up to depth 5 as REAL\8 FUNCTIONS: DOUBLE PRECISION FUNCTION SUM1(I1,N) $`\mathrm{}`$ DOUBLE PRECISION FUNCTION SUM5(I1,I2,I3,I4,I5,N) The input parameters I$`i`$ may be either positive (non–alternating summation) or negative (alternating summation) integers. Example : $`S_{1,1,3}(N)`$ is calculated by SUM3(1,-1,3,N). #### 6.5.6 Polylogarithms The REAL\8 FUNCTIONS DOUBLE PRECISION FUNCTION FLI$`i`$(x), i=2,3,4 calculate the polylogarithms $`\text{Li}_k(x),k=2,3,4`$ for $`xϵ[1,1]`$. #### 6.5.7 Nielsen Integral $`𝐒_{\mathrm{𝟏𝟐}}(𝐱)`$ The Nielsen integral $`S_{1,2}(x)`$, cf. Eq. (4), is calculated by DOUBLE PRECISION FUNCTION S12$`(x)`$. #### 6.5.8 $`𝐈_\mathrm{𝟏}(𝐱)`$ The function $`I_1(x)`$, Eq. (25) is provided by DOUBLE PRECISION FUNCTION YI1$`(x)`$. #### 6.5.9 Numerical Integration REAL\8 FUNCTION DAIND(A,B,FUN,EPS,KEY,MAX,KOU,EST) The function DAIND yields the integral over the function FUN(x) from A to B. The integrand has to be declared as EXTERNAL in the calling routine. The corresponding algorithm was published in Ref. . We recommend to use the integrator setting KEY = 2. Here EPS denotes the demanded relative accuracy of the integral and MAX $``$ 10000 the number of points at which FUN(x) is calculated by this adaptive integration. The output parameters KOU and EST refer to the number of points being used and the estimated accuracy reached. ## 7 Summary We have calculated semi-analytic representations for the analytic continuations of the Mellin transforms of the set of basic functions through which the Wilson coefficients and splitting functions which occur in hard scattering processes in massless field theories as QED and QCD can be expressed up to two–loop order. Here we aimed on high–precision representations which were performed using widely the algebraic and analytic relations of the Nielsen integrals and related functions being considered as well as their Mellin transforms limiting the necessary approximative representations to as few cases as possible. The expressions obtained are compared both to the Mellin moments at positive integer Mellin index with the values obtained by numerical integration and the representations in terms of harmonic sums. The Mellin inversion to $`x`$ space provides a further test on the numerical accuracy of the expressions derived for the analytic continuations of the Mellin transforms for complex values. The FORTRAN–code ANCONT is provided. With the help of these representations the Mellin transforms for all two–loop quantities for the different space– and time–like hard scattering processes in the massless limit can be assembled. A̱cknowledgement. This work was supported in part by EU contract FMRX-CT98-0194 (DG 12 - MIHT).
warning/0003/math0003023.html	ar5iv	text	# Poisson sigma models and symplectic groupoids ## 1. Introduction The notion of a symplectic groupoid , was introduced as part of a program to quantize Poisson manifolds. It is modeled on the following basic construction. Let $`𝔤`$ be a finite dimensional real Lie algebra. Then its dual space $`𝔤^{}`$ carries a Poisson structure, the Kirillov–Kostant structure. It is characterized by the property that the Poisson bracket of linear functions coincides with the Lie bracket of the corresponding elements of $`𝔤`$. Let $`G`$ be any Lie group whose Lie algebra is $`𝔤`$, and let $`T^{}G`$ be its cotangent bundle, with its canonical symplectic structure. Then $`𝔤^{}`$ may be embedded as the cotangent space at the identity, a Lagrangian submanifold of $`T^{}G`$. The Poisson structure on $`𝔤^{}`$ is the one that makes the right-invariant projection $`l:T^{}G𝔤^{}`$ a Poisson map. Then $`T^{}G`$ may be canonically quantized: the algebra of differential operators on $`G`$ is a quantization of the Poisson algebra of functions on $`T^{}G`$ and right-invariant differential operators form a subalgebra which is a quantization of the Poisson algebra of (polynomial) functions on $`𝔤^{}`$. For a general Poisson manifold $`M`$, the program is to embed $`M`$ as a Lagrangian submanifold of a symplectic manifold $`𝒢`$ in such a way that (deformation, geometric, $`\mathrm{})`$ quantization of $`𝒢`$ descends to a quantization of $`M`$. The manifold $`𝒢`$ is supposed to be a symplectic groupoid, an abstraction of the algebraic and geometric properties of $`T^{}G`$. See Sect. 4 for the definition of symplectic groupoids. The difficulties with this program are, on one side, that symplectic groupoids do not always exist as smooth manifolds. On the other side, it does not seem to be completely clear in general how to quantize $`𝒢`$ in such a way that the quantization descends to a quantization of $`M`$. In the meantime, Kontsevich found a different approach to deformation quantization and constructed star products for general Poisson manifolds. In this paper, we show that (with hindsight) the program of deformation quantization based on symplectic groupoids works, albeit in a rather indirect way. For each Poisson manifold $`M`$ we construct a canonical object $`𝒢`$, the phase space of the Poisson sigma model with target space $`M`$. The latter is a classical topological field theory. In its Hamiltonian formulation it is given by an infinite dimensional Hamiltonian system with constraints. The constraints generate Hamiltonian vector fields forming an integrable distribution of tangent subspaces of codimension $`2\mathrm{dim}(M)`$ on the constraint surface. The phase space $`𝒢`$ is then the space of leaves of the corresponding foliation. It carries a natural structure of groupoid, and also of a symplectic groupoid, in the sense that “the axioms would be fulfilled if $`𝒢`$ were a manifold”. The relation with the deformation quantization of $`M`$ is that the same Poisson sigma model, in its perturbative path integral quantization yields Kontsevich’s deformation quantization formula, as was shown in . We study several examples where $`𝒢`$ is a manifold and also an example, suggested by Weinstein, where it is not: the latter is a rotation invariant Poisson structure on $`^3`$ whose symplectic leaves are spheres centered at the origin. If the symplectic area as a function of the radius is not constant but has a critical point, it is known that no symplectic groupoid can exist, since it would contradict a theorem of Dazord . We show how conical singularities of $`𝒢`$ develop in this case. In general, the singularities of $`𝒢`$ stem from the global structure of the foliation. However, if we work in the setting of formal power series, taking a Poisson structure of the form $`ϵ\alpha `$ with $`ϵ`$ a formal parameter, then a symplectic groupoid may be constructed : it is a formal deformation of the cotangent bundle of $`M`$. We also note that our $`𝒢`$ is related to the “local phase space” of Karasev , . His construction is based on first order equations which are essentially our constraint equation with a special choice of gauge, valid near the identity elements of the groupoid. Technically, to work with infinite dimensional manifolds, we use the framework of manifolds modeled on a Banach(able) space, for which we refer to . The plan of the paper is as follows. In Sect. 2 we introduce the Poisson sigma model and recall its relation with deformation quantization. The construction of the phase space of this model is explained in Sect. 3. In Sect. 4 we describe the groupoid structure of the phase space. We then turn to examples. In Sect. 5 the basic examples are presented: in the case of a symplectic manifold $`M`$ our symplectic groupoid is the fundamental groupoid of $`M`$ and in the case of the dual of a Lie algebra it is the cotangent bundle of the connected, simply connected Lie group with the given Lie algebra. In Sect. 6 we examine a counterexample. In the last section we study the case of a two-dimensional Poisson domain and give a construction of a smooth symplectic groupoid in this case. ###### Acknowledgment. We are grateful to L. Tomassini for useful comments, to M. Bordemann for interesting discussions and to A. Weinstein for useful explanations and references to the literature. ## 2. Poisson sigma model Let $`M`$ be a smooth paracompact finite-dimensional manifold. A Poisson structure on $`M`$ is a bivector field $`\alpha C^{\mathrm{}}(M,^2TM)`$ so that $`\{f,g\}=\alpha (\mathrm{d}f,\mathrm{d}g)`$ defines a Lie algebra structure on the space of smooth functions on $`M`$. We assume that a Poisson structure on $`M`$ is given. Let $`\mathrm{\Sigma }`$ be a two-dimensional oriented compact manifold with boundary. We consider an action functional on the space of vector bundle morphisms $`\widehat{X}:T\mathrm{\Sigma }T^{}M`$ from the tangent bundle of $`\mathrm{\Sigma }`$ to the cotangent bundle of $`M`$. Such a map is given by its base map $`X:\mathrm{\Sigma }M`$ and a section $`\eta `$ of $`\mathrm{Hom}(T\mathrm{\Sigma },X^{}(TM))`$: for $`u\mathrm{\Sigma },vT_u\mathrm{\Sigma }`$, $`\widehat{X}(u,v)=(X(u),\eta (u)v)`$. The action functional is defined on maps obeying the boundary condition (2.1) $$\eta (u)v=0,u\mathrm{\Sigma },vT(\mathrm{\Sigma }).$$ Denote by $`,`$ the pairing between the cotangent and tangent space at a point of $`M`$. If $`X`$ is a map from $`\mathrm{\Sigma }`$ to $`M`$, then this pairing induces a pairing between the differential forms on $`\mathrm{\Sigma }`$ with values in the pull-back $`X^{}(T^{}M)`$ and the differential forms on $`\mathrm{\Sigma }`$ with values in $`X^{}TM`$. It is defined as the pairing of the values and the exterior product of differential forms, and takes values in the differential forms on $`\mathrm{\Sigma }`$. Then the action functional is $$S(X,\eta )=_\mathrm{\Sigma }\eta ,\mathrm{d}X+\frac{1}{2}\eta ,(\alpha X)\eta .$$ Here $`\eta `$, $`\mathrm{d}X`$ are viewed as one-forms on $`\mathrm{\Sigma }`$ with values in the pull-back of the (co)tangent bundle and $`\alpha (x)`$ is viewed as a linear map $`T_x^{}MT_xM`$: $`\alpha (x)=\xi _i\zeta _i`$ is identified with the map $`\beta (\xi _i\beta ,\zeta _i\zeta _i\beta ,\xi _i)`$. A natural space of vector bundle morphisms in our setting consists of pairs $`(X,\eta )`$ with $`X`$ continuously differentiable and $`\eta `$ continuous, obeying the boundary condition (2.1). This model was first considered (in the case of surfaces $`\mathrm{\Sigma }`$ without boundary) in ,. The Feynman path integral for this model with $`\mathrm{\Sigma }`$ a disk was studied in : if $`p,q,r`$ are three distinct points on the boundary of the disk, then the semiclassical expansion of the path integral $$fg(x)=_{X(r)=x}f\left(X(p)\right)g\left(X(q)\right)e^{\frac{i}{\mathrm{}}S(X,\eta )}𝑑X𝑑\eta $$ around the critical point $`X(u)=x`$, $`\eta =0`$ gives Kontsevich’s star product formula. This action functional is invariant under an infinite dimensional space of infinitesimal symmetries and the above integral has to be properly gauge fixed. Here we want to study the classical part of this story and formulate the model in the Hamiltonian formalism. ## 3. The phase space of the Poisson sigma model The Hamiltonian formulation of the Poisson sigma model is obtained by taking $`\mathrm{\Sigma }`$ to be a rectangle $`[T,T]\times I`$ with coordinates $`(t,u)`$ (times and space). The action can then be written as $`S=_\mathrm{\Sigma }(\eta _u,_tX+_uX+\alpha \eta _u,\eta _t)dudt`$. The boundary conditions for $`\eta _t`$ are $`\eta _t=0`$ on $`[T,T]\times I`$. According to the rules of Hamiltonian mechanics, the first part of this action defines a symplectic structure on the space of vector bundle morphisms $`TIT^{}M`$ and the coefficient of the Lagrange multiplier $`\eta _t`$ is a system of constraints that generate a distribution of subspaces spanned by Hamiltonian vector fields. The phase space of the Poisson sigma model is then obtained by Hamiltonian reduction, as the set of integral manifolds of this distribution contained in the set of zeros of the constraints. It may also be expressed as Marsden–Weinstein reduction for the symplectic action of an infinite dimensional Lie algebra on (a version of) the cotangent bundle of the path space $`PM`$ of maps $`IM`$. ### 3.1. The cotangent bundle of $`PM`$ Let $`I`$ be the interval $`[0,1]`$ and $`PM`$ the space of continuously differentiable maps $`IM`$. The tangent space at $`XPM`$ is the space of maps $`V:ITM`$ with $`V(u)T_{X(u)}M`$. Let $`T^{}PM`$ be the space of continuous vector bundle morphisms $`(X,\eta ):TIT^{}M`$ with continuously differentiable base map $`X:IM`$. Then $`T^{}PM`$ is a vector bundle over $`PM`$. The fiber $`T_X^{}PM`$ at $`X`$ may be thought of as the space of continuous 1-forms on $`I`$ with values in $`X^{}(T^{}M)`$. The vector bundle $`T^{}PM`$ may be thought of as the cotangent bundle of $`PM`$, since we have a non-degenerate pairing $`(\eta ,V)_0^1\eta (u),V(u)`$ between $`T_X^{}PM`$ and $`T_XPM`$. The canonical symplectic form $`\omega `$ on $`T^{}PM`$ is defined as the differential of the 1-form $`\theta _{(X,\eta )}(V)=_0^1\eta (u),p_{}V(u)`$, $`VT_{(X,\eta )}T^{}PM`$, where $`p:T^{}PMPM`$ is the bundle projection. In local coordinates $`\widehat{X}`$ is described by $`n=\mathrm{dim}(M)`$ functions $`X^iC^1(I)`$ and $`n`$ 1-forms $`\eta _iC^0(I)du`$ on $`I`$. The symplectic form reads then (3.1) $$\omega _{\widehat{X}}(\delta _1\widehat{X},\delta _2\widehat{X})=_0^1\left(\delta _1X^i\delta _2\eta _i\delta _2X^i\delta _1\eta _i\right).$$ We use here and below the Einstein summation convention and do not write the summation signs for sums over repeated indices. ### 3.2. The constraint manifold Let $`𝒞`$ be the space of solutions of the constraint equation (“Gauss law”) (3.2) $$\mathrm{d}X(u)+\alpha (X(u))\eta (u)=0,$$ with $`X`$ continuously differentiable and $`\eta `$ continuous. This space can be made into a smooth manifold modeled on the Banach space $`^n\times C^0(I,^n)`$: pick a Riemannian metric on $`M`$. Then for each $`x_0M`$ and $`\eta _0`$ a continuous 1-form on $`I`$ with values in $`T_{x_0}^{}M`$ sufficiently small, there exists a unique solution of (3.2) such that $`X(0)=x_0`$ and $`\eta (u)`$ is obtained from $`\eta _0(u)`$ by parallel translation for the Levi-Civita connection along the path $`X`$. All solutions of (3.2) may be obtained this way. Thus, upon choosing local coordinates on a neighborhood $`UM`$ of a point and an orthonormal basis in each tangent space, we have a chart $`𝒞𝒰^n\times C^0(I,^n)`$. ### 3.3. An integrable distribution of subspaces Let $`\widehat{X}=(X,\eta )`$ be a vector bundle morphism $`TIT^{}M`$ and suppose $`\beta `$ is a continuously differentiable function $`IT^{}M`$ such that $`\beta (u)T_{X(u)}^{}M`$, $`uI`$ and $`\beta (0)=\beta (1)=0`$. In other words, $`\beta `$ is in the Banach space $`C_0^1(I,X^{}(T^{}M))`$ of $`C^1`$ sections of the pull-back bundle $`X^{}(T^{}M)`$, vanishing at the endpoints. Let $$H_\beta =_0^1\mathrm{d}X+\alpha \eta ,\beta .$$ If we vary $`\widehat{X}`$ in some open set and let $`\beta `$ depend on $`\widehat{X}`$ then $`H_\beta `$ defines a Hamiltonian vector field $`\xi _\beta `$ (“the infinitesimal gauge transformation with gauge parameter $`\beta `$”) on this open set by the rule $$\iota _{\xi _\beta }\omega =\mathrm{d}H_\beta .$$ Here $`\iota `$ denotes interior multiplication. This rule makes sense if the dependence of $`\beta `$ on $`\widehat{X}`$ is such that $`\mathrm{d}H_\beta `$ is in the image of $`\omega `$. We show below a way to extend any given $`\beta C_0^1(I,X^{}(T^{}M))`$ in such a way that this holds. If $`\widehat{X}𝒞`$, then $`H_\beta `$ vanishes and the value of $`\mathrm{d}H_\beta `$ at $`\widehat{X}`$ only depends on $`\beta `$ at $`\widehat{X}`$. Therefore we have for each solution $`\widehat{X}`$ of (3.2) a subspace of the tangent space to the space of vector bundle morphisms $`TIT^{}M`$ at $`\widehat{X}`$ spanned by the vectors $`\xi _\beta `$, $`\beta C_0^1(I,X^{}(T^{}M))`$. A formula for $`\xi _\beta `$ is the following. Let $`^{TM}`$ be a torsion-free connection on $`TM`$. This connection induces connections $`^{T^{}M}`$, $`^{^2TM}`$, $`^{X^{}(T^{}M)}`$ on the vector bundles $`T^{}M`$, $`^2TM`$ over $`M`$, and $`X^{}(T^{}M)`$ over $`I`$, respectively. For $`xM`$, let $`pT_x^{}M`$, $`h_{(x,p)}`$ denote the horizontal lift homomorphism $`T_x^{}MT_{(x,p)}(T^{}M)`$. It maps the tangent vector to a curve $`\gamma `$ through $`x`$ to the tangent vector of the curve $`\widehat{\gamma }`$ through $`(x,p)`$ obeying the geodesic equation $`^{\gamma ^{}(T^{}M)}\widehat{\gamma }=0`$. If $`\widehat{X}=(X,\eta )T^{}PM`$, then $`\xi _\beta (\widehat{X})`$ is the vector bundle morphism $`TIT(T^{}M)`$ $`\xi _\beta (\widehat{X})(u,v)`$ $`=`$ $`h_{(X(u),\eta (u)v)}\left(\alpha (X(u))\beta (u)\right)+_v^{X^{}(T^{}M)}\beta (u)`$ $`\beta (u),(\alpha )(X(u))\eta (u)v,uI,vT_uI.`$ The last two terms are in $`T_{X(u)}^{}M`$ which is identified with the vertical tangent space at $`(X(u),\eta (u)v)T^{}M`$. If $`\widehat{X}`$ solves (3.2), then this expression is independent of the choice of the connection. It may be more illuminating to write $`\xi _\beta `$ in local coordinates: applying $`\xi _\beta `$ to the coordinate maps $`\widehat{X}X^i(u)`$, $`\widehat{X}\eta _i(u)`$, with respect to some choice of coordinates on $`M`$, gives $`\xi _\beta X^i(u)`$ $`=`$ $`\alpha ^{ij}(X(u))\beta _j(u)`$ $`\xi _\beta \eta _i(u)`$ $`=`$ $`\mathrm{d}_u\beta _i(u)+_i\alpha ^{jk}(X(u))\eta _j(u)\beta _k(u)`$ ###### Theorem 3.1. Let $`\widehat{X}=(X,\eta )𝒞`$. Then the subspace of $`T_{\widehat{X}}PM`$ spanned by $`\xi _\beta `$, $`\beta C_0^1(I,X^{}(T^{}M))`$, is a closed subspace of codimension $`2\mathrm{dim}(M)`$. Proof. For simplicity, we present the proof for $`M`$ a domain in $`^n`$ and work with coordinates. A general tangent vector at a point $`(X,\eta )`$ of $`𝒞`$ is a solution $`(\dot{X},\dot{\eta })`$ of the linearization $$\dot{X}^i(u)+_k\alpha ^{ij}\left(X(u)\right)\dot{X}^k(u)\eta _j(u)+\alpha ^{ij}\left(X(u)\right)\dot{\eta }_j(u)=0,$$ of the constraint equation. With our conditions on differentiability, $`(\dot{X},\dot{\eta })C^1(I,^n)C^0(I,^n)`$, the map $`\beta \xi _\beta `$ is a continuous linear map from the Banach space $`C_0^1(I,^n)`$ to $`C^1(I,^n)C^0(I,^n)`$. It is injective, since $`\xi _\beta =0`$ implies that $`\beta `$ obeys a homogeneous linear first order differential equation with zero initial condition, and thus vanishes identically. Let us describe the image of $`\xi `$. If $`(\dot{X},\dot{\eta })`$ is in the image then $`\dot{X}(0)=0`$ and $`\dot{\eta }`$ is of the form (3.4) $$\dot{\eta }(u)=\mathrm{d}_u\beta (u)+A(u)\beta (u),$$ for some $`\beta C_0^1(I,^n)`$, where $`A(u)`$ is the matrix $`(_i\alpha ^{kj}\eta _j)_{i,j=1\mathrm{},n}`$. If $`V(u)`$ is the solution of $`d_uV(u)=V(u)A(u)`$ with $`V(0)=1`$, then (3.4) reads $`V(u)\dot{\eta }(u)=\mathrm{d}_u\left(V(u)\beta (u)\right)`$. Since $`\beta (u)`$ vanishes at the endpoints, we see that $`_IV(u)\dot{\eta }(u)=0`$. Conversely, if $`(\dot{X},\dot{\eta })`$ obey (3.5) $$\dot{X}(0)=0,_IV(u)\dot{\eta }(u)=0,$$ then $`(\dot{X},\dot{\eta })=\xi _\beta `$, with $`\beta (u)=V(u)^1_0^uV(u^{})\eta (u^{})`$. The image is thus described as the common kernel (3.5) of $`2n`$ linearly independent continuous linear functions, and is thus closed of codimension $`2n`$. $`\mathrm{}`$ The next step is to show that the distribution of subspaces in the tangent bundle to the space of solutions of $`(\text{3.2})`$ is integrable and thus defines a foliation of codimension $`2\mathrm{dim}(M)`$. This is best seen by interpreting the leaves as orbits of a gauge group which we introduce in the next section. ### 3.4. The Lie algebra and its action on the cotangent bundle The Lie algebra acting on $`T^{}PM`$ is obtained from the Lie algebra of 1-forms $`\mathrm{\Omega }^1(M)`$ with the Koszul Lie bracket. This bracket is defined by $$[\beta ,\gamma ]=\mathrm{d}\beta ,\alpha \gamma \iota _{\alpha \beta }\mathrm{d}\gamma +\iota _{\alpha \gamma }\mathrm{d}\beta ,$$ for any $`\beta ,\gamma \mathrm{\Omega }^1(M)`$. In local coordinates $`\alpha =\frac{1}{2}\alpha ^{ij}\frac{}{x^i}\frac{}{x^j}`$, $`\beta =beta_i\mathrm{d}x^i`$, $`\gamma =\gamma _i\mathrm{d}x^i`$, with $`_i=/x^i`$, $$[\beta ,\gamma ]=\left(_i\alpha ^{jk}\beta _j\gamma _k+\alpha ^{jk}_j\beta _i\gamma _k+\alpha ^{jk}\beta _j_k\gamma _i\right)\mathrm{d}x^i.$$ This bracket obeys the Jacobi identity as a consequence of the Jacobi identity for $`\alpha `$. Let $`P_0\mathrm{\Omega }^1(M)`$ be the Lie algebra of continuously differentiable maps $`I\mathrm{\Omega }^1(M)`$ such that $`\beta (0)=\beta (1)=0`$, with bracket $`[\beta ,\gamma ](u)=[\beta (u),\gamma (u)]`$. If $`\beta P_0\mathrm{\Omega }^1(M)`$, let $$H_\beta (X,\eta )=_I\mathrm{d}X(u)+\alpha (X(u))\eta (u),\beta (X(u),u).$$ Recall that if $`H`$ is a smooth function on a symplectic manifold, then a vector field $`\xi `$ is called Hamiltonian vector field generated by $`H`$ if $`\iota _\xi \omega =\mathrm{d}H`$. Such a vector field, it it exists, is unique. In the infinite dimensional setting existence is not guaranteed in general. ###### Theorem 3.2. 1. For each $`\beta P_0\mathrm{\Omega }^1(M)`$ there exists a Hamiltonian vector field $`\xi _\beta `$ generated by $`H_\beta `$. 2. The Lie algebra $`P_0\mathrm{\Omega }^1(M)`$ acts on $`T^{}PM`$ by the Hamiltonian vector fields $`\xi _\beta `$, i.e., $`\beta \xi _\beta `$ is a Lie algebra homomorphism. 3. The map $`\mu :T^{}PMP_0\mathrm{\Omega }^1(M)^{}`$ with $`\mu (X,\eta ),\beta =H_\beta (X,\eta )`$ is an equivariant moment map for this action. Proof. (i) By using a partition of unity, we may restrict ourselves to $`\beta `$ with support in a coordinate neighborhood of $`M`$, and use local coordinates. If $`\dot{X}^i(u),\dot{\eta }_i(u)`$ are the coordinates of a vector field $`\zeta `$ on $`T^{}PM`$ then, with the abbreviations $`\beta _i=\beta _i(X(u),u)`$, $`\alpha ^{ij}=\alpha ^{ij}(X(u))`$, $`\mathrm{d}H_\beta (\zeta )`$ $`=`$ $`{\displaystyle _0^1}\left((\mathrm{d}_u\dot{X}^i+_k\alpha ^{ij}\dot{X}^k\eta _j+\alpha ^{ij}\dot{\eta }_j)\beta _i+C^j_i\beta _j\dot{X}^i\right)`$ $`=`$ $`{\displaystyle _0^1}\dot{X}^i\left(\mathrm{d}_u\beta _i+_i\alpha ^{jk}\eta _k\beta _j+C^j_i\beta _j\right)+{\displaystyle _0^1}\dot{\eta }_i\alpha ^{ji}\beta _j.`$ The term with $`C^j=\mathrm{d}_uX^i+\alpha ^{ij}\eta _j`$ vanishes on $`𝒞`$. Here $`\mathrm{d}_u`$ is the (total) differential with respect to the coordinate $`u`$ on the interval. We may then read off the coordinates $`\delta _\beta X^i,\delta _\beta \eta _i`$ of $`\xi _\beta `$, and at the same time show that they exist, from the defining relation $`\omega (\xi _\beta ,\zeta )=\mathrm{d}H_\beta (\zeta )`$, where $`\omega `$ is given by (3.1). We obtain $`\delta _\beta X^i(u)`$ $`=`$ $`\alpha ^{ij}(X(u))\beta _j(X(u),u)`$ $`\delta _\beta \eta _i(u)`$ $`=`$ $`\mathrm{d}_u\beta _i(X(u),u)+_i\alpha ^{jk}(X(u))\eta _j(u)\beta _k(X(u),u)`$ $`C^j(u)_i\beta _j(X(u),u).`$ (ii) is a consequence of (iii) (iii) The statement amounts to the identity $`H_{[\beta ,\gamma ]}=\xi _\beta H_\gamma `$, which we may again check in local coordinates. We have $`\xi _\beta H_\gamma (X,\eta )`$ $`=`$ $`\xi _\beta {\displaystyle _0^1}C^i\gamma _i`$ $`=`$ $`{\displaystyle _0^1}(C^i_k\gamma _i\alpha ^{kl}\beta _l_k\alpha ^{ij}\mathrm{d}_uX^k\beta _j\gamma _i`$ $`_k\alpha ^{ij}\alpha ^{kl}\beta _k\eta _j\gamma _i+\alpha ^{ij}_j\alpha ^{kl}\eta _j\beta _l\alpha ^{ij}C^k_j\beta _k\gamma _i).`$ By combining terms with the Jacobi identity, we arrive at the formula $$\xi _\beta H_\gamma =_0^1C^i(_i\alpha ^{jk}\beta _j\gamma _k+\alpha ^{jk}_j\beta _i\gamma _k+\alpha ^{jk}\beta _j_k\gamma _i)=H_{[\beta ,\gamma ]}.$$ $`\mathrm{}`$ ### 3.5. The phase space The set $`\mu ^1(0)`$ of zeros of the moment map is the constraint manifold $`𝒞`$. One would like to define the phase space as the Marsden–Weinstein symplectic quotient $`T^{}PM//H=𝒞/H`$. The gauge group $`H`$ is the group of symplectic diffeomorphisms generated by the flows of the Hamiltonian vector fields $`\xi _\beta `$. The trouble is that not only the manifold is infinite dimensional, but the action of the group is far from being nice, and one should not expect to have a good quotient. However, locally the orbits form a smooth foliation: ###### Theorem 3.3. The distribution of tangent subspaces of $`𝒞`$ spanned by $`\xi _\beta `$, $`\beta C_0^1(I,X^{}(T^{}M))`$ is integrable. Its integral manifolds are smooth of codimension $`2\mathrm{dim}(M)`$ and are the orbits of $`H`$. Proof. We present the proof in the case where $`M`$ is a domain in $`^n`$. The general case is treated in a similar but more cumbersome way. Let $`V_{(X,\eta )}`$ be the subspace of $`T_{(X,\eta )}𝒞`$. spanned by $`\xi _\beta `$, $`\beta C_0^1(I,X^{}(T^{}M))`$. These vector spaces form a smooth subbundle of the tangent bundle: locally over a neighborhood $`𝒰𝒞`$, this subbundle is the image of the smooth vector bundle morphism $`𝒰\times C_0^1(I,^n)`$ $``$ $`T𝒰`$ $`((X,\eta ),\beta )`$ $``$ $`((X,\eta ),\xi _\beta (X,\eta )).`$ By Theorem 3.1, in each fiber this is an injective linear continuous map with closed image of codimension $`2n`$. Now the integrability follows from the Frobenius theorem (see , Chapter VI, for a proof valid in the infinite dimensional setting): every $`\beta C_0^1(I,X^{}(T^{}M))`$ may be extended to an element of $`P_0\mathrm{\Omega }^1(M)`$: in coordinates, take $`\beta _i(x,u)`$ independent of $`x`$. It then follows from Theorem 3.2 (ii), that $`[\xi _\beta ,\xi _\gamma ]=\xi _{[\beta ,\gamma ]}`$, which implies the Frobenius integrability criterion. The fact that the integral manifolds are orbits of $`H`$, follows from the fact that $`V_{(X,\eta )}`$ coincides with the space spanned by the restriction to $`(X,\eta )`$ of Hamiltonian vector fields generated by $`H_\beta `$, $`\beta P_0\mathrm{\Omega }^1(M)`$. $`\mathrm{}`$ ## 4. The symplectic groupoid structure on $`T^{}PM//H`$ A symplectic manifold $`𝒢`$ with symplectic form $`\omega _𝒢`$ is called symplectic groupoid for a Poisson manifold $`M`$ if we have an injection $`j:M𝒢`$, two surjections $`l,r:𝒢M`$, a composition law $`g,hg{}_{}{}^{_{}}h`$ defined if $`g,h𝒢`$ and $`r(g)=l(h)`$ obeying a set of axioms. The first axioms say that $`𝒢`$ is a groupoid, i.e., denoting $`𝒢_{x,y}=l^1(x)r^1(y)`$, 1. $`lj=rj=\mathrm{id}_M`$. 2. If $`g𝒢_{x,y}`$ and $`h𝒢_{y,z}`$, then $`g{}_{}{}^{_{}}h𝒢_{x,z}`$. 3. $`j(x){}_{}{}^{_{}}g=g{}_{}{}^{_{}}j(y)=g`$, if $`g𝒢_{x,y}`$. 4. To each $`g𝒢_{x,y}`$ there exists an inverse $`g^1𝒢_{y,x}`$ such that $`g{}_{}{}^{_{}}g_{}^{1}=j(x)`$. 5. The composition law is associative: $`(g{}_{}{}^{_{}}h){}_{}{}^{_{}}k=g{}_{}{}^{_{}}(h{}_{}{}^{_{}}k)`$ whenever defined. In the language of categories, these axioms say that $`𝒢`$ is the set of morphisms of a category in which all morphisms are isomorphisms. $`M`$ is the set of objects and $`j(M)`$ the set of identity morphisms. It follows from the axioms that $`g^1`$ is uniquely determined by $`g`$ and that $`g^1{}_{}{}^{_{}}g=j(y)`$, if $`g𝒢_{x,y}`$. The next axioms relate to the symplectic and Poisson structure. A smooth map $`\varphi :M_1M_2`$ between Poisson manifolds is called Poisson if $`\{f,g\}\varphi =\{f\varphi ,g\varphi \}`$, for all $`f,gC^{\mathrm{}}(M_2)`$. It is called anti-Poisson if $`\{f,g\}\varphi =\{f\varphi ,g\varphi \}`$. Then the remaining axioms are: 1. $`j`$ is a smooth embedding, $`l,r`$ are smooth submersions, the composition and inverse maps are smooth. 2. $`j(M)`$ is a Lagrangian submanifold. In particular $`\mathrm{dim}(𝒢)=2\mathrm{dim}(M)`$. 3. $`l`$ is a Poisson map and $`r`$ is an anti-Poisson map. 4. Let $`P:𝒢_0𝒢\times 𝒢𝒢`$ be the composition law on $`𝒢_0=\{(g,h)𝒢\|r(g)=l(h)\}`$, and $`\pi _1,\varphi _2:𝒢\times 𝒢𝒢`$ denote the projections onto the first and second factor. Then $`P^{}\omega _𝒢=\pi _1^{}\omega _𝒢+\pi _2^{}\omega _𝒢`$. 5. $`gg^1`$ is an anti-Poisson map. The basic example is the following: ###### Example 4.1. Let $`M=𝔤^{}`$ be the dual space to a Lie algebra $`𝔤`$ with Kirillov–Kostant Poisson structure. For any Lie group $`G`$ with Lie algebra isomorphic to $`𝔤`$, we have the inclusion $`j:𝔤^{}T^{}G`$ of $`𝔤^{}`$ as the cotangent space at the identity $`e`$ and projections $`l,r:T^{}G𝔤^{}`$ sending the cotangent space at each point to the cotangent space at the identity by left (right) translation. If $`L_g,R_g:GG`$, with $`L_g(h)=gh`$, $`R_g(h)=hg`$, denote the left and right translation by $`g`$, we have $`l(g,a)=\mathrm{d}R_g(e)^{}a`$, $`r(g,a)=\mathrm{d}L_g(e)^{}b`$, ($`gG`$, $`aT_g^{}G`$). If $`r(g,a)=l(h,b)`$, the composition law is $`(g,a){}_{}{}^{_{}}(h,b)=(gh,c)`$ with $`c=(\mathrm{d}R_h(g)^{})^1a=(\mathrm{d}L_g(h)^{})^1b`$. A more explicit description is obtained by identifying $`T^{}G`$ with $`𝔤^{}\times G`$ via $`(g,a)(\mathrm{d}R_g(e)^{}a,g)`$, see 5.4 below. ### 4.1. The groupoid structure The algebraic groupoid structure of $`𝒢=𝒞/H`$ can be naturally defined in terms of composition of paths. We have an inclusion $`j:M𝒢`$ sending a point $`x`$ to the class of the constant solution $`X(u)=x,\eta (u)=0`$. Let $`l,r:T^{}PMM`$ be the maps (4.1) $$l(X,\eta )=X(0),r(X,\eta )=X(1).$$ These maps are $`H`$-invariant, since the symmetries preserve the endpoints, hence they descend to maps $`l,r:𝒢M`$, and it is clear that axiom (i) holds. ### 4.2. Composition law and inverses To define the composition law we need to do some adjustments at the endpoints: Let $`H_0`$ be the subgroup of $`H`$ generated by the flows of the vector fields $`\xi _\beta `$ such that $`\mathrm{d}\beta (0)=\mathrm{d}\beta (1)=0`$. ###### Lemma 4.2. In each equivalence class $`[(X,\eta )]`$ in $`𝒢=𝒞/H`$ there exists a representative with $`\eta (0)=\eta (1)=0`$. Any two representatives with this property can be related by an element of $`H_0`$. Proof. Let $`(X,\eta )𝒞`$. To obtain a representative with $`\eta (0)=0`$ we perform a gauge transformation obtained as the flow of a vector field $`\xi _\beta `$ with $`\beta `$ supported on a small neighborhood $`I_0`$ of $`0I`$. Small means here that $`X(u)`$ lies in a coordinate neighborhood $`U`$ of $`M`$ for $`uI_0`$. Then the gauge transformation may be described in local coordinates. The problem is then to find continuously differentiable functions $`\beta _i(u)`$ supported on $`I_0`$ with $`\beta _i(0)=0`$, so that the solution to the system $`{\displaystyle \frac{}{s}}X^i(u,s)`$ $`=`$ $`\alpha ^{ij}\left(X(u,s)\right)\beta _j(u),`$ $`{\displaystyle \frac{}{s}}\eta _i(u,s)`$ $`=`$ $`\mathrm{d}_u\beta _i(u)+_i\alpha ^{jk}\left(X(u,s)\right)\eta _j(u,s)\beta _k(u),`$ with initial condition $`(X,\eta )`$ at $`s=0`$ (a) exists with $`X`$ in $`U`$ for all $`s[0,1]`$, and (b) obeys $`\eta _i(0,1)=0`$. A sufficient condition for (a) is that $`\|\beta _i(u)\|<\delta `$ for some $`\delta >0`$ and all $`uI_0`$: if this bound holds with sufficiently small $`\delta `$, then the first equation has a solution which is close to $`X`$ and thus remains in $`U`$. Given $`X(u,s)`$, the second equation is linear for $`\eta `$ and thus has a solution for all $`s[0,1]`$. To achieve (b), let $`a_i=\eta _i(0)`$ and choose $`\beta _i(u)`$ so that $`\beta _i(u)=a_iu+O(u^2)`$. Then $`\eta _i(0,s)=(1s)a_i`$ vanishes for $`s=1`$. The same procedure may be applied at the other end of $`I`$. Suppose now that $`\widehat{X}^{(0)}=(X^{(0)},\eta ^{(0)})`$ and $`\widehat{X}^{(1)}=(X^{(1)},\eta ^{(1)})`$ are two representatives of a class in $`𝒢`$, obeying the condition $`\eta ^{(0)}(u)=\eta ^{(1)}(u)=0`$ for $`u=0,1`$. These representatives are related by an element of $`H`$, which is a product of a finite number $`k`$ of flows of vector fields of the form $`\xi _\gamma `$, $`\gamma P_0\mathrm{\Omega }^1(M)`$. Let us first assume that $`k=1`$. Then we have a smooth path $`s\widehat{X}_s`$ in $`𝒞`$, such that $`\widehat{X}_{s=0}=\widehat{X}^{(0)}`$, $`\widehat{X}_{s=1}=\widehat{X}^{(1)}`$, and $`\mathrm{d}\widehat{X}_s/\mathrm{d}s=\xi _\gamma (\widehat{X}_s)`$. We now repeat the procedure of the first part of the proof, for each $`s[0,1]`$. Let $`x_0=\widehat{X}_s(0)`$, which is independent of $`s`$, and choose coordinates in a neighborhood $`UM`$ of $`x_0`$. Let $`\beta _sP_0\mathrm{\Omega }^1(U)`$ be such that (a) $`\beta _{s,i}(x,u)=a_{s,i}u+O(u^2)`$ $`(u0`$), where $`\eta _{s,i}(0)=a_{s,i}\mathrm{d}u`$; (b) $`\beta _s(x,u)=0`$ if $`u\delta ^{}`$, for some sufficiently small $`\delta ^{}>0`$, (c)$`\mathrm{max}\|\beta _{s,i}\|`$ is sufficiently small. Then the flow of $`\xi _{\beta _s}`$ exists on $`𝒰=\{\widehat{X}𝒰\|\widehat{X}(u)U,`$ for $`u\delta ^{}\}`$ for all times in $`[0,1]`$. Applying this flow to $`\widehat{X}_s`$ we obtain a two-parameter family $`\widehat{X}_{s,\sigma }`$, $`(s,\sigma )[0,1]^2`$, in $`𝒞`$, differing from $`\widehat{X}_s`$ only in some small neighborhood of $`0I`$, such that $`\widehat{X}_{s,0}=\widehat{X}_s`$ and $`\widehat{X}_{s,\sigma }/\sigma =\xi _{\beta _s}(\widehat{X}_{s,\sigma })`$. By construction, we have 1. $`\widehat{X}_{0,\sigma }=\widehat{X}^{(0)}`$, $`\widehat{X}_{1,\sigma }=\widehat{X}^{(1)}`$, for all $`\sigma [0,1]`$. 2. $`\widehat{X}_{s,1}=(X_{s,1},\eta _{s,1})`$ with $`\eta _{s,1}(0)=0`$. Since $`\xi _{\beta _s}`$, $`\xi _\gamma `$ belong to an integrable distribution of tangent subspaces, any curve $`t\widehat{X}_{s(t),\sigma (t)}`$ is in the integral manifold passing through $`\widehat{X}^{(0)}`$. In particular, $`s\widehat{X}_{s,1}`$ defines by (i), (ii) a curve of points related by a transformation of $`H`$. Since $`\eta _{s,1}(0)=0`$, this transformation is in $`H_0`$. The same argument applies to the other endpoint, and for $`k1`$ one applies this construction $`k`$ times. $`\mathrm{}`$ Then we may define the composition law $`[(X,\eta )]=[(X_1,\eta _1)]{}_{}{}^{_{}}[(X_2,\eta _2)]`$ in $`𝒢`$ by choosing representatives as in Lemma 4.2 and setting (4.2) $$\begin{array}{cc}\hfill X(u)& =\{\begin{array}{cc}\hfill X_1(2u),& 0u\frac{1}{2},\hfill \\ \hfill X_2(2u1),& \frac{1}{2}u1,\hfill \end{array}\hfill \\ \hfill \eta (u)& =\{\begin{array}{cc}\hfill 2\eta _{1u}(2u)\mathrm{d}u,& 0u\frac{1}{2},\hfill \\ \hfill 2\eta _{2u}(2u1)\mathrm{d}u,& \frac{1}{2}u1.\hfill \end{array}\hfill \end{array}$$ (we write $`\eta _i(u)=\eta _{iu}(u)\mathrm{d}u`$, $`i=1,2`$), provided $`X_1(1)=X_2(0)`$. Lemma 4.2 ensures that $`\eta `$ is continuous. $`X`$ is continuously differentiable since the derivatives of $`X_1`$, $`X_2`$ at the endpoints match — they vanish, as a consequence of (3.2). It is immediate to check that $`(X,\eta )`$ obeys the constraint equation if $`(X_1,\eta _1)`$, $`(X_2,\eta _2)`$ do. Therefore the composition is well-defined at the level of representatives. By the second part of Lemma 4.2, the class of $`(X,\eta )`$ is independent of the choice of representatives: infinitesimal transformations of $`(X_1,\eta _1)`$, $`(X_2,\eta _2)`$ associated to $`\beta _1,\beta _2:IT^{}M`$ and obeying $`\mathrm{d}\beta _1(1)=\mathrm{d}\beta _2(0)`$ match at the end points to give to give an infinitesimal transformation of $`(X,\eta )`$ associated to $$\beta (u)=\{\begin{array}{cc}\hfill \beta _1(2u),& 0u\frac{1}{2},\hfill \\ \hfill \beta _2(2u1),& \frac{1}{2}u1,\hfill \end{array}$$ which is a differentiable function $`\beta :IT^{}M`$ with $`\beta (u)T_{X(u)}^{}M`$. For $`u[0,1]`$, let $`\theta (u)=1u`$. If $`\widehat{X}=(X,\eta )`$ obeys the constraint equation (3.2), then $`\theta ^{}\widehat{X}=(X\theta ,\theta ^{}\eta )`$ also does. Moreover the endpoints of the path $`X`$ are interchanged under this map. If $`\beta `$ is a section of $`X^{}(T^{}M)`$ then $`\beta \theta `$ is a section of $`(X\theta )^{}(T^{}M)`$, and $`\theta ^{}(\xi _\beta \widehat{X})=\xi _{\beta \theta }\theta ^{}\widehat{X}`$. Therefore $$[(X,\eta )][(X,\eta )]^1=[(X\theta ,\theta ^{}\eta )]$$ is a well-defined map from $`𝒢`$ to $`𝒢`$. ###### Theorem 4.3. $`𝒢`$ obeys axioms (i)–(v) The idea of the proof of the associativity is based on the fact that the composition law is associative up to reparametrization of $`I`$. But it turns out that reparametrizations are special gauge transformations: indeed, if an infinitesimal reparametrization is given by a vector field $`ϵ`$ on $`I`$ vanishing at the endpoints, then the variation of a solution $`\widehat{X}`$ of the constraint equation is an infinitesimal transformation with parameter $`\beta (u)=\eta (u)ϵ(u)`$, provided $`\eta `$ is differentiable. Similarly, to prove that $`g{}_{}{}^{_{}}g_{}^{1}=j(x)`$ one uses the fact that $`g{}_{}{}^{_{}}g_{}^{1}`$ is the class of a point $`(X,\eta )𝒞`$ so that $`(X\theta ,\theta ^{}\eta )=(X,\eta )`$. Let $$\beta (u)=\beta _{\widehat{X}}(u)=\{\begin{array}{cc}\hfill u\eta _u(u),& 0u\frac{1}{2},\hfill \\ \hfill (1u)\eta _u(u),& \frac{1}{2}u1.\hfill \end{array}$$ Here again, we have to assume that $`\eta `$ is differentiable. Then the solution of $`\frac{\mathrm{d}\widehat{X}_s}{\mathrm{d}s}=\xi _\beta `$ with initial condition $`\widehat{X}_0=\widehat{X}`$ is a solution $`\widehat{X}_s`$ of (3.2) which is also symmetric with respect to $`\theta `$. For $`s=1`$, $`X_1(u)=x`$ is constant. Thus $`\eta `$ is in the kernel of $`\alpha (x)`$ and obeys $`\theta ^{}\eta =\eta `$. But Ker$`\alpha (x)`$ is naturally a Lie algebra (the Lie bracket $`[\beta ,\gamma ]`$ is the Koszul bracket for any two 1-forms coinciding with $`\beta `$, $`\gamma `$ at $`x`$). Then $`\mathrm{d}+\eta `$ has the interpretation of a connection on a trivial vector bundle over $`I`$. Infinitesimal gauge transformations preserving the condition $`X(u)=x`$ are infinitesimal gauge transformations in the usual gauge theory sense. In particular a connection with $`\theta ^{}\eta =\eta `$ is gauge equivalent (with a gauge transformation which is trivial at the endpoints) to the trivial connection. Technically, these operations are possible thanks to the ###### Lemma 4.4. In every class $`[(X,\eta )]𝒢`$ there exists a representative so that $`X`$ and $`\eta `$ are smooth maps. Proof. Let $`(X,\eta )𝒞`$. Let us divide the interval $`I`$ into an odd number $`k3`$ of parts $`I_j=[j/k,(j+1)/k]`$, $`0jk1`$, in such a way that $`X(I_j)`$ is contained in a coordinate neighborhood of $`M`$. To find a smooth representative, we perform a sequence of gauge transformations. Each of these gauge transformation is generated by parameters $`\beta `$ with support in a small neighborhood of an interval $`I_j`$. Such gauge transformations only affect $`(X,\eta )`$ in such a neighborhood. Therefore we may describe them in local coordinates by the formula (3.3). Let $`\eta ^{\mathrm{smooth}}C^{\mathrm{}}(I,X^{}(T^{}M))\mathrm{\Omega }^1(I)`$ be a smooth section, $`C^0`$-close to $`\eta `$. As a first step, we show that $`\eta `$ may be taken to be equal to $`\eta ^{\mathrm{smooth}}`$ on $`I_0`$. Let for $`s[0,1]`$, $`\eta _i(u,s)=s\eta _{}^{\mathrm{smooth}}{}_{i}{}^{}(u)+(1s)\eta _i(u)`$, $`1i\mathrm{dim}(M)`$, $`uI_0`$. Let $`X(u,s)`$ be the solution of the constraint equation on $`I_0`$ $$\frac{}{u}X^i(u,s)+\alpha ^{ij}(X(u,s))\eta _j(u,s)=0,$$ with $`X^i(0,s)=X^i(0)`$. This equations has a unique solution on $`I_0`$ if $`\eta ^{\mathrm{smooth}}`$ is sufficiently close to $`\eta `$. Let $`\beta `$ be the solution of the linear differential equation (4.3) $$\frac{}{u}\beta _i(u,s)+_i\alpha ^{jk}(X(u,s))\eta _j(u,s)\beta _k(u,s)=\frac{}{s}\eta _i(u,s),$$ on $`I_0`$ with initial condition $`\beta _i(0,s)=0`$. Extend $`\beta _i(u,s)`$ to a function on $`I`$ vanishing outside some small neighborhood of $`I_0`$. Then $`\beta _i(,s)`$ is the local coordinate expression of a section $`\beta _sC_0^1(I,X^{}(T^{}M))`$ with support in a neighborhood of $`I_0`$. It may be taken to be small in the $`C^1`$-topology if $`\eta ^{\mathrm{smooth}}`$ is close to $`\eta `$ in the $`C^0`$-topology. The flow of the vector field $`\xi _{\beta _s}`$, $`0s1`$, is then a gauge transformation that transforms $`(X,\eta )`$ into a solution $`(\stackrel{~}{X},\stackrel{~}{\eta })`$ coinciding with $`(X,\eta )`$ outside a neighborhood of $`I_0`$ and such that $`\stackrel{~}{\eta }=\eta ^{\mathrm{smooth}}`$ on $`I_0`$. This step may be repeated on $`I_1`$, $`I_2`$, …, until we get a representative $`(\stackrel{~}{X},\stackrel{~}{\eta })`$ with $`\stackrel{~}{\eta }=\eta ^{\mathrm{smooth}}`$ on $`[0,(k+1)/2k]`$. We then repeat the same step integrating (4.3) backwards, starting from the last interval $`I_{k1}`$, and continuing with $`I_{k2},\mathrm{}`$, until we arrive at the middle interval $`I_{(k1)/2}`$. At this point the representative $`(\stackrel{~}{X},\stackrel{~}{\eta })`$ has $`\stackrel{~}{\eta }=\eta ^{\mathrm{smooth}}`$ except on some small interval $`I^{}`$ in the middle of $`I`$. We apply once more our step to a slightly bigger interval $`I^{\prime \prime }`$ including $`I^{}`$. Then the solution $`\beta `$ of (4.3) is a smooth function of $`u`$ on $`I^{\prime \prime }I^{}`$ and may be extended to a section in $`C_0^1(I,X^{}(T^{}M))`$ which is smooth outside $`I^{}`$. The resulting representative $`(\overline{X},\overline{\eta })`$ of the class $`[(X,\eta )]`$ has then $`\overline{\eta }`$ smooth. Then also $`\overline{X}`$, as a solution of (3.2), is smooth. $`\mathrm{}`$ ### 4.3. Symplectic structure To formulate axioms (vi)–(x) we need $`𝒢`$ to be a manifold, which is not always the case, as we shall see below. So we assume that $`𝒢`$ is a manifold, or more precisely: ###### Assumption 4.5. There exists a smooth manifold $`𝒢`$ and a smooth submersion $`\pi :𝒞𝒢`$ whose fibers are the $`H`$-orbits. Below we give examples where this assumption holds and examples where it does not. The symplectic structure $`\omega _𝒢`$ on $`𝒢`$ is constructed in the usual way: the point is that the symplectic 2-form $`\omega `$ on $`T^{}PM`$ restricts to an $`H`$-invariant closed 2-form on $`𝒞`$ whose null spaces are the tangent spaces to the orbits. This implies that $$\omega _𝒢(x)(\xi ,\zeta )=\omega (\widehat{X})(\widehat{\xi },\widehat{\zeta }),\xi ,\zeta T_x𝒢,$$ is independent of the choice of $`\widehat{X}𝒞`$ such that $`\pi (\widehat{X})=x`$ or of $`\widehat{\xi },\widehat{\zeta }T_{\widehat{X}}𝒞`$ projecting to $`\xi ,\zeta `$, and defines a symplectic 2-form on $`𝒢`$. ###### Theorem 4.6. Under Assumption 4.5, $`𝒢`$ is a symplectic groupoid for $`M`$. Proof. The non-trivial assertions are (viii)-(x). Let us prove that the left projection $`l`$ is a Poisson map. Let $`𝒰`$ be some small neighborhood in $`T^{}PM`$ of a point $`\widehat{X}_0𝒞`$, and choose local coordinates on $`M`$ around $`X_0(0)`$. Then it is sufficient to show that the coordinates $`l^i`$ of $`l`$ obey $`\{l^i,l^j\}_𝒢=\alpha ^{ij}l`$. Let $`\psi (u)\mathrm{d}u`$ be any smooth 1-form on $`I`$ with support in a small neighborhood of $`0`$ and such that $`_0^1\psi (u)du=1`$. Then the function $`𝒰`$ $$l_\psi ^i:\widehat{X}_0^1\left(X^i(u)+_0^u\alpha ^{ij}\left(X(v)\right)\eta _j(v)\right)\psi (u)du$$ restrict to $`l^i`$ on $`𝒰𝒞`$ (with the support condition on $`\psi `$, this local coordinate expression makes sense). The main property of this extension of $`l^i`$ is that its differential lies in the image of $`\omega `$ and thus generates a local Hamiltonian vector field $`\xi ^i`$. Therefore we may compute $`\{l^i,l^j\}_𝒢`$ as the Poisson bracket $`\{l_\psi ^i,l_\psi ^j\}`$ on $`T^{}PM`$, which is then $`\xi ^il^j`$. The differential of $`l_\psi ^i`$ applied to a vector field $`\zeta `$ with coordinates $`\dot{X}^j,\dot{\eta }_j`$ is $`\mathrm{d}l_\psi ^i(\zeta )`$ $`=`$ $`{\displaystyle _0^1}\dot{X}^j(u)\left(\delta _{ij}\mathrm{d}u+_j\alpha ^{ik}\left(X(u)\right)\eta _k(u){\displaystyle _u^1}\psi (v)dv\right)`$ $`+{\displaystyle _0^1}\dot{\eta }_j(u)\alpha ^{ij}\left(X(u)\right){\displaystyle _u^1}\psi (v)dv.`$ The vector field $`\xi ^i`$, solution of $`\omega (\xi ^i,\zeta )=\mathrm{d}l^i(\zeta )`$ has then coordinates $`\delta ^iX^j,\delta ^i\eta _j`$ with $$\delta ^iX^j(u)=\alpha ^{ij}\left(X(u)\right)_u^1\psi (v).$$ We do not need $`\delta ^i\eta _j`$. Then $$\{l^i,l^j\}_𝒢([\widehat{X}])=\xi ^il^j(\widehat{X})=\delta ^iX^j(0)=\alpha ^{ij}(l(\widehat{X})),$$ as was to be shown. To prove (ix) we notice that the integral defining $`\omega `$ at the product of two solutions is the sum of the integrals defining $`\omega `$ at the two solutions. Axiom (x) follows from the fact that the inversion changes the sign of the symplectic form, as can be seen by changing $`u`$ to $`1u`$ in the integral defining $`\omega `$. $`\mathrm{}`$ ## 5. Basic examples In this and in the next sections we discuss some examples. To fix the notations, we will always denote by $`u`$ the variable in our space interval $`I=[0,1]`$. When considering a flow generated by symmetries, we will denote the flow parameter by $`s`$. Finally, we will use a prime to indicate derivatives w.r.t. $`u`$ and a dot for derivatives w.r.t. $`s`$. ### 5.1. Trivial Poisson structures Let us consider a manifold $`M`$ with Poisson bivector field $`\alpha =0`$. In this case, the “Gauss law” selects the constant maps $`X:IM`$. Let $`X(u)=\xi M`$ be such a solution. Then the corresponding bundle map $`\widehat{X}`$ is given by $`X`$ and a continuous one-form $`\eta `$ on $`I`$ that takes value in $`T_\xi ^{}M`$. The infinitesimal symmetries are given by $$\delta \eta =\mathrm{d}\beta ,$$ with $`\beta :IT_\xi ^{}M`$, $`\beta \|_I=0`$. If we define $`\pi :=_I\eta T_\xi ^{}M`$, then for a given solution we have the well-defined map $`i:𝒢T^{}M`$, which maps $`\widehat{X}`$ into $`(X(0),\pi )`$. We can invert this mapping by defining $`j:T^{}M𝒢`$ as follows: $`j(g)`$, $`gT^{}M`$, is the constant morphism $`\widehat{X}(u)=g,uI`$. An immediate check shows that $`ij=\mathrm{id}`$. We can also prove that $`ji=\mathrm{id}`$. In fact, let $`\widehat{X}𝒢`$. Then $`\widehat{\stackrel{~}{X}}:=ji(\widehat{X})`$ is a solution with $`\stackrel{~}{X}=X`$, and $`_I\stackrel{~}{\eta }=_I\eta `$. Denoting by $`I^u`$ the path in $`I`$ from the lower boundary till a point $`u`$, we can then define $`\beta (u)=_{I^u}(\stackrel{~}{\eta }\eta )`$, which is an allowed symmetry generator. Next we consider the following path of $`T_X^{}M`$-valued one-forms $$\eta _s:=s\stackrel{~}{\eta }+(1s)\eta ,s[0,1].$$ Finally, we have $$\dot{\eta }_s=\mathrm{d}\beta ,$$ so that $`\widehat{\stackrel{~}{X}}`$ is equivalent to $`\widehat{X}`$. We have then proved the following ###### Theorem 5.1. The phase space $`𝒢`$ for a trivial Poisson structure on $`M`$ is diffeomorphic to $`T^{}M`$. It is well-known that $`T^{}M`$ is a symplectic groupoid for $`M`$. The two projections $`l`$ and $`r`$ coincide with the natural projection $`T^{}MM`$, while the product is given by $$(\xi ,\pi ){}_{}{}^{_{}}(\xi ,\pi ^{})=(\xi ,\pi +\pi ^{}).$$ ### 5.2. The symplectic case Since now the Poisson bivector field is nondegenerate, the Gauss law allows to completely determine the bundle morphism $`\widehat{X}:TIT^{}M`$ in terms of the base map $`X`$: $$\widehat{X}=\alpha ^1(\mathrm{d}X).$$ As for $`X`$, the infinitesimal symmetries are now all infinitesimal diffeomorphisms of the target that fix the endpoints of $`X(I)`$. Thus, the space of solutions modulo symmetries coincides with the fundamental groupoid of $`M`$. In the case when $`M`$ is simply connected we can further identify $`𝒢`$ with $`M\times \overline{M}`$, where $`\overline{M}`$ denotes $`M`$ with opposite symplectic structure. The product is then simply $$(x,y){}_{}{}^{_{}}(y,z)=(y,z).$$ In the general case, a point in $`𝒢`$ is given by a pair of points $`x`$ and $`y`$ in $`M`$ together with a class $`c`$ of homotopic paths with fixed endpoints in $`x`$ and $`y`$. The product is then $$(x,y,c){}_{}{}^{_{}}(y,z,c^{})=(y,z,cc^{}),$$ where $`cc^{}`$ is the class of paths defined by glueing $`c`$ and $`c^{}`$ together. ### 5.3. Constant Poisson structures This example combines the two previous ones. Let us assume that $`M=^n`$ with a constant Poisson structure $`\alpha `$. It is then possible to assume, if necessary after a linear change of coordinates, that $`\alpha `$ has the following block form: $$\alpha ^{I\mu }=\alpha ^{\mu \nu }=0,I=1,\mathrm{},r,\mu ,\nu =r+1,\mathrm{},n,$$ $$(\alpha ^{IJ})\text{ invertible},I,J=1,\mathrm{},r,$$ where $`r`$ is the rank. The the “Gauss law” reads $$\mathrm{d}X^I+\alpha ^{IJ}\eta _J=0,\mathrm{d}X^\mu =0,$$ and the infinitesimal symmetries are $`\delta X^I`$ $`=\alpha ^{IJ}\beta _J,`$ $`\delta X^\mu `$ $`=0,`$ $`\delta \eta _J`$ $`=\mathrm{d}\beta _J,`$ $`\delta \eta _\mu `$ $`=\mathrm{d}\beta _\mu .`$ Thus, we can split $`^n`$ into $`^r`$ with symplectic structure $`(\alpha ^{IJ})^1`$ and $`^{nr}`$ with trivial Poisson structure. By the previous two examples we get then $$𝒢=^r\times \overline{}^r\times T^{}^{nr}$$ with product $$(x,y,\xi ,\pi ){}_{}{}^{_{}}(y,z,\xi ,\pi ^{})=(x,z,\xi ,\pi +\pi ^{}).$$ ### 5.4. The dual of a Lie algebra Let $`𝔤^{}`$ be the dual of a Lie algebra $`𝔤`$ with structure constants in a given basis denoted by $`f_k^{ij}`$. The Kirillov–Kostant Poisson structure on $`𝔤^{}`$ is then given by the bivector field $$\alpha ^{ij}(x)=f_k^{ij}x^k.$$ In this case the Gauss law reads $$\mathrm{d}X^i+f_k^{ij}X^k\eta _j=0,$$ where $`X`$ is a map $`I𝔤^{}`$ and $`\eta \mathrm{\Omega }^1(I,𝔤)`$. Let then $`\beta `$ be a map $`I𝔤`$ that vanishes on the boundary of $`I`$. The infinitesimal symmetries read $`\delta X^i`$ $`=f_k^{ij}X^k\beta _j,`$ $`\delta \eta _i`$ $`=\mathrm{d}\beta _i+f_i^{jk}\eta _j\beta _k.`$ We can rewrite the above equations in a more recognizable form if we consider $`\eta `$ as the connection one-form for a $`G`$-bundle over $`I`$, where $`G`$ is a Lie group whose Lie algebra is $`𝔤`$. The Gauss law becomes (5.1) $$\mathrm{d}_\eta X=0,$$ while the infinitesimal symmetries now read (5.2) $`\delta X`$ $`=\mathrm{ad}_\beta ^{}X,`$ $`\delta \eta `$ $`=\mathrm{d}_\eta \beta ,`$ that is, $`\beta `$ is an infinitesimal gauge transformation. We define $`𝒢`$ as the space of solutions of (5.1) modulo gauge transformations connected to the identity. We have then the following ###### Theorem 5.2. The phase space $`𝒢`$ is diffeomorphic to $`T^{}G`$, where $`G`$ is the connected, simply connected Lie group whose Lie algebra is $`𝔤`$. The symplectic groupoid structure on $`T^{}G`$ is the one described in Example 4.1. #### 5.4.1. Proof of Theorem 5.2 We first recall that $`T^{}G`$ is isomorphic to $`𝔤^{}\times G`$. We then define a map $$i:\begin{array}{ccc}𝒢& & 𝔤^{}\times G\\ (X,\eta )& & (X(0),\mathrm{Hol}(\eta ))\end{array}$$ where $`0`$ denotes the lower boundary of $`I`$, and $`\mathrm{Hol}(\eta )`$ is the parallel transport from the lower to the upper boundary of $`I`$. Next we want to define an inverse to $`i`$. Let us then consider $`(\xi ,g)𝔤^{}\times G`$. Since $`G`$ is connected, there is a path $`h:IG`$ from the identity to $`g`$. For such a path, we define $$\eta _{[h]}:=h\mathrm{d}h^1.$$ We then define $`X_{\xi ,[h]}`$ as the solution of (5.1) with initial condition $`X_{\xi ,[h]}(0)=\xi `$ determined by $`\eta _{[h]}`$. More precisely, $$X_{\xi ,[h]}=\mathrm{Ad}_{h^1}^{}\xi .$$ So $`(X_{\xi ,[h]},\eta _{[h]})`$ is an element of $`𝒢`$. ###### Lemma 5.3. Let $`h`$ and $`l`$ be two paths connecting the identity to the same element $`gG`$. Then $`(X_{\xi ,[h]},\eta _{[h]})=(X_{\xi ,[l]},\eta _{[l]})`$ in $`𝒢`$. Proof. Let us consider the map $`\gamma :=hl^1:IG`$. Since $`\gamma `$ is the identity at the boundaries of $`I`$ and it is in the connected component of the identity (as a consequence of the fact that $`G`$ is simply connected), this is an allowed gauge transformation. Moreover, an easy computation proves that $`(\eta _{[h]})^\gamma `$ $`=\gamma ^1\eta _{[h]}\gamma +\gamma ^1\mathrm{d}\gamma =\eta _{[l]},`$ $`(X_{\xi ,[h]})^\gamma `$ $`=\mathrm{Ad}_\gamma ^{}X_{\xi ,[h]}=X_{\xi ,[l]}.`$ Therefore, $`(X_{\xi ,[h]},\eta _{[h]})`$ and $`(X_{\xi ,[l]},\eta _{[l]})`$ define the same element in $`𝒢`$. $`\mathrm{}`$ As a consequence we have the following well-defined map: $$j:\begin{array}{ccc}𝔤^{}\times G& & 𝒢\\ (\xi ,g)& & (X_{\xi ,g},\eta _g)\end{array}$$ with $`(X_{\xi ,g},\eta _g):=(X_{\xi ,[h]},\eta _{[h]})`$ for any path $`h`$ from the identity to $`g`$. We then have the following ###### Lemma 5.4. The maps $`i`$ and $`j`$ are inverse to each other. Proof. Since $`\eta _{[h]}`$ is obtained from the trivial connection by the gauge transformation $`h^1`$ (which is not one of the symmetries we allow since $`h`$ at the boundary is not the identity), we see immediately that $`\mathrm{Hol}(\eta _{[h]})=g`$. Moreover, $`X_{\xi ,g}(0)=\xi `$ by definition. So $`ij=\mathrm{id}`$. Next we take a solution $`(X,\eta )`$ of (5.1). Let $`l(u):=\mathrm{Hol}^u(\eta )`$ be the parallel transport from the lower boundary of $`I`$ till the point $`u`$. Notice that $`l`$ is a path from the identity to $`\mathrm{Hol}(\eta )`$. Since moreover $`\eta `$ is equal to $`l\mathrm{d}l^1`$, we see that $`(X,\eta )=(X_{X(0),[l]},\eta _{[l]})`$. But, from the previous Lemma, we get then $`(X,\eta )=(X_{X(0),\mathrm{Hol}(\eta )},\eta _{\mathrm{Hol}(\eta )})=ji(X,\eta )`$. $`\mathrm{}`$ To conclude the proof of Theorem 5.2, we briefly discuss the induced groupoid structure on $`𝔤^{}\times G`$. Recalling that for us the the left and right projections correspond to the boundary values $`X(0)`$ and $`X(1)`$, we obtain $$l(\xi ,g)=\xi ,r(\xi ,g)=\mathrm{Ad}_{g^1}^{}\xi .$$ The product is given by composition of solutions as in (4.2), and under this operation the parallel transports also compose. So we get $$(\xi ,g){}_{}{}^{_{}}(\mathrm{Ad}_{g^1}^{}\xi ,h)=(\xi ,gh).$$ After identifying $`𝔤^{}\times G`$ with $`T^{}G`$ by the map described in Example 4.1, we recover the groupoid structure described there. ## 6. A singular phase space We want to discuss here an example proposed by Weinstein of a regular Poisson manifold that does not admit a symplectic groupoid and show what singularities arise in the construction of the phase space of the corresponding Poisson sigma model. Let $`M=^3\{0\}`$ with Poisson bivector field $$\alpha ^{ij}(x)=f(\|x\|)ϵ_k^{ij}x^k,f(R)0R>0.$$ where $`\|\|`$ is the standard Euclidean norm. For $`f`$ constant this Poisson manifold is equivalent to $`𝔰𝔲(2)\{0\}`$ with the Kirillov–Kostant Poisson structure, and the corresponding phase space is $`(𝔰𝔲(2)^{}\{0\})\times SU(2)`$, as described in the previous section. If we introduce a non constant $`f`$, however, some problems may arise. Observe first that, in any case, the symplectic leaves are the same as in the case of $`𝔰𝔲(2)\{0\}`$, i.e., spheres centered at the origin. The symplectic form on these spheres is however rescaled by a factor $`f`$, and the symplectic area $`A`$ of the sphere with radius $`R`$ is $$A(R)=\frac{4\pi R}{f(R)}.$$ Then the observation of Weinstein , based on theorem of Dazord , is that such a Poisson manifold cannot have a symplectic groupoid if $`A(R)`$ is non constant has critical points. We want to see now how this condition arises in the construction of the phase space. Namely, we have the following ###### Theorem 6.1. The phase space $`𝒢`$ corresponding to $`(M,\alpha )`$ as above is singular iff $`A`$ is non constant and has critical points. ### 6.1. Proof In order to discuss the phase space $`𝒢`$, it is convenient to use a vector notation; viz., we identify $`(^3)^{}`$ and $`^3`$ using the Euclidean scalar product. Moreover, we fix the volume form $`du`$ on the interval $`I=[0,1]`$. Then both our fields $`X`$ and $`\eta `$ can be identified with functions from $`I`$ to $`^3`$ that we denote by $`𝐗`$ and $`𝜼`$. With these notations the Gauss law reads (6.1) $$𝐗^{}+f(\|𝐗\|)𝜼\times 𝐗=0,$$ where $`\times `$ denotes the cross product. The infinitesimal symmetries can also be written in vector notation after identifying $`\beta `$ with a map $`𝜷:I^3`$: (6.2) $$\begin{array}{cc}\hfill \dot{𝐗}& =f(\|𝐗\|)𝜷\times 𝐗,\hfill \\ \hfill \dot{𝜼}& =𝜷^{}+f(\|𝐗\|)𝜼\times 𝜷+\frac{f^{}(\|𝐗\|)}{\|𝐗\|}(𝐗𝜼\times 𝜷)𝐗,\hfill \end{array}$$ where $``$ is the Euclidean scalar product. Given a map $`𝐯:I^3`$ (e.g., $`𝜼`$ or $`𝜷`$), we define its radial component $`v_r`$ and its tangential part $`𝐯_t`$ w.r.t. $`𝐗`$ by: (6.3) $$v_r(u):=\frac{𝐯(u)𝐗(u)}{\|𝐗(u)\|},𝐯_t(u):=𝐯(u)v_r(u)\frac{𝐗(u)}{\|𝐗(u)\|}.$$ Then we have the following: ###### Lemma 6.2. With the decomposition in (6.3), the Gauss law reads $$𝐗^{}+f(\|𝐗\|)𝜼_t\times 𝐗=0,$$ while the symmetries can be written as $$\begin{array}{cc}\hfill \dot{𝐗}& =f(\|𝐗\|)𝜷_t\times 𝐗,\hfill \\ \hfill \dot{\eta }_r& =\beta _r^{}\frac{f(\|𝐗\|)}{\|𝐗\|}(1C(\|𝐗\|))(𝐗𝜼_t\times 𝜷_t),\hfill \\ \hfill \dot{𝜼}_t& =𝜷_t^{}+f(\|𝐗\|)𝜼_t\times 𝜷_t+\frac{f(\|𝐗\|)}{\|𝐗\|^2}(𝐗𝜼_t\times 𝜷_t)𝐗,\hfill \end{array}$$ with $$C(R)=\frac{Rf^{}(R)}{f(R)}=1\frac{f(R)A^{}(R)}{4\pi }.$$ Proof. The Gauss law and the symmetry transformation for $`𝐗`$ simply follow from the fact that in a cross product or in a triple product containing $`𝐗`$ only tangential components of other vectors contribute. For the symmetry transformation of $`𝜼`$, first of all we observe that $`\|𝐗\|^{}=\|𝐗\|^{}=0`$. Then we obtain by (6.3), (6.1) and (6.2) the following identities: $`\dot{\eta }_r`$ $`=(\dot{𝜼})_r{\displaystyle \frac{f(\|𝐗\|)}{\|𝐗\|}}𝐗𝜼_t\times 𝜷_t,`$ $`\beta _r^{}`$ $`=(𝜷^{})_r+{\displaystyle \frac{f(\|𝐗\|)}{\|𝐗\|}}𝐗𝜼_t\times 𝜷_t.`$ These yield immediately the symmetry equation for $`\eta _r`$. To obtain the symmetry equation for $`𝜼_t`$, we first observe that $$𝜷_t^{}=𝜷^{}\beta _r^{}\frac{𝐗}{\|𝐗\|}+\frac{f(\|𝐗\|)}{\|𝐗\|}\beta _r𝜼_t\times 𝐗.$$ Then we get $$\begin{array}{c}\dot{𝜼}_t=\dot{𝜼}\dot{\eta }_r\frac{𝐗}{\|𝐗\|}+\frac{f(\|𝐗\|)}{\|𝐗\|}\eta _r𝜷_t\times 𝐗=\hfill \\ \hfill =𝜷_t^{}+f(\|𝐗\|)𝜼\times 𝜷\frac{f(\|𝐗\|)}{\|𝐗\|}\beta _r𝜼_t\times 𝐗+\frac{f(\|𝐗\|)}{\|𝐗\|}\eta _r𝜷_t\times 𝐗+\\ \hfill +\frac{f(\|𝐗\|)}{\|𝐗\|^2}(𝐗𝜼_t\times 𝜷_t)𝐗,\end{array}$$ which, after using again (6.3), leads to the desired identity. $`\mathrm{}`$ Observe now that the original case of $`𝔰𝔲(2)`$ is recovered by setting $`f1`$ and $`C0`$ in the equations displayed in Lemma 6.2. On the other hand, the critical case $`A^{}(R)=0`$ corresponds to $`C(R)=1`$. Let us begin considering solutions with $`C(\|𝐗\|)1`$. In this case, we can define new variables as follows: (6.4) $`a_r`$ $`={\displaystyle \frac{f(\|𝐗\|)}{1C(\|𝐗\|)}}\eta _r,`$ $`𝐚_t`$ $`=f(\|𝐗\|)𝜼_t,`$ (6.5) $`b_r`$ $`={\displaystyle \frac{f(\|𝐗\|)}{1C(\|𝐗\|)}}\beta _r,`$ $`𝐛_t`$ $`=f(\|𝐗\|)𝜷_t.`$ Then we obtain the Gauss law in the form $$𝐗^{}+𝐚_t\times 𝐗=0,$$ while the symmetries read now $$\begin{array}{cc}\hfill \dot{𝐗}& =𝐛_t\times 𝐗,\hfill \\ \hfill \dot{a}_r& =b_r^{}\frac{1}{\|𝐗\|}(𝐗𝐚_t\times 𝐛_t),\hfill \\ \hfill \dot{𝐚}_t& =𝐛_t^{}+𝐚_t\times 𝐛_t+\frac{1}{\|𝐗\|^2}(𝐗𝐚_t\times 𝐛_t)𝐗.\hfill \end{array}$$ Thus we have recovered, in the new variables, the case of $`𝔰𝔲(2)`$. Proceeding now as in the proof of Theorem 5.2 (namely, taking holonomies of $`𝐚`$ as coordinates), we find that the fiber of the left projection of $`𝒢`$ over a point $`𝐱M`$ with $`C(\|𝐱\|)1`$ is diffeomorphic to $`SU(2)`$. On the other hand, when $`C(\|𝐱\|)=1`$, the above change of variables is not defined. In this case we may however choose the following complete set of invariant functions: $$𝐱:=𝐗(0),𝐲:=\frac{𝐗(1)}{\|𝐱\|},\pi :=_I\eta _r.$$ That is, the fiber of the left projection of $`𝒢`$ over $`𝐱`$ with $`C(\|𝐱\|)=1`$ is diffeomorphic to $`S^2\times `$. If $`C1`$—i.e., if $`A`$ is constant—then $`𝒢`$ is the smooth manifold $`^+\times S^2\times S^2\times `$. To better visualize the singularities in the general case, let us pick up a neighborhood $`U`$ of a point in $`^+`$ where $`A^{}`$ vanishes but $`A`$ is non constant. Let $`V`$ be a neighborhood of a point in $`S^2`$. We want to show that $`𝒢_{UV}:=l^1(U\times V)`$ is not a manifold. We can describe $`𝒢_{UV}`$ as follows. Given a solution $`(𝐗,𝜼)`$, we can always transform it into a solution with $`\eta _r`$ constant (just take a transformation generated by $`𝜷`$ with $`𝜷_t=0`$). Under small gauge transformations such a solution is characterized by the values of $`𝐗`$ at the endpoints and the value of $`\eta _r`$. If $`C(\|𝐗\|)=1`$, there is no way of changing $`\eta _r`$ into another constant. On the other hand, if $`C(\|𝐗\|)1`$, large gauge transformations can send $`\eta _r`$ into another constant that differs from the former by a multiple of $`4\pi [1C(R)]/f(R)`$ (we are trivializing the Hopf bundle $`SU(2)S^2`$ over $`V`$ taking into account the rescaling (6.4)). Therefore, $`𝒢_{UV}=V\times V\times Q`$, where $`Q`$ is the quotient of $`U\times `$ by the equivalence relation $$(R,p)(R,p+\frac{4\pi [1C(R)]}{f(R)}).$$ ## 7. The phase space of the Poisson sigma model with two-dimensional target Let $`U`$ be a domain in $`^2`$ with a given Poisson bivector field $`\alpha `$. After choosing coordinates, it is always possible to write $$\alpha ^{ij}(x^1,x^2)=ϵ^{ij}\varphi (x^1,x^2),\varphi C^{\mathrm{}}(U).$$ We also fix the volume form $`du`$ on $`I`$ and then identify $`\eta `$ with a map $`I^2`$. With these choices, the “Gauss law”simply reads (7.1) $`(X^1)^{}+\varphi (X^1,X^2)\eta _2`$ $`=0,`$ $`(X^2)^{}\varphi (X^1,X^2)\eta _1`$ $`=0.`$ The infinitesimal symmetries read then (7.2) $`\delta X^1`$ $`=\varphi \beta _2,`$ $`\delta X^2`$ $`=\varphi \beta _1,`$ $`\delta \eta _1`$ $`=\beta _1^{}+_1\varphi (\eta _1\beta _2\eta _2\beta _1),`$ $`\delta \eta _2`$ $`=\beta _2^{}+_2\varphi (\eta _1\beta _2\eta _2\beta _1),`$ where $`_i\varphi `$ is a shorthand notation for $`\varphi /x^i`$, and the infinitesimal generators $`\beta _i`$ are continuously differentiable maps $`I^2`$ with the conditions (7.3) $$\beta _i(0)=\beta _i(1)=0,i=1,2.$$ We will denote by $`\stackrel{~}{𝒢}`$ the phase space of solutions of (7.1) modulo the symmetries generated by (7.2). More precisely, we first introduce the Banach spaces $`C^0(I,^2)`$, $`C^1(I,^2)`$ and $`C_0^1(I,^2)`$. Then we consider the Banach manifold $`C^1(I,U)`$ modeled on $`C^1(I,^2)`$. With these notations we can finally write $$\stackrel{~}{𝒢}:=\frac{\{(X,\eta )C^1(I,U)\times C^0(I,^2)\|(X,\eta )\text{ satisfies (}\text{7.1}\text{)}\}}{\{\text{symmetries (}\text{7.2}\text{) with }\beta C_0^1(I,^2)\}}.$$ In the rest of this section we will study $`\stackrel{~}{𝒢}`$. Namely, in subsection 7.1 we will give an equivalent but easier description of $`\stackrel{~}{𝒢}`$, and in subsection 7.2 we will show that the latter is diffeomorphic to a submanifold $`𝒢`$ of $`^4`$, at least if all the symplectic leaves of $`U`$ are simply connected; in subsection 7.4 we will describe the product structure for $`𝒢`$ induced from the composition of paths $`X:IU`$; finally, in subsection 7.5 we will derive the symplectic structure for $`𝒢`$ from the symplectic structure on the space of fields $`(X,\eta )`$. ### 7.1. An equivalent description of the phase space ¿From now on, by abuse of notation, we will write $`\varphi `$ for $`\varphi X`$. The Gauss law (7.1) implies an equation for $`\varphi `$, viz., (7.4) $$\varphi ^{}=T\varphi ,$$ with (7.5) $$T:=_2\varphi \eta _1_1\varphi \eta _2.$$ The solution of (7.4) is simply given by (7.6) $$\varphi (X^1(u),X^2(u))=\varphi _0H(u),$$ where $`\varphi _0`$ is a shorthand notation for $`\varphi (X^1(0),X^2(0))`$ and (7.7) $$H(u):=\mathrm{exp}_0^uT(v)𝑑v.$$ It is then useful to define the following change of variables: (7.8) $$E_i:=\eta _iH,$$ Notice that the map $`(X,\eta )(X,E)`$ is a smooth map from $`C^1(I,U)\times C^0(I,^2)`$ into itself. With these new variables, we can rewrite the Gauss law (7.1) as (7.9) $`(X^1)^{}+\varphi _0E_2`$ $`=0,`$ $`(X^2)^{}\varphi _0E_1`$ $`=0.`$ Notice that every solution of (7.1) determines a solution of (7.9) via (7.8). The converse, however, is not true in general. Assume in fact that $`(X,E)`$ is a solution of (7.9). Then we get the following equation for $`\varphi `$: (7.10) $$\varphi ^{}=\widehat{T}\varphi _0,$$ with (7.11) $$\widehat{T}:=_2\varphi E_1_1\varphi E_2.$$ Comparing the solution of (7.10) with (7.6), we get (7.12) $$H(u)=1+_0^u\widehat{T}(v)𝑑v.$$ By comparison with (7.7), we see that a solution $`(X,E)`$ of (7.9) determines a solution $`(X,\eta )`$ of (7.1) iff the following condition is satisfied: (7.13) $$H(u)>0,uI.$$ So we have the following ###### Lemma 7.1. Solutions of (7.1) are mapped by (7.8) into solutions of (7.9) satisfying (7.13) and vice versa. As for the symmetries acting on $`(X,E)`$, we introduce (7.14) $$e_i:=\beta _iH.$$ Observe here that the map $`(X,\eta ,\beta )e`$ is a smooth map from $`C^1(I,U)\times C^0(I,^2)\times C_0^1(I,^2)C_0^1(I,^2)`$. Then we have the following: ###### Lemma 7.2. Under the infinitesimal symmetry (7.2), the variables $`(X,E)`$ defined via (7.8) in terms of a solution $`(X,\eta )`$ of (7.1) change as follows: (7.15) $`\delta X^1`$ $`=\varphi _0e_2,`$ $`\delta X^2`$ $`=\varphi _0e_1,`$ $`\delta E_1`$ $`=e_1^{},`$ $`\delta E_2`$ $`=e_2^{}.`$ Conversely, if $`(X,E)`$ is a solution of (7.9) satisfying to (7.13), then the infinitesimal symmetry (7.15) implies the infinitesimal symmetry (7.2) on the variables $`(X,\eta )`$ obtained by inverting (7.8). Proof. The first two equations are immediately obtained by (7.6) and (7.14). As for the two other equations, we first observe that $$\delta T=\tau ^{},$$ with $$\tau :=_2\varphi \beta _1_1\varphi \beta _2.$$ In fact, $$\begin{array}{c}\delta T=\delta (ϵ^{ij}\eta _i_j\varphi )=\hfill \\ \hfill =ϵ^{ij}\beta _i^{}_j\varphi +ϵ^{ij}_i\varphi ϵ^{kl}\eta _k\beta _l_j\varphi ϵ^{ij}\eta _i\varphi _{jk}ϵ^{kl}\varphi \beta _l=\\ \hfill =\frac{\mathrm{d}}{\mathrm{d}u}(ϵ^{ij}\beta _i_j\varphi )=\tau ^{},\end{array}$$ where we have made use of (7.1) and (7.2). From this we get $$\delta H=\tau H.$$ Finally, $$\delta E_i=\delta \eta _iH+\eta _i\tau H=\frac{\mathrm{d}}{\mathrm{d}u}(\beta _iH)(\beta _iT\eta _i\tau _i\varphi ϵ^{kl}\eta _k\beta _l)H.$$ A direct computation shows that the terms in the second brackets cancel, so we have proved the first part of the Lemma. As for the second part, we observe that $$\delta \widehat{T}=\widehat{\tau }^{},$$ with $`\widehat{\tau }=_2\varphi e_1_1\varphi e_2`$. As a consequence, $`\delta H=\widehat{\tau }`$. Observing then that $`\widehat{T}=HT`$ and $`\widehat{\tau }=H\tau `$, we get $$\begin{array}{c}\delta \eta _i=\delta \left(\frac{E_i}{H}\right)=\frac{\delta E_i}{H}\frac{E_i}{H^2}\delta H=\frac{e_i^{}}{H}\frac{E_i}{H^2}\widehat{\tau }=\hfill \\ \hfill =\frac{\frac{\mathrm{d}}{\mathrm{d}u}(\beta _iH)}{H}\eta _i\tau =\beta _i^{}+\beta _iT\eta _i\tau ,\end{array}$$ from which (7.2) follows. $`\mathrm{}`$ We now define a new phase space: $$\stackrel{~}{\stackrel{~}{𝒢}}:=\frac{\{(X,E)C^1(I,U)\times C^0(I,^2)\|(X,E)\text{ satisfies (}\text{7.9}\text{) and (}\text{7.13}\text{)}\}}{\{\text{symmetries (}\text{7.15}\text{) with }eC_0^1(I,^2)\}}.$$ Then the preceding discussion, and in particular the two Lemmata, prove the following ###### Proposition 7.3. If $`\stackrel{~}{\stackrel{~}{𝒢}}`$ is a smooth manifold, then $`\stackrel{~}{𝒢}`$ and $`\stackrel{~}{\stackrel{~}{𝒢}}`$ are diffeomorphic. In the next subsection we will prove that $`\stackrel{~}{\stackrel{~}{𝒢}}`$ is actually a smooth $`4`$-manifold, at least under the following ###### Assumption 7.4. We assume that all the symplectic leaves of $`(U,\alpha )`$ are simply connected. Observe that for example $`^2`$ with $`\varphi =(x^1)^2+(x^2)^2`$ will not be allowed. ### 7.2. The phase space is a smooth manifold We begin by defining some invariants of $`\stackrel{~}{\stackrel{~}{𝒢}}`$. The first are the initial conditions of $`X`$, viz., (7.16) $$x^i:=X^i(0);$$ the others are the following integrals: (7.17) $$\pi _i:=_0^1E_i(u)𝑑u.$$ All of them are invariant under (7.15) since $`eC_0^1(I,^2)`$. In this way we get a well-defined, smooth map $`i:\stackrel{~}{\stackrel{~}{𝒢}}U\times ^2`$. This map is however not surjective because of (7.13). We want then to define an appropriate domain in $`^4`$ so that $`i`$ becomes a diffeomorphism. We first define the continuous map $`x_f:U\times ^2^2`$ by (7.18) $$x_f^i=x^i\varphi (x^1,x^2)ϵ^{ij}\pi _j,$$ and then $$V:=\{pU\times ^2\|x_f(p)U\}.$$ Remark that $`x_f`$ can also be seen as the final point of a solution $`X`$ of (7.9), with $`x`$ and $`\pi `$ given by (7.16) and (7.17). Next we define the map $`h:V`$ by (7.19) $$h(x^1,x^2,\pi _1,\pi _2):=\{\begin{array}{cc}\frac{\varphi (x_f^1,x_f^2)}{\varphi (x^1,x^2)}\hfill & \text{if }\varphi (x^1,x^2)0,\hfill \\ 1\pi _2_1\varphi (x^1,x^2)+\pi _1_2\varphi (x^1,x^2)\hfill & \text{if }\varphi (x^1,x^2)=0.\hfill \end{array}$$ Then we define $$𝓖:=\{pV\|h(p)>0\},$$ and finally we denote by $`𝒢`$ the connected component of $`𝓖`$ containing $`U\times \{(0,0)\}`$. ###### Lemma 7.5. $`𝒢`$ is a 4-manifold. Proof. We just have to prove that $`h`$ is continuous. To do this, we observe that the two definitions for $`h`$ are continuous when restricted to the appropriate subset. Since the zero locus of $`\varphi `$ is closed, we only have to check that, for any sequence in the complement of the zero locus that approaches a point in the zero locus, the limit of the first expression yields the second expression. This can be easily proved by Taylor expanding the numerator. To prove that the connected component we are interested in is not empty, it is enough to observe that $`h(x^1,x^2,0,0)=1`$, $`(x^1,x^2)U`$. $`\mathrm{}`$ ###### Example 7.6 (Semiclassical quantum plane). Let $`U=^2`$ and $`\varphi (x^1,x^2)=x^1x^2`$. In this case the map $`h`$ simply reads $$h(x^1,x^2,\pi _1,\pi _2)=(1x^2\pi _2)(1+x^1\pi _1).$$ In the connected component of $`h^1(^+)`$ both factors must be positive. So we get $$𝒢=\{(x^1,x^2,\pi _1,\pi _2)^4\|x^1\pi _1>1,x^2\pi _2<1\}.$$ The 2-dimensional symplectic leaves are the four open quadrants. Over each point $`(x^1,x^2)`$ of one of these leaves, the fiber is given by those vectors $`(\pi _1,\pi _2)`$ such that the linear trajectory with constant velocity $$(\varphi (x^1,x^2)\pi _2,\varphi (x^1,x^2)\pi _1)$$ is entirely contained in the same symplectic leaf for all times $`t1`$. Over points in the zero locus of $`\varphi `$, i.e., the axes, the fiber is the whole of $`^2`$, for the velocity here is zero. Observe that this simple description of $`𝒢`$ is possible whenever all symplectic leaves in $`U`$ are convex; e.g., when $`U=^2`$ and $`\varphi (x^1,x^2)=(x^1)^r(x^2)^s`$, $`r,s0`$. The central result of this section is the following: ###### Theorem 7.7. Under Assumption 7.4, the phase spaces $`\stackrel{~}{𝒢}`$ and $`\stackrel{~}{\stackrel{~}{𝒢}}`$ are diffeomorphic to $`𝒢`$. ### 7.3. Proof of Theorem 7.7 In view of Proposition 7.3, we have only to prove that $`𝒢`$ is diffeomorphic to $`\stackrel{~}{\stackrel{~}{𝒢}}`$. The idea is to show that the mapping given by (7.16) and (7.17) defines a diffeomorphism. ###### Lemma 7.8. There is a well-defined smooth map $$i:\begin{array}{ccc}\stackrel{~}{\stackrel{~}{𝒢}}& & 𝒢\\ [(X,E)]& & (X(0),_0^1E(u)𝑑u)\end{array}$$ Proof. Since (7.16) and (7.17) are invariant under the symmetries (7.15), this map descends to $`\stackrel{~}{\stackrel{~}{𝒢}}`$. We want then to show that the image of this map is contained in $`𝒢`$. First of all, we observe that $`x_f=X(1)`$; so automatically $$x_f(x^1,x^2,\pi _1,\pi _2)U.$$ Then we want to show that (7.13) implies $`h>0`$. Consider first the case when $`\varphi _00`$. Then, by (7.6), $`\varphi (u)/\varphi _0=H(u)>0,uI`$. In particular, for $`u=1`$, this implies $`\varphi (x_f)/\varphi (x)>0`$. The other possibility is when $`\varphi _0=0`$. In this case $`X`$ is constant, and by (7.12) we get that $`H(1)=h`$. Since for every solution $`(X,E)`$ of (7.9) and (7.13) there is a solution with the same initial condition and $`\stackrel{~}{E}=E/\lambda `$, $`\lambda >0`$, then $`\stackrel{~}{\stackrel{~}{𝒢}}`$ is connected, so its image is contained in a connected component of the set $`h^1(^+)`$. The existence of constant solutions with $`E=0`$ implies that this connected component contains $`U\times \{(0,0)\}`$. So we have proved that $`[(X,E)](x,\pi )`$ is a well-defined map from $`\stackrel{~}{\stackrel{~}{𝒢}}`$ to $`𝒢`$. $`\mathrm{}`$ We now want to define an inverse map $$j:\begin{array}{ccc}𝒢& & \stackrel{~}{\stackrel{~}{𝒢}}\\ (x,\pi )& & [(X,E)]\end{array}$$ We consider two cases: 1. $`\varphi (x)0`$. We take $`X`$ equal to any path joining $`x`$ to $`x_f(x,\pi )`$ that is completely contained in the symplectic leaf. Then we set $`E_i(u)=ϵ_{ij}(X^j)^{}(u)/\varphi (x)`$. 2. $`\varphi (x)=0`$. We set $`X(u)=x`$ and $`E(u)=\pi `$ $`uI`$. It is not difficult to see that the image of $`j`$ is a solution of (7.9) and (7.13). ###### Lemma 7.9. Let $`(X,E)`$ and $`(\stackrel{~}{X},\stackrel{~}{E})`$ be two solutions determined as above. Then, under Assumption 7.4, they define the same element in $`\stackrel{~}{\stackrel{~}{𝒢}}`$. Proof. In the case when $`\varphi (x)=0`$, we completely specified the solution; so $`(X,E)=(\stackrel{~}{X},\stackrel{~}{E})`$. Let us consider then the case $`\varphi (x)0`$. Since any symplectic leaf is simply connected by Assumption 7.4, there is a path $`X(u,s)`$ connecting $`X`$ to $`\stackrel{~}{X}`$. More precisely, $`X(u,0)=X(u)`$, $`X(u,1)=\stackrel{~}{X}(u)`$, $`X(0,s)=x`$, $`X(1,s)=x_f`$, and $`X({}_{}{}^{_{}},s)`$ is entirely contained in the symplectic leaf of $`x`$. We set then $$e_i=ϵ_{ij}\frac{\dot{X}^j}{\varphi (x)}$$ and integrate the infinitesimal symmetry (7.15) obtaining $$E_i(u,s)=E_i(u)+_0^se_i^{}(u,\sigma )𝑑\sigma .$$ $`\mathrm{}`$ As a consequence the map $`j`$ is well-defined, and it is not difficult to prove that it is smooth. Moreover, we have the following ###### Lemma 7.10. Under Assumption 7.4, the maps $`i`$ and $`j`$ are inverse to each other. Proof. The identity $`ij=\mathrm{id}`$ is trivial. We want to prove that also $`ji=\mathrm{id}`$. Let us begin with the case $`\varphi _00`$. In this case $`ji[(X,E)]`$ is a solution $`(\stackrel{~}{X},\stackrel{~}{E})`$ so that $`\stackrel{~}{X}`$ has the same end-points of $`X`$. Thus, as in the proof of the previous Lemma, there is a symmetry that relates them. In the case when $`\varphi _0=0`$, we must prove that, given a solution $`(X,E)`$, then $`(\stackrel{~}{X},\stackrel{~}{E}):=ji(X,E)`$ is an equivalent solution. First we observe that $`X=\stackrel{~}{X}`$ since both are constant solutions with the same starting point. We have then to find a symmetry that sends $`E`$ to $`\stackrel{~}{E}`$. To do so, we define the following element of $`C_0^1(I,^2)`$: $$e_i(u):=_0^u(\stackrel{~}{E}_iE_i(v))𝑑v.$$ Then we consider the path in $`C^0(I,^2)`$ given by $$E_i(u,s):=s\stackrel{~}{E}_i+(1s)E_i(u),s[0,1].$$ We have then $`E_i(u,0)`$ $`=E_i(u),`$ $`E_i(u,1)`$ $`=\stackrel{~}{E}_i,`$ $`\dot{E}_i(u,s)`$ $`=e_i^{}(u).`$ So we can go from $`E`$ to $`\stackrel{~}{E}`$ via a symmetry transformation (7.15). Since $`\varphi _0=0`$, the corresponding path of paths $`X(u,s)`$ is constant and equal to $`X(0)`$. To complete the proof, we have only to check that condition (7.13) is satisfied for any intermediate value $`u[0,1]`$. To do this we just observe that, by definition, $$\begin{array}{c}H(u,s)=\hfill \\ \hfill =1+_2\varphi [s\stackrel{~}{E}_1+(1s)E_1(u)]_1\varphi [s\stackrel{~}{E}_2+(1s)E_2(u)]=\\ \hfill =sA+(1s)B(u),\end{array}$$ with $`A=1+_2\varphi \stackrel{~}{E}_1_1\varphi \stackrel{~}{E}_2`$ and $`B(u)=1+_2\varphi E_1(u)_1\varphi E_2(u)`$. Since $`A>0`$ and $`B(u)>0uI`$, we get that $`H(u,s)>0(u,s)I\times [0,1]`$. $`\mathrm{}`$ This concludes the proof of Theorem (7.7). ### 7.4. The product on $`𝒢`$ We will describe $`𝒢`$ in terms of local coordinates $`(x,\pi )`$ with $`xU`$, $`\pi ^2`$. We define the two projections $`r,l:𝒢U`$, by $$l(x,\pi ):=x,r(x,\pi ):=x_f(x,\pi ),$$ which correspond to the initial and final point of the given solution in $`\stackrel{~}{𝒢}`$ as prescribed by (4.1). Let us consider now another point $`(\stackrel{~}{x},\stackrel{~}{\pi }_1)𝒢`$, with $`\stackrel{~}{x}=x_f(x,\pi )`$. Then we look for the solutions $`(X,\eta )`$ and $`(\stackrel{~}{X},\stackrel{~}{\eta })`$ in $`\stackrel{~}{𝒢}`$ that correspond to the points in $`𝒢`$ described above. In particular we choose the solutions so that the tangent at $`X(1)`$ is equal to the tangent at $`\stackrel{~}{X}(0)`$. So we can compose the solutions in a differentiable way as in (4.2). We now want to compute the point $`(\widehat{x},\widehat{\pi })𝒢`$ corresponding to the new solution $`(\widehat{X},\widehat{\eta })`$. We immediately get (7.20) $$\widehat{x}=x.$$ As for $`\widehat{\pi }`$, we use (7.17) and (7.8) obtaining $$\begin{array}{c}\widehat{\pi }=_0^1\widehat{H}(u)\widehat{\eta }(u)𝑑u=_0^1\widehat{H}\left(\frac{u}{2}\right)\eta (u)𝑑u+_0^1\widehat{H}\left(\frac{u+1}{2}\right)\stackrel{~}{\eta }(u)𝑑u.\hfill \end{array}$$ By (7.7) we get then $$\widehat{H}\left(\frac{u}{2}\right)=H(u),H\left(\frac{u+1}{2}\right)=H(1)\stackrel{~}{H}(u),$$ for $`u[0,1]`$. Thus, we get (7.21) $$\widehat{\pi }=\pi +h(x,\pi )\stackrel{~}{\pi }.$$ ### 7.5. The symplectic structure on $`𝒢`$ As in the general description, we consider the constant symplectic structure on $`C^1(I,U)\times C^0(I,^2)`$ determined by the action $`_0^1\eta _i(u)(X^i)^{}(u)`$; viz., $$\omega ((\alpha ,\beta ),(\stackrel{~}{\alpha },\stackrel{~}{\beta })):=_0^1[\stackrel{~}{\alpha }(u)\beta (u)\alpha (u)\stackrel{~}{\beta }(u)]𝑑u,$$ for $`(\alpha ,\beta ),(\stackrel{~}{\alpha },\stackrel{~}{\beta })T(C^1(I,U)\times C^0(I,^2))`$. In order to perform the computations of this subsection, it is however more convenient to work with the corresponding Poisson structure that we write $$\{\eta _i(u),X^j(v)\}=\delta _i^j\delta (uv),$$ while all other brackets vanish. As usual in infinite dimensional cases, the Poisson bracket is defined only for a certain class of functions. We now want to determine the induced Poisson structure on $`𝒢`$. By the general argument we get (7.22) $$\{x^1,x^2\}=\varphi (x^1,x^2).$$ An easy computation yields $$\{T(u),X^i(0)\}=ϵ^{ij}_j\varphi _0\delta (u).$$ Thus we get (7.23) $$\begin{array}{c}\{x^i,\pi _j\}=\{X^i(0),_0^1H(u)\eta _j(u)𝑑u\}=\hfill \\ \hfill =\delta _i^j+_0^1H(u)\eta _j(u)\{X^i(0),_0^uT(v)𝑑v\}=\\ \hfill =\delta _i^j\pi _jϵ^{ik}_k\varphi (x^1,x^2).\end{array}$$ Finally, we have the most complicated bracket, that is, $`\{\pi _1,\pi _2\}`$. ###### Lemma 7.11. Let us consider the function $`\psi :𝒢`$ defined by (7.24) $$\psi :=\{\begin{array}{cc}\frac{1+\pi _1_2\varphi \pi _2_1\varphi h}{\varphi }\hfill & \text{if }\varphi 0,\hfill \\ \pi _1\pi _2_1_2\varphi \frac{1}{2}(\pi _1)^2(_2)^2\varphi \frac{1}{2}(\pi _2)^2(_1)^2\varphi \hfill & \text{if }\varphi =0,\hfill \end{array}$$ with $`h`$ defined in (7.19). Then $`\psi `$ is smooth and (7.25) $$\{\pi _1,\pi _2\}=\psi (x^1,x^2,\pi _1,\pi _2).$$ Proof. The smoothness of $`\psi `$ is proved by Taylor expanding $`h`$ in the first expression. As for the second assertion, we first observe the following useful identities: $`\{\eta _i(u),T(v)\}`$ $`=\delta (uv)ϵ^{kl}\eta _k(u)_i_l\varphi (u),`$ $`\{T(u),T(v)\}`$ $`=0,`$ which imply $`\{\eta _i(u),H(v)\}`$ $`=\theta (vu)ϵ^{kl}H(v)\eta _k(u)_i_l\varphi (u),`$ $`\{H(u),H(v)\}`$ $`=0.`$ Then a straightforward computation yields $$\begin{array}{c}\{\pi _1,\pi _2\}=_0^1du_0^1dv[E_1(u)E_2(v)(\theta (uv)_1_2\varphi _v+\theta (vu)_1_2\varphi _u)+\hfill \\ \hfill \theta (uv)E_1(u)E_1(v)(_2)^2\varphi _v\theta (vu)E_2(u)E_2(v)(_1)^2\varphi _v],\end{array}$$ where $`\varphi _v`$ is a short-hand notation for $`\varphi (X^1(v),X^2(v))`$. In the case when $`\varphi _0=0`$, the solution $`X`$ is constant. So we can take all the terms of the form $`_i_j\varphi `$ out of the integral. What is left, thanks to (7.17), yields the second formula in (7.24). If however $`\varphi _00`$, we multiply both sides by $`(\varphi _0)^2`$ and then use (7.9), obtaining $$\begin{array}{c}(\varphi _0)^2\{\pi _1,\pi _2\}=_0^1du_0^1dv[\theta (uv)(X^2)^{}(u)\frac{\mathrm{d}}{\mathrm{d}v}_2\varphi _v+\hfill \\ \hfill +\theta (vu)(X^1)^{}(v)\frac{\mathrm{d}}{\mathrm{d}u}_1\varphi _u].\end{array}$$ A simple integration yields then $$(\varphi _0)^2\{\pi _1,\pi _2\}=(X^2(1)X^2(0))_2\varphi _0+(X^1(1)X^1(0))_1\varphi _0\varphi _1+\varphi _0,$$ which is the first formula in (7.24) thanks to (7.16), to (7.18) and to the identity $`h=\varphi _1/\varphi _0`$. $`\mathrm{}`$ The brackets of the coordinates define a bivector field on $`𝒢`$, which we will denote by $`𝔓`$, through the relation (7.26) $$\{f,g\}=𝔓(\mathrm{d}f,\mathrm{d}g).$$ Locally, in the basis corresponding to the coordinates $`x^1,x^2,\pi _1,\pi _2`$, we write this bivector field in matrix form as (7.27) $$𝖯=\left(\begin{array}{cccc}0& \varphi & 1\pi _1_2\varphi & \pi _2_2\varphi \\ \varphi & 0& \pi _1_1\varphi & 1+\pi _2_1\varphi \\ 1+\pi _1_2\varphi & \pi _1_1\varphi & 0& \psi \\ \pi _2_2\varphi & 1\pi _2_1\varphi & \psi & 0\end{array}\right).$$ This matrix is always invertible thanks to the condition $`h>0`$. We will exhibit the corresponding 2-form $`\omega _𝒢`$ in the next subsection. ### 7.6. Summary We started with a 2-dimensional domain $`U`$ with a bivector field $`\alpha ^{ij}=ϵ^{ij}\varphi `$, $`\varphi C^{\mathrm{}}(U)`$, such that Assumption 7.4 holds. Using $`\varphi `$ we defined the 4-dimensional domain $`𝒢`$ as at the beginning of subsection 7.2, viz., as the connected component containing $`U\times \{(0,0)\}`$ of the set $$\{(x,\pi )U\times ^2\|x_f(x,\pi )U,h(x,\pi )>0\},$$ with $`x_f`$ and $`h`$ defined in (7.18) and (7.19). Next we obtained the left and right projections $`l,r:𝒢U`$ by $`l(x,\pi )`$ $`=x,`$ $`r(x,\pi )`$ $`=x_f(x,\pi )=x\alpha \pi .`$ Given two points $`(x,\pi )`$ and $`(\stackrel{~}{x},\stackrel{~}{\pi })`$ in $`𝒢`$ with $`r(x)=l(\stackrel{~}{x})`$, we got the product $$(x,\pi ){}_{}{}^{_{}}(\stackrel{~}{x},\stackrel{~}{\pi })=(x,\pi +h(x,\pi )\stackrel{~}{\pi }).$$ Finally, in (7.27) we defined a bivector field $`𝔓`$ whose inverse exists and is given by the following 2-form: $$\begin{array}{c}\omega _𝒢=[\psi \mathrm{d}x^1\mathrm{d}x^2+(1\pi _2_1\varphi )\mathrm{d}x^1\mathrm{d}\pi _1+\pi _1_1\varphi \mathrm{d}x^2\mathrm{d}\pi _2+\hfill \\ \hfill \pi _2_2\varphi \mathrm{d}x^2\mathrm{d}\pi _1+(1+\pi _1_2\varphi )\mathrm{d}x^2\mathrm{d}\pi _2\varphi \mathrm{d}\pi _1\mathrm{d}\pi _2]/h.\end{array}$$ From the general results of Section 4, we get then the following: ###### Theorem 7.12. $`(𝒢,r,l,{}_{}{}^{_{}},\omega _𝒢)`$ is a symplectic groupoid for $`(U,\alpha )`$. ###### Remark 7.13. It is interesting to note that the above theorem holds also without Assumption 7.4, as can be proved directly. However, $`𝒢`$ as we have defined it is not the phase space for $`(U,\alpha )`$ in the general case, the missing information being a class of homotopic paths inside a symplectic leaf joining the given endpoints.
warning/0003/hep-ph0003171.html	ar5iv	text	# 1 Introduction ## 1 Introduction We shall shall base our study of the hierarchy problems on a finite quantum field theory (FQFT) which is gauge invariant, finite and unitary to all orders of perturbation theory \[1-16\]. In contrast to superstring theory and membrane theory , our FQFT formalism does not require the existence of extra dimensions to guarantee its consistency, so we shall develop our field theory in a four-dimensional spacetime. This has the advantage over string and membrane theory that we do not have to concern ourselves with the problems of compactification. In string theory there are numerous ways to compactify to lower dimensions and it is difficult at present to justify any unique method for achieving this compactification. Moreover, as has been learned during the recent past, the string scale is not uniquely fixed at the Planck scale . String theory is not a quantum field theory in the usual sense. All attempts to formulate a string field theory have so far been unsuccessful. Such theories demonstrate instabilities equivalent to those met with in $`\varphi ^3`$ field theories, in which the Hamiltonian is unbounded from below . Our quantum field theory has strictly local tree graphs and nonlocal quantum loop graphs. The FQFT gauge formalism is applied to guarantee a self-consistent quantum gravity theory coupled to the Yang-Mills, Higgs and spinor fields. The formalism is free of tachyons and unphysical ghosts and satisfies unitarity to all orders of perturbation theory. It could incorporate supersymmetry if required, in the form of a supergravity theory, but we shall not do so here, in order to aim for as minimal a scheme as possible. No attempt is made to unify gauge fields with gravity, or to extend the standard model, for we wish to focus on the hierarchy problems. It is commonly believed that local quantum field theory is the only way to guarantee a consistent Poincaré invariant quantum mechanics . If one is willing to give up the notion of a strictly local observable, then this belief can be shown to be incorrect. The issue depends on the support of the field operators and for nonlocal field theories, it can be shown that it is impossible to construct observables whose support is a compact set. Such theories emerge as “quasi-local” field theories whose local behavior only acts at distances much larger than a certain length scale $`\mathrm{}`$. Nonlocal field theories were the subject of considerable study in the 60s and 70s, because it was thought that they could cure the problems of non-renormalizable field theories. Recently, there has been renewed interest in nonlocal field theories in connection with string theory, M-theory and Little String Theory (LST) . The LSTs are generated by decoupling gravity and other bulk modes from five-branes. One takes, for example, $`N`$ coincident five-branes and considers the limit $`g_s\mathrm{},M_s\mathrm{}`$, where $`g_s`$ is the string coupling at spatial infinity and $`M_s=1/\sqrt{\alpha ^{}}`$ is the string scale. It was shown by Kapustin that LSTs do not possess local observables and that due to the exponentially increasing density of states, the Wightman functions are not polynomially bounded. However, he showed that the nonlocality can be accomodated by choosing a space of test functions different from the usual Schwartz space. We shall consider these issues in more detail in Sect. 4. The standard gauge symmetry of local quantum field theory is generalized to a nonlocal transformation consisting of an inhomogeneous term, which preserves the local, quadratic part of the action, and a nonlocal homogeneous term, which generates a variation of the free field action that cancels the inhomogeneous variation of the nonlocal action. This generalized gauge transformation is similar to the nonlocal gauge transformation of string field theory . The key to the success of string theory lies in the nature of this generalized gauge invariance. Its existence guarantees the raison d’$`\widehat{e}`$tre of gauge symmetry in quantum field theory, which is to decouple unphysical vector and tensor quanta while maintaining Poincaré invariance. The fundamental gauge hierarchy problem is resolved because the finite scalar Higgs self-energy loop graphs are damped exponentially at high energies above the physical Higgs scale $`\mathrm{\Lambda }_H`$ set by the FQFT formalism and by choosing $`\mathrm{\Lambda }_H1`$ TeV. The constant $`\mathrm{\Lambda }_H`$ enters naturally, because it sets the physical non-localizable energy scale of the Higgs particle quantum loop graphs. The cosmological constant problem is considered to be the most severe hierarchy problem in modern physics . The problem arises because, in contrast to classical Newtonian gravity theory, the Einstein gravitational Lagrangian $`_{\mathrm{grav}}`$ is not invariant under the translation $`_{\mathrm{grav}}_{\mathrm{grav}}+C`$, where $`C`$ is a constant identified with the cosmological constant $`\lambda `$. Many attempts to solve this hierarchy problem have been made , and most recently there has been a proposal to solve the problem by postulating a composite graviton connected to string theory . A model based on (3+1) branes and a five-dimensional bulk has also recently been proposed. Solving the cosmological constant problem appears to demand a low-energy mechanism to cancel soft photon loop contributions. How can we obtain such a mechanism in the low-energy framework without destroying the familar successes of the standard model? We shall propose a quantum gravity solution to the problem based on FQFT. We can define an effective cosmological constant $$\lambda _{\mathrm{eff}}=\lambda +\lambda _{\mathrm{vac}},$$ (1) where $`\lambda `$ is the “bare” cosmological constant in Einstein’s classical field equations, and $`\lambda _{\mathrm{vac}}`$ is the contribution that arises from the vacuum density $`\lambda _{\mathrm{vac}}=8\pi G\rho _{\mathrm{vac}}`$. Already at the standard model electroweak scale $`10^2`$ GeV, a calculation of the vacuum density $`\rho _{\mathrm{vac}}`$ results in a discrepancy with the observational bound $$\rho _{\mathrm{vac}}<10^{47}(\mathrm{GeV})^4,$$ (2) of order $`10^{55}`$, resulting in a a severe fine tuning problem, since the virtual quantum fluctuations giving rise to $`\lambda _{\mathrm{vac}}`$ and the “bare” cosmological constant $`\lambda `$ must cancel to an unbelievable degree of accuracy. If we choose the quantum gravity scale $`\mathrm{\Lambda }_G10^4`$ eV, then our quantum gravity theory leads to an exponential damping of gravitational vacuum polarisation for $`p^2\mathrm{\Lambda }_G^2`$, where $`p^2`$ is the square of the Euclidean graviton momentum. This suppresses the cosmological constant $`\lambda _{\mathrm{vac}}`$ below the observational bound (2). Since the graviton tree graphs in our FQFT are identical to the standard point like, local tree graphs of perturbative gravity, we retain classical, causal GR and Newtonian gravity theory, and the measured value of the gravitational constant $`G`$. Only the quantum gravity loop graphs are suppressed above energies $`10^4`$ eV. Thus, at very low energies or large distances, the point-like, local graviton dominates giving rise to classical Newtonian and GR dynamics. The scales $`\mathrm{\Lambda }_H`$ and $`\mathrm{\Lambda }_G`$ are determined by the quantum non-localizable nature of the Higgs particle and the gauge particles $`W`$ and $`Z`$ of the standard model as compared to the graviton. The Higgs particle radiative corrections have a nonlocal scale at $`\mathrm{}_H10^{16}`$ cm, whereas the graviton radiative corrections are localizable down to a large length scale $`\mathrm{}_G1`$ cm. Thus, the fundamental energy scales in the theory are determined by the underlying physical nature of the particles and fields and do not correspond to arbitrary cut-offs, which destroy the gauge invariances of the field theory. The underlying explanation of these physical scales must be sought in a more fundamental theory. In Section 2, we describe the basic local action of the theory and in Section 3, we provide a review of FQFT as a perturbative quantum field theory. The nonlocal quantum behavior of the theory is considered in detail in Section 4, and in Section 5, we discuss a possible exerimental test of the onset of nonlocality by detecting CPT asymmetries. In Sections 6 and 7, we develop the formalism for Yang-Mills gauge theory and quantum gravity. In Section 8, we turn our attention to the resolution of the gauge hierarchy problem in the Higgs sector, while in Section 9, we analyze the results of gluon and gravitational vacuum polarization calculations. In Section 10, we use FQFT quantum gravity to resolve the cosmological constant problem and in Section 11, we end with concluding remarks. ## 2 The Action We shall begin with the four-dimensional action $$W=W_{\mathrm{grav}}+W_{YM}+W_\mathrm{H}+W_{\mathrm{Dirac}}+W_M,$$ (3) where $$W_{\mathrm{grav}}=\frac{2}{\kappa ^2}d^4x\sqrt{g}(R+2\lambda ),$$ (4) $$W_{\mathrm{YM}}=\frac{1}{4}d^4x\sqrt{g}\mathrm{Tr}(F^2),$$ (5) $$W_\mathrm{H}=\frac{1}{2}d^4x\sqrt{g}[D_\mu \varphi ^iD^\mu \varphi ^i+V(\varphi ^2)],$$ (6) $$W_{\mathrm{Dirac}}=\frac{1}{2}d^4x\sqrt{g}\overline{\psi }\gamma ^ae_a^\mu [_\mu \psi \omega _\mu \psi 𝒟(A_{i\mu })\psi ]+h.c.$$ (7) Here, we use the notation: $`\mu ,\nu =0,1,2,3`$, $`g=\mathrm{det}(g_{\mu \nu })`$ and the metric signature of Minkowski spacetime is $`\eta _{\mu \nu }=\mathrm{diag}(1,+1,+1,+1)`$. The Riemann tensor is defined such that $$R_{}^{\lambda }{}_{\mu \nu \rho }{}^{}=_\rho \mathrm{\Gamma }_{\mu \nu }^{}{}_{}{}^{\lambda }_\nu \mathrm{\Gamma }_{\mu \rho }^{}{}_{}{}^{\lambda }+\mathrm{\Gamma }_{\mu \nu }^{}{}_{}{}^{\alpha }\mathrm{\Gamma }_{\rho \alpha }^{}{}_{}{}^{\lambda }\mathrm{\Gamma }_{\mu \rho }^{}{}_{}{}^{\alpha }\mathrm{\Gamma }_{\nu \alpha }^{}{}_{}{}^{\lambda }.$$ (8) Moreover, h.c. denotes the Hermitian conjugate, $`\overline{\psi }=\psi ^{}\gamma ^0`$, and $`e_a^\mu `$ is a vierbein, related to the metric by $$g_{\mu \nu }=\eta _{ab}e_\mu ^ae_\nu ^b,$$ (9) where $`\eta _{ab}`$ is the four-dimensional Minkowski metric tensor associated with the flat tangent space with indices a,b,c… Moreover, $`F^2=F_{i\mu \nu }F^{i\mu \nu }`$, $`R`$ denotes the scalar curvature, $`\lambda `$ is the cosmological constant and $$F_{i\mu \nu }=_\nu A_{i\mu }_\mu A_{i\nu }ef_{ikl}A_{k\mu }A_{l\nu },$$ (10) where $`A_{i\mu }`$ are the gauge fields of the Yang-Mills group with generators $`f_{ikl}`$, $`e`$ is the coupling constant and $`\kappa ^2=32\pi G`$ with $`c=1`$. We denote by $`D_\mu `$ the covariant derivative operator $$D_\mu \varphi ^i=_\mu \varphi ^i+ef^{ikl}A_\mu ^k\varphi ^l.$$ (11) The Higgs potential $`V(\varphi ^2)`$ is of the form leading to spontaneous symmetry breaking $$V(\varphi ^2)=\frac{1}{4}g(\varphi ^i\varphi ^iK^2)^2+V_0,$$ (12) where $`V_0`$ is an adjustable constant and the coupling constant $`g>0`$. The spinor field is minimally coupled to the gauge potential $`A_{i\mu }`$, and $`𝒟`$ is a matrix representation of the gauge group $`SO(3,1)`$. The spin connection $`\omega _\mu `$ is $$\omega _\mu =\frac{1}{2}\omega _{\mu ab}\mathrm{\Sigma }^{ab},$$ (13) where $`\mathrm{\Sigma }^{ab}=\frac{1}{4}[\gamma ^a,\gamma ^b]`$ is the spinor matrix associated with the Lorentz algebra $`SO(3,1)`$. The components $`\omega _{\mu ab}`$ satisfy $$_\mu e_a^\sigma +\mathrm{\Gamma }_{\mu \nu }^{}{}_{}{}^{\sigma }e_a^\nu \omega _{\mu a}^{}{}_{}{}^{\rho }e_\rho ^\sigma =0,$$ (14) where $`\mathrm{\Gamma }_{\mu \nu }^{}{}_{}{}^{\sigma }`$ is the Christoffel symbol. The field equations for the gravity-Yang-Mills-Higgs-Dirac sector are $$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R\lambda g_{\mu \nu }=\frac{1}{4}\kappa ^2T_{\mu \nu },$$ (15) $$g^{\rho \mu }_\rho F_{\mu \nu }^i=g^{\rho \mu }(_\rho F_{\mu \nu }^i\mathrm{\Gamma }_{\rho \mu }^{}{}_{}{}^{\sigma }F_{\sigma \nu }^i$$ $$\mathrm{\Gamma }_{\rho \nu }^{}{}_{}{}^{\sigma }F_{\mu \sigma }^i+[A_\rho ,F_{\mu \nu }]^i)=0,$$ (16) $$\frac{1}{\sqrt{g}}D_\mu [\sqrt{g}g^{\mu \nu }D_\nu \varphi ^i]=\left(\frac{V}{\varphi ^2}\right)\varphi ^i,$$ (17) $$\gamma ^ae_a^\mu [_\mu \omega _\mu 𝒟(A_\mu )]\psi =0.$$ (18) The energy-momentum tensor is $$T_{\mu \nu }=T_{\mu \nu }^{\mathrm{YMH}}+T_{\mu \nu }^{\mathrm{Dirac}}+T_{\mu \nu }^\mathrm{M},$$ (19) where $$T_{\mu \nu }^{\mathrm{YMH}}=\mathrm{Tr}(F_{\mu \sigma }F_\nu ^\sigma )+D_\mu \varphi ^iD_\nu \varphi ^i$$ $$\frac{1}{2}g_{\mu \nu }\left[\frac{1}{2}\mathrm{Tr}(F^2)+D_\sigma \varphi ^iD^\sigma \varphi ^i+V(\varphi ^2)\right],$$ (20) $$T_{\mu \nu }^{\mathrm{Dirac}}=\overline{\psi }\gamma _\mu [_\nu \omega _\nu 𝒟(A_{i\nu })]\psi ,$$ (21) and $`T_{\mu \nu }^\mathrm{M}`$ is the energy-momentum tensor of non-field matter. ## 3 Finite Quantum Field Theory Formalism An important development in nonlocal FQFT was the discovery that gauge invariance and unitarity can be restored by adding series of higher interactions. The resulting theory possesses a nonlinear, field representation dependent gauge invariance which agrees with the original local symmetry on shell but is larger off shell. Quantization is performed in the functional formalism using an analytic and convergent measure factor which retains invariance under the new symmetry. An explicit calculation was made of the measure factor in QED , and it was obtained to lowest order in Yang-Mills theory . Kleppe and Woodard obtained an ansatz based on the derived dimensionally regulated result when $`\mathrm{\Lambda }\mathrm{}`$, which was conjectured to lead to a general functional measure factor in FQFT gauge theories. In contrast to string theory, we can achieve a genuine quantum field theory, which allows vertex operators to be taken off the mass shell. The finiteness draws from the fact that factors of $`\mathrm{exp}[𝒦(p^2)/2\mathrm{\Lambda }^2]`$ are attached to propagators which suppress any ultraviolet divergences in Euclidean momentum space, where $`\mathrm{\Lambda }`$ is an energy scale factor. An important feature of FQFT is that only the quantum loop graphs have nonlocal properties; the classical tree graph theory retains full causal and local behavior. A convenient formalism which makes the FQFT construction transparent is based on shadow fields . We shall consider the 4-dimensional spacetime to be approximately flat Minkowski spacetime. Let us denote by $`f_i`$ a generic local field and write the standard local action as $$W[f]=W_F[f]+W_I[f],$$ (22) where $`W_F`$ and $`W_I`$ denote the free part and the interaction part of the action, respectively, and $$W_F[f]=\frac{1}{2}d^4xf_i𝒦_{ij}f_j.$$ (23) In a gauge theory $`W`$ would be the Becchi, Rouet, Stora, Tyutin (BRST) gauge-fixed action including ghost fields in the invariant action required to fix the gauge. The kinetic operator $`𝒦`$ is fixed by defining a Lorentz-invariant distribution operator $$\mathrm{exp}\left(\frac{𝒦}{2\mathrm{\Lambda }^2}\right)$$ (24) and the shadow operator: $$𝒪^1=\frac{𝒦}{^21}.$$ (25) Every local field $`f_i`$ has an auxiliary counterpart field $`h_i`$, and they are used to form a new action $$W[f,h]W_F[\widehat{f}]P[h]+W_I[f+h],$$ (26) where $$\widehat{f}=^1f,P[h]=\frac{1}{2}d^4xh_i𝒪_{ij}^1h_j.$$ By iterating the equation $$h_i=𝒪_{ij}\frac{\delta W_I[f+h]}{\delta h_j}$$ (27) the shadow fields can be determined as functions, and the regulated action is derived from $$\widehat{W}[f]=W[f,h(f)].$$ (28) We recover the original local action when we take the limit $`\mathrm{\Lambda }\mathrm{}`$ and $`\widehat{f}f,h(f)0`$. The expression (27) can be developed into a series expansion for $`h_i[f]`$. The regularized action is found by substituting into it the classical solution $`h_i[f]`$. Expanding $`\widehat{W}`$ in powers of $`f`$ gives the kinetic term $`W_F[\widehat{f}]`$, together with an infinite series of interaction terms the first of which is just $`W_I[f]`$. Since $`𝒪`$ is an entire function of $`𝒦`$ the higher interactions are also entire functions of $`𝒦`$. This is important for preserving unitarity. Quantization is performed using the definition $$0\|T^{}(O[f])\|0_{}=[Df]\mu [f](\mathrm{gauge}\mathrm{fixing})O[\widehat{f}]\mathrm{exp}(i\widehat{W}[f]).$$ (29) On the left-hand side we have the regulated vacuum expectation value of the $`T^{}`$-ordered product of an arbitrary operator $`O[f]`$ formed from the local fields $`f_i`$. The subscript $``$ signifies that a regulating Lorentz distribution has been used. Moreover, $`\mu [f]`$ is a measure factor and there is a gauge fixing factor, both of which are needed to maintain perturbative unitarity in gauge theories. The new Feynman rules for FQFT are obtained as follows: The vertices remain unchanged within the regularized action, but every leg of a diagram is connected either to a regularized propagator, $$\frac{i^2}{𝒦+iϵ}=i_1^{\mathrm{}}\frac{d\tau }{\mathrm{\Lambda }^2}\mathrm{exp}\left(\tau \frac{𝒦}{\mathrm{\Lambda }^2}\right),$$ (30) or to a shadow propagator, $$i𝒪=\frac{i(1^2)}{𝒦}=i_0^1\frac{d\tau }{\mathrm{\Lambda }^2}\mathrm{exp}\left(\tau \frac{𝒦}{\mathrm{\Lambda }^2}\right).$$ (31) We shall also attach a factor $`(p^2)`$ to every external leg connected to a loop, which is unity on shell. The formalism is set up in Minkowski spacetime and loop integrals are formally defined in Euclidean space by performing a Wick rotation. This facilitates the analytic continuation; the whole formalism could from the outset be developed in Euclidean space. In FQFT renormalization is carried out as in any other field theory. The bare parameters are calculated from the renormalized ones and $`\mathrm{\Lambda }`$, such that the limit $`\mathrm{\Lambda }\mathrm{}`$ is finite for all noncoincident Green’s functions, and the bare parameters are those of the local theory. The regularizing interactions are determined by the local operators. The regulating Lorentz distribution function $``$ must be chosen to perform an explicit calculation in perturbation theory. We do not know the unique choice of $``$. However, once a choice for the function is made, then the theory and the perturbative calculations are uniquely fixed. A standard choice in early FQFT papers is : $$_m=\mathrm{exp}\left(\frac{^2m^2}{2\mathrm{\Lambda }^2}\right).$$ (32) An explicit construction for QED was given using the Cutkosky rules as applied to FQFT whose propagators have poles only where $`𝒦=0`$ and whose vertices are entire functions of $`𝒦`$. The regulated action $`\widehat{W}[f]`$ satisfies these requirements which guarantees unitarity on the physical space of states. The local action is gauge fixed and then a regularization is performed on the BRST theory. The infinitesimal transformation $$\delta f_i=T_i(f)$$ (33) generates a symmetry of $`W[f]`$, and the infinitesimal transformation $$\widehat{\delta }f_i=_{ij}^2T_j(f+h[f])$$ (34) generates a symmetry of the regulated action $`\widehat{W}[f]`$. To see this consider the transformations $$\delta f_i=_{ij}^2T_j[f+h],\delta h_i=(1^2)_{ij}T_j[f+h].$$ (35) Adding these two transformations gives $$\delta (f+h)_i=T_i[f+h].$$ (36) Then, (35) is a symmetry of the action $`W[f,h]`$. We have $$\delta W[f,h]=d^4x\left\{(f_i+h_i)𝒦_{ij}T_j[f+h]\right\}$$ $$+\frac{\delta W_I[f+h]}{\delta f_i}T_i[f+h]=\delta W[f+h].$$ (37) It follows that $`\delta W[f,h]=0`$ is a consequence of the assumed invariance $`\delta W[f+h]=0`$. Now $`\widehat{W}[f]`$ is invariant under (34), for we have $$\widehat{\delta }h_i[f]=(1_{ij}^2)T_j[f+h[f]]L_{ij}[f+h[f]]\frac{\delta T_j}{\delta f_k}[f+h[f]]_{kl}^2\frac{\delta \widehat{W}[f]}{\delta f_l},$$ where $$L_{ij}^1=𝒪_{ij}^1\frac{\delta ^2W_I[f]}{\delta f_i\delta f_j}.$$ It follows that FQFT regularization preserves all continuous symmetries including supersymmetry. The quantum theory will preserve symmetries provided a suitable measure factor can be found such that $$\widehat{\delta }([Df]\mu [f])=0.$$ (38) Moreover, the interaction vertices of the measure factor must be entire functions of the operator $`𝒦`$ and they must not destroy the FQFT finiteness. In FQFT tree order, Green’s functions remain local except for external lines which are unity on shell. It follows immediately that since on shell tree amplitudes are unchanged by the regularization, $`\widehat{W}`$ preserves all symmetries of $`W`$ on shell. Also all loops contain at least one regularizing propagator and therefore are ultraviolet finite. Shadow fields are eliminated at the classical level, for functionally integrating over them would produce divergences from shadow loops. Since shadow field propagators do not contain any poles there is no need to quantize the shadow fields. In FQFT, the on shell tree amplitudes agree with the local, unregulated action, while the loop amplitudes disagree. This seems to contradict the Feynman tree theorem , which states that loop amplitudes of local field theory can be expressed as sums of integrals of tree diagrams. If two local theories agree at the tree level, then the loop amplitudes agree as well. However, the tree theorem does not apply to nonlocal field theories. The tree theorem is proved by using the propagator relation $$D_F=D_R+D^+$$ (39) to expand the Feynman propagator $`D_F`$ into a series in the on shell propagator $`D^+`$. This decomposes all terms with even one $`D^+`$ into trees. The term with no $`D^+s`$ is a loop formed with the retarded propagator and vanishes for local interactions. But for nonlocal interactions, this term generally survives and new physical effects occur in loop amplitudes, which cannot be predicted from the local on shell tree graphs. ## 4 Quantum Nonlocal Behavior in FQFT It appears on general grounds that interacting strings are nonlocal . Nonlocality in open string theory can arise from the non-commutativity of spacetime coordinates $$[x^\mu ,x^\nu ]=i\theta ϵ^{\mu \nu }.$$ (40) This nonlocality in string theory is closely associated with the string uncertainty principle $$\mathrm{\Delta }x\mathrm{\Delta }t\alpha ^{}.$$ (41) Nonlocality has also been associated with the formation of black hole horizons and the lack of commutativity of spatial coordinates and time . The horizon responds to incoming matter before it comes in. Kapustin has recently shown that LSTs are quasi-local field theories whose infrared limit can approach local field theories in the large . The exponential growth of Wightman functions (Green’s functions) in momentum space is a characteristic feature of nonlocal field theories. The corresponding test functions in x-space are real analytic and cannot possess compact support. The Wightman functions or vacuum expectations values of products of field operators $`\varphi (p)`$: $$W_n(q_1,\mathrm{},q_{n1})=0\|\varphi (q_1)\varphi (q_2)\mathrm{}\varphi (q_{n1})\|0,$$ (42) grow exponentially with momenta for nonlocal field theories. By the positivity of energy, $`W_n`$ vanishes when any of its arguments are outside the forward light cone. Inside the forward light cone $`W_n`$ is bounded by $$\mathrm{exp}[\mathrm{}(\|q_1\|+\mathrm{}\|q_{n1}\|)],$$ (43) where $`\|q\|=\sqrt{q^2}`$ and $`\mathrm{}`$ is a length scale. In the case of LST models, the length scale is given by $`\mathrm{}\sqrt{N}/M_s`$ where $`N`$ is the number of coincident five-branes. Jaffe defined a test function space $`\stackrel{~}{S}_g`$ in momentum space, which is convenient to use when discussing nonlocal field theories, in which all functions are infinitely differentiable and for which all the norms are finite. Given a positive function $`g(t)`$ which is entire, Jaffe showed that if $`g(t)`$ satisfies $$_0^{\mathrm{}}𝑑t\frac{\mathrm{ln}g(t^2)}{1+t^2}<\mathrm{},$$ (44) then the Fourier transform of $`\stackrel{~}{S}_g`$ has functions with compact support, strictly local quantum fields can be defined and a local quantum field theory can be formulated. On the other hand, if (44) is not satisfied, then there are no test functions with compact support and we have a nonlocal quantum field theory. Our choice of the entire function $`(p^2)`$ in the factor, Eq. (24), will not lead to a test function space that satisfies the condition (44). We can choose a function $`(p^2)`$ which will provide a test function space that leads to a quasi-local quantum field theory, as defined by Kapustin, and in the earlier work by Iofa and Fainberg. In the present work, we have chosen $`𝒦(p^2)=(p^2+m^2)`$, because it leads to a simplification of calculations in perturbation theory. But this is purely a technical issue, and we can certainly adopt entire functions $`𝒦(p^2)`$ which lead to quasi-local field operators, which only violate locality at short distances. The commutator for a scalar field operator $`\varphi (x)`$: $$[\varphi (x),\varphi (y)]=W_2(xy)W_2(yx)$$ (45) in our theory will not vanish outside the light cone for space-like separations $`(xy)^2>0`$. Indeed, it will satisfy $$[\varphi (x),\varphi (y)]\delta [(xy)^2\mathrm{}^2]\mathrm{sign}(x_0y_0).$$ (46) In FQFT, it can be argued that the extended objects that replace point particles (the latter are obtained in the limit $`\mathrm{\Lambda }\mathrm{}`$) cannot be probed because of a Heisenberg uncertainty type of argument. The FQFT nonlocality only occurs at the quantum loop level, so there is no noncausal classical behavior. In FQFT the strength of a signal propagated over an invariant interval $`\mathrm{}^2`$ outside the light cone would be suppressed by a factor $`\mathrm{exp}(\mathrm{}^2\mathrm{\Lambda }^2)`$. Nonlocal field theories can possess non-perturbative instabilities. These instabilities arise because of extra canonical degrees of freedom associated with higher time derivatives. If a Lagrangian contains up to $`N`$ time derivatives, then the associated Hamiltonian is linear in $`N1`$ of the corresponding canonical variables and extra canonical degrees of freedom will be generated by the higher time derivatives. A nonlocal theory can be viewed as the limit $`N\mathrm{}`$ of an Nth derivative Lagrangian. Unless the dependence on the extra solutions is arbitrarily choppy in the limit, then the higher derivative limit will produce instabilities . The condition for the smoothness of the extra solutions is that no invertible field redefinition exists which maps the nonlocal field equations into the local ones. String theory does satisfy this smoothness condition as can be seen by inspection of the S-matrix tree graphs. In FQFT the tree amplitudes agree with those of the local theory, so the smoothness condition is not obeyed. It was proved by Kleppe and Woodard that the solutions of the nonlocal field equations in FQFT are in one-to-one correspondence with those of the original local theory. The relation for a generic field $`v_i`$ is $$v_i^{\mathrm{nonlocal}}=_{ij}^2v_j^{\mathrm{local}}.$$ (47) Also the actions satisfy $$W[v]=\widehat{W}[^2v].$$ (48) Thus, there are no extra classical solutions. The solutions of the regularized nonlocal Euler-Lagrange equations are in one-to-one correspondence with those of the local action. It follows that the regularized nonlocal FQFT is free of higher derivative solutions, so FQFT can be a stable theory. Since only the quantum loop graphs in the nonlocal FQFT differ from the local field theory, then FQFT can be viewed as a non-canonical quantization of fields which obey the local equations of motion. Provided the functional quantization in FQFT is successful, then the theory does maintain perturbative unitarity. ## 5 Experimental Tests of Nonlocality In order to solve the Higgs and cosmological constant radiative stability, hierarchy problems, we have relaxed the assumption of microcausal locality in our FQFT. A scale of nonlocality $`\mathrm{\Lambda }`$ is set for the graviton ($`\mathrm{\Lambda }_G10^3`$ eV), the Higgs particle ($`\mathrm{\Lambda }_H1`$ TeV) and the standard model gauge particles ($`\mathrm{\Lambda }_{GP}1`$ TeV). We do not understand the fundamental physics which is the source of these nonlocality scales but, as we shall see, given these scales we can potentially solve the radiative stability problems in a fully gauge invariant, finite and unitary fashion, including the gravitational stability of the cosmological constant. Supersymmetry and technicolor models have been proposed to solve the Higgs gauge hierarchy problem. The mass scales for supersymmetry are set ’by hand’, so to speak, according to when we expect supersymmetry breaking to set in, allowing super-partners to be detected, and when technicolor fermions form condensates, allowing us to detect technicolor particles. No known fundamental physics tells us what these mass scales are. We can only guess their magnitude above certain obvious intermediate energy bounds. Experiments already tend to disfavour techniclor models, and if the large hadron colliders do not detect super-partners below 2-3 TeV, then this would kill the possibility of using supersymmetric models to explain the radiative stability of the Higgs particle. Can we experimentally detect the onset of nonlocality? We could do this by checking dispersion relations for scattering amplitudes at high energies. We expect that the non-vanishing of commutators of field operators outside the light cone will decrease exponentially with the spacelike distance, so violations of nonlocality will be small, and changes of analyticity of the scattering amplitudes from the standard microcausal analyticity properties will correspondingly be small. Another possible signature of nonlocality is a violation of CPT invariance. This is a fundamental theorem of local quantum field theory . There have been suggestions that CPT invariance could be broken in quantum gravity . Moreover, there have been several studies of meson decays with the prospects of detecting CPT invariance breaking at K-meson and B-meson factories . Let us investigate how CPT invariance could be violated by nonlocality. Consider a complex, nonlocal Heisenberg-picture scalar field operator $`\mathrm{\Phi }(x)`$. The Källen-Lehmann representation is given by the vacuum expectation value $$0\|\mathrm{\Phi }(x)\mathrm{\Phi }^{}(y)\|0=_0^{\mathrm{}}𝑑\mu ^2\rho (\mu ^2)\stackrel{~}{\mathrm{\Delta }}_+(xy;\mu ^2),$$ (49) where $$\stackrel{~}{\mathrm{\Delta }}_+(xy;\mu ^2)=\frac{1}{(2\pi )^3}_0^{\mathrm{}}d^4p\mathrm{exp}[ip(xy)]\mathrm{\Pi }(xy)\theta (p^0)\delta (p^2+\mu ^2),$$ (50) and $`\mathrm{\Pi }(xy)`$ is an entire analytic function with $`\mathrm{\Pi }(x)>0`$ for real $`x`$ . The spectral function $`\rho `$ is defined by $$\underset{n}{}\delta ^4(pp_n)\|0\|\mathrm{\Phi }(0)\|n\|^2=\frac{1}{(2\pi )^3}\theta (p^0)\rho (p^2)$$ (51) with $`\rho (p^2)=0`$ for $`p^2<0`$. We also have $$0\|\mathrm{\Phi }^{}(y)\mathrm{\Phi }(x)\|0=_0^{\mathrm{}}𝑑\mu ^2\overline{\rho }(\mu ^2)\stackrel{~}{\mathrm{\Delta }}_+(yx;\mu ^2),$$ (52) where $$\underset{n}{}\delta ^4(pp_n)\|n\|\mathrm{\Phi }^{}(0)\|0\|^2=\frac{1}{(2\pi )^3}\theta (p^0)\overline{\rho }(p^2).$$ (53) Let us define $$\alpha (\mu ^2)=\rho (\mu ^2)\overline{\rho }(\mu ^2).$$ (54) The vacuum expectation value of the commutator is $$0\|[\mathrm{\Phi }(x),\mathrm{\Phi }^{}(y)]\|0$$ $$=_0^{\mathrm{}}𝑑\mu ^2\{\rho (\mu ^2)[\stackrel{~}{\mathrm{\Delta }}_+(xy;\mu ^2)\stackrel{~}{\mathrm{\Delta }}_+(yx;\mu ^2)]+\alpha (\mu ^2)\stackrel{~}{\mathrm{\Delta }}_+(yx;\mu ^2)\}.$$ (55) For spacelike separations $`(xy)^2>0`$, the function $`\stackrel{~}{\mathrm{\Delta }}_+(xy;\mu ^2)=\stackrel{~}{\mathrm{\Delta }}_+(yx;\mu ^2)`$ and it does not vanish. For (5) to vanish for spacelike separations, we must have $`\alpha (\mu ^2)=0`$. This is a nonperturbative proof of the CPT theorem, for states with $`p^2=\mu ^2`$ have the quantum numbers of the particle associated with $`\mathrm{\Phi }`$, and there must be corresponding states with $`p^2=\mu ^2`$ that have the quantum numbers of the anti-particle described by the operator $`\mathrm{\Phi }^{}`$ . For strictly local field operators $`\mathrm{\Phi }`$, the commutator $$[\mathrm{\Phi }(x),\mathrm{\Phi }(y)]=0$$ (56) for spacelike separations $`(xy)^2>0`$. However, we assumed that the $`\mathrm{\Phi }(x)`$ were nonlocal field operators, so there will be a violation of the CPT theorem when $`\alpha (\mu ^2)0`$, and we have for spacelike separation $$0\|[\mathrm{\Phi }(x),\mathrm{\Phi }^{}(y)]\|0=_0^{\mathrm{}}𝑑\mu ^2\alpha (\mu ^2)\stackrel{~}{\mathrm{\Delta }}_+(yx;\mu ^2).$$ (57) The vacuum expectation value of the time-ordered product is $$0\|T\left\{\mathrm{\Phi }(x)\mathrm{\Phi }^{}(y)\right\}\|0$$ $$=i_0^{\mathrm{}}𝑑\mu ^2\rho (\mu ^2)\stackrel{~}{\mathrm{\Delta }}_F(xy;\mu ^2)+i_0^{\mathrm{}}𝑑\mu ^2\alpha (\mu ^2)\theta (y^0x^0)\stackrel{~}{\mathrm{\Delta }}_+(yx;\mu ^2),$$ (58) where $`\mathrm{\Delta }_F`$ is the Feynman propagator $$i\stackrel{~}{\mathrm{\Delta }}_F(xy;\mu ^2)=\theta (x^0y^0)\stackrel{~}{\mathrm{\Delta }}_+(xy;\mu ^2)\theta (y^0x^0)\stackrel{~}{\mathrm{\Delta }}_+(yx;\mu ^2).$$ (59) For a nonlocal interaction $$V_{\mathrm{NL}}=d^3x_{\mathrm{NL}}(\stackrel{}{x},0),$$ (60) the commutator $`[CPT,V_{\mathrm{NL}}]`$ will not in general vanish. The masses and decay rates of particles and anti-particles will not be equal for CPT invariance violating processes. For the discrete symmetries of nature, violations have been observed for C, P and the combined CP symmetries. Two types of CP symmetry violation have been observed for K-mesons. An active pursuit to detect CPT asymmetries in meson decays is presently underway. ## 6 Finite Quantum Yang-Mills Theory Let us now review the finite quantization of the Yang-Mills sector in four-dimensional Minkowski flat space. The gauge field strength $`F_{i\mu \nu }`$ is invariant under the familiar transformations: $$\delta A_{i\mu }=_\mu \theta _i+ef_{ikl}A_{k\mu }\theta _l.$$ (61) To regularize the Yang-Mills sector, we identify the kinetic operator $$𝒦_{ik}^{\mu \nu }=\delta _{ik}(^2\eta ^{\mu \nu }^\mu ^\nu ).$$ The regularized action is given by $$\widehat{W}_{YM}[A]=\frac{1}{2}d^4x\left\{\widehat{A}_{i\mu }𝒦_{ik}^{\mu \nu }\widehat{A}_{k\nu }B_{i\mu }[A](𝒪_{ik}^{\mu \nu })^1B_{k\nu }[A]\right\}$$ $$+W_{YM}^I[A+B[A]],$$ (62) where $`B_{i\mu }`$ is the Yang-Mills shadow field, which satisfies the expansion $$B_i^\mu [A]=𝒪_{ik}^{\mu \nu }\frac{\delta W_{YM}^I[A+B]}{\delta B_k^\nu }$$ $$=𝒪_{ik}^{\mu \nu }ef_{klm}[A_{\nu l}_\sigma A_m^\sigma +A_{l\sigma }_\nu A_m^\sigma 2A_{l\sigma }^\sigma A_{\nu _m}]+O(e^2A^3).$$ (63) The regularized gauge symmetry transformation is $$\widehat{\delta }_\theta A_i^\mu =(_{ik}^{2\mu \nu })\left\{_\nu \theta _k+ef_{klm}(A_{l\nu }+B_{l\nu }[A])\theta _m\right\}.$$ The extended gauge transformation is neither linear nor local. We functionally quantize the Yang-Mills sector using $$0\|T^{}(O[A])\|0_{}=[DA]\mu [A](\mathrm{gauge}\mathrm{fixing})O[\widehat{A}]\mathrm{exp}(i\widehat{W}_{\mathrm{YM}}[A]).$$ (64) To fix the gauge we use Becchi-Rouet-Stora-Tyutin (BRST) invariance. The ghost structure of the BRST action comes from exponentiating the Faddeev-Popov determinant. Since the FQFT algebra fails to close off-shell, we need to introduce higher ghost terms into both the action and the BRST transformation. In Feynman gauge, the local BRST Lagrangian is $$_{YMBRST}=\frac{1}{2}_\mu A_{i\nu }^\mu A_i^\nu ^\mu \overline{\eta _i}_\mu \eta _i+ef_{ikl}^\mu \overline{\eta }_iA_{k\mu }\eta _l$$ $$+ef_{ikl}_\mu A_{i\nu }A_k^\mu A_l^\nu \frac{1}{4}e^2f_{ikl}f_{lmn}A_{i\mu }A_{k\nu }A_m^\mu A_n^\nu .$$ (65) It is invariant under the global symmetry transformation: $$\delta A_{i\mu }=(_\mu \eta _ief_{ikl}A_{k\mu }\eta _l)\delta \zeta ,$$ $$\delta \eta _i=\frac{1}{2}ef_{ikl}\eta _k\eta _l\delta \zeta ,$$ $$\delta \overline{\eta }_i=_\mu A_i^\mu \delta \zeta ,$$ where $`\zeta `$ is a constant anticommuting c-number. The gluon and ghost kinetic operators are $$𝒦_{ik}^{\mu \nu }=\delta _{ik}\eta ^{\mu \nu }^2,𝒦_{ik}=\delta _{ik}^2,$$ (66) The gluon propagator and the shadow gluon propagator are given by $$D_{ik}^{\mu \nu }(p^2)=\frac{i\delta _{ik}\eta ^{\mu \nu }}{p^2iϵ}\mathrm{exp}\left(p^2/\mathrm{\Lambda }_{\mathrm{YM}}^2\right),$$ (67) $$D_{ik}^{\mathrm{shad}\mu \nu }(p^2)=\frac{i\delta _{ik}\eta ^{\mu \nu }}{p^2iϵ}\left[1\mathrm{exp}\left(p^2/\mathrm{\Lambda }_{\mathrm{YM}}^2\right)\right],$$ (68) where $`\mathrm{\Lambda }_{YM}`$ denotes the FQFT Yang-Mills energy scale. The regularized BRST action is $$\widehat{W}_{YM}[A,\overline{\eta },\eta ]=d^4x\{\frac{1}{2}_\nu \widehat{A}_{i\mu }^\nu \widehat{A}_i^\mu \frac{1}{2}B_{i\mu }\overline{𝒪}^1B_i^\mu $$ $$^\mu \widehat{\overline{\eta }}_i_\mu \widehat{\eta }_i\overline{\chi }_i\overline{𝒪}^1\chi _i\}+W^I_{\mathrm{YM}}[A+B,\overline{\eta }+\overline{\chi },\eta +\chi ],$$ (69) where $`\chi `$ is the ghost shadow field. The regularizing, nonlocal BRST symmetry transformation is $$\widehat{\delta }A_{i\mu }=\overline{}^2\left\{(_\mu \eta _i+_\mu \chi _i)ef_{ikl}(A_{k\mu }+B_{k\mu })(\eta _l+\chi _l)\right\}\delta \zeta ,$$ $$\widehat{\delta }\eta _i=\frac{1}{2}ef_{ikl}\overline{}^2(\eta _k+\chi _k)(\eta _l+\chi _l)\delta \zeta ,$$ $$\widehat{\delta }\overline{\eta }_i=\overline{}^2(_\mu A_i^\mu +_\mu B_i^\mu )\delta \zeta .$$ (70) The full functional, gauge fixed quantization is now given by $$0\|T^{}(O[A,\overline{\eta },\eta ])\|0_{}=[DA][D\overline{\eta }][D\eta ]\mu [A,\overline{\eta },\eta ]O[\widehat{A},\widehat{\overline{\eta }},\widehat{\eta }]$$ $$\times \mathrm{exp}(i\widehat{W}_{\mathrm{YM}}[A,\overline{\eta },\eta ]).$$ (71) Kleppe and Woodard have obtained the invariant measure factor for the regularized Yang-Mills sector to first order in the coupling constant $`e`$: $$\mathrm{ln}(\mu [A,\overline{\eta },\eta ])=\frac{1}{2}e^2f_{ilm}f_{klm}d^4xA_{i\mu }A_k^\mu +O(e^3),$$ (72) where $$=\frac{1}{16\pi ^2}_0^1𝑑\tau \frac{\mathrm{\Lambda }^2}{(\tau +1)^2}\mathrm{exp}\left(\frac{\tau }{\tau +1}\frac{^2}{\mathrm{\Lambda }^2}\right)\left\{\frac{2+6\tau }{\tau +1}3\right\}.$$ (73) The existence of a suitable invariant measure factor implies that the necessary Slavnov-Taylor identities also exist. ## 7 Finite Perturbative Quantum Gravity As is well know, the problem with perturbative quantum gravity based on a point-like graviton and a local field theory formalism is that the theory is not renormalizable . Due to the Gauss-Bonnet theorem, it can be shown that the one-loop graviton calculation is renormalizable but two-loop is not . Moreover, gravity-matter interactions are not renormalizable at any loop order. We shall now formulate the gravitational sector in more detail as a FQFT. This problem has been considered previously in the context of four-dimensional GR . We shall expand the gravity sector about flat Minkowski spacetime. In fact, FQFT can be formulated as a perturbative theory by expanding around any fixed, classical metric background $$g_{\mu \nu }=\overline{g}_{\mu \nu }+h_{\mu \nu },$$ (74) where $`\overline{g}_{\mu \nu }`$ is any smooth background metric field, e.g. a de Sitter spacetime metric. For the sake of simplicity, we shall only consider expansions about flat spacetime. Since the gravitational field is weak up to the Planck energy scale, this expansion is considered justified; even at the standard model energy scale $`E_{\mathrm{SM}}10^2`$ GeV, we have $`\kappa ^2E_{\mathrm{SM}}^210^{33}`$. Also, at these energy scales the curvature of spacetime is very small. However, if we wish to include the cosmological constant $`\lambda `$, then we cannot strictly speaking expand about flat spacetime, because such an expansion of the Einstein field equations will lead to the result that $`\lambda =0`$. This is to be expected, because the cosmological constant produces a curved spacetime even when the energy-momentum tensor $`T_{\mu \nu }=0`$. Therefore, we should in this case use the expansion (74). But for energy scales encountered in particle physics, the curvature is very small, so we can approximate the perturbation caculation by using the flat spacetime expansion and trust that the results are valid in general for curved spacetime backgrounds including the cosmological constant. As in ref. , we will regularize the GR equations using the covariant shadow field formalism. Let us define $`𝐠^{\mu \nu }=\sqrt{g}g^{\mu \nu }`$. It can be shown that $`\sqrt{g}=\sqrt{𝐠}`$, where $`𝐠=\mathrm{det}(𝐠^{\mu \nu })`$ and $`_\rho 𝐠=𝐠_{\alpha \beta }_\rho 𝐠^{\alpha \beta }𝐠`$. We can then write the local gravitational action $`W_{\mathrm{grav}}`$ in the form : $$W_{\mathrm{grav}}=d^4x_{\mathrm{grav}}=\frac{1}{2\kappa ^2}d^4x[(𝐠^{\rho \sigma }𝐠_{\lambda \mu }𝐠_{\kappa \nu }$$ $$\frac{1}{2}𝐠^{\rho \sigma }𝐠_{\mu \kappa }𝐠_{\lambda \nu }2\delta _\kappa ^\sigma \delta _\lambda ^\rho 𝐠_{\mu \nu })_\rho 𝐠^{\mu \kappa }_\sigma 𝐠^{\lambda \nu }$$ $$\frac{1}{\alpha \kappa ^2}_\mu 𝐠^{\mu \nu }_\kappa 𝐠^{\kappa \lambda }\eta _{\nu \lambda }+\overline{C}^\nu ^\mu X_{\mu \nu \lambda }C^\lambda ],$$ (75) where we have added a gauge fixing term with the parameter $`\alpha `$, $`C^\mu `$ is the Fadeev-Popov ghost field and $`X_{\mu \nu \lambda }`$ is a differential operator. We expand the local interpolating graviton field $`𝐠^{\mu \nu }`$ as $$𝐠^{\mu \nu }=\eta ^{\mu \nu }+\kappa \gamma ^{\mu \nu }+O(\kappa ^2).$$ (76) Then, $$𝐠_{\mu \nu }=\eta _{\mu \nu }\kappa \gamma _{\mu \nu }+\kappa ^2\gamma _{\mu }^{}{}_{}{}^{\alpha }\gamma _{\alpha }^{}{}_{\nu }{}^{}+O(\kappa ^3).$$ (77) The gravitational Lagrangian density is expanded as $$_{\mathrm{grav}}=^{(0)}+\kappa ^{(1)}+\kappa ^2^{(2)}+\mathrm{}.$$ (78) We obtain $$^{(0)}=\frac{1}{2}_\sigma \gamma _{\lambda \rho }^\sigma \gamma ^{\lambda \rho }_\lambda \gamma ^{\rho \kappa }_\kappa \gamma _\rho ^\lambda \frac{1}{4}_\rho ^\rho \gamma $$ $$\frac{1}{\alpha }_\rho \gamma _\lambda ^\rho _\kappa \gamma ^{\kappa \lambda }+\overline{C}^\lambda _\sigma ^\sigma C_\lambda ,$$ (79) $$^{(1)}=\frac{1}{4}(4\gamma _{\lambda \mu }^\rho \gamma ^{\mu \kappa }_\rho \gamma _\kappa ^\lambda +2\gamma _{\mu \kappa }^\rho \gamma ^{\mu \kappa }_\rho \gamma $$ $$+2\gamma ^{\rho \sigma }_\rho \gamma _{\lambda \nu }_\sigma \gamma ^{\lambda \nu }\gamma ^{\rho \sigma }_\rho \gamma _\sigma \gamma +4\gamma _{\mu \nu }_\lambda \gamma ^{\mu \kappa }_\kappa \gamma ^{\nu \lambda })$$ $$+\overline{C}^\nu \gamma _{\kappa \mu }^\kappa ^\mu C_\nu +\overline{C}^\nu ^\mu \gamma _{\kappa \mu }^\kappa C_\nu \overline{C}^\nu ^\lambda ^\mu \gamma _{\mu \nu }C_\lambda \overline{C}^\nu ^\mu \gamma _{\mu \nu }^\lambda C_\lambda ,$$ (80) $$^{(2)}=\frac{1}{4}(4\gamma _{\kappa \alpha }\gamma ^{\alpha \nu }^\rho \gamma ^{\lambda \kappa }_\rho \gamma _{\nu \lambda }+(2\gamma _{\lambda \mu }\gamma _{\kappa \nu }\gamma _{\mu \kappa }\gamma _{\nu \lambda })^\rho \gamma ^{\mu \kappa }_\rho \gamma ^{\nu \lambda }$$ $$2\gamma _{\lambda \alpha }\gamma _\nu ^\alpha ^\rho \gamma ^{\lambda \nu }_\rho \gamma 2\gamma ^{\rho \sigma }\gamma _\nu ^\kappa _\rho \gamma _{\lambda \kappa }_\sigma \gamma ^{\nu \lambda }$$ $$+\gamma ^{\rho \sigma }\gamma ^{\nu \lambda }_\sigma \gamma _{\nu \lambda }_\rho \gamma 2\gamma _{\mu \alpha }\gamma ^{\alpha \nu }^\lambda \gamma ^{\mu \kappa }_\kappa \gamma _{\nu \lambda }),$$ (81) where $`\gamma =\gamma _{}^{\alpha }{}_{\alpha }{}^{}`$. In the limit $`\alpha \mathrm{}`$, the Lagrangian density $`_{\mathrm{grav}}`$ is invariant under the gauge transformation $$\delta \gamma _{\mu \nu }=X_{\mu \nu \lambda }\xi ^\lambda ,$$ (82) where $`\xi ^\lambda `$ is an infinitesimal vector quantity and $$X_{\mu \nu \lambda }=\kappa (_\lambda \gamma _{\mu \nu }+2\eta _{(\mu \lambda }\gamma _{\kappa \nu )}^\kappa )+(\eta _{(\mu \lambda }_{\nu )}\eta _{\mu \nu }_\lambda ).$$ (83) However, for the quantized theory it is more useful to require the BRST symmetry. We choose $`\xi ^\lambda =C^\lambda \sigma `$, where $`\sigma `$ is a global anticommuting scalar. Then, the BRST transformation is $$\delta \gamma _{\mu \nu }=X_{\mu \nu \lambda }C^\lambda \sigma ,\delta \overline{C}^\nu =_\mu \gamma ^{\mu \nu }\left(\frac{2\sigma }{\alpha }\right),\delta C_\nu =\kappa C^\mu _\mu C_\nu \sigma .$$ (84) We now substitute the operators $$\gamma _{\mu \nu }\widehat{\gamma }_{\mu \nu },C_\lambda \widehat{C}_\lambda ,\overline{C}_\nu \widehat{\overline{C}}_\nu ,$$ (85) where $$\widehat{\gamma }_{\mu \nu }=^1\gamma _{\mu \nu },\widehat{C}_\lambda =^1C_\lambda ,\widehat{\overline{C}}_\lambda =^1C_\lambda .$$ (86) As in the case of the Yang-Mills sector, the on shell propagators are unaltered from their local antecedents, while virtual particles are nonlocal. This destroys the gauge invariance of e.g. graviton-graviton scattering and requires an iteratively defined series of “stripping” vertices to ensure the decoupling of all unphysical modes. Moreover, the local gauge transformations have to be extended to nonlinear, nonlocal gauge transformations to guarantee the over-all invariance of the regularized amplitudes. Cornish has derived the primary graviton vertices and the BRST symmetry relations for the regularized $`\widehat{W}_{\mathrm{grav}}`$ , using the nonlinear, nonlocal extended gauge transformations suitable for the perturbative gravity equations. The regularized graviton propagator in the fixed de Donder gauge $`\alpha =1`$ is given by $$D_{\mu \nu \rho \sigma }^{\mathrm{grav}}(x)=(\eta _{\mu \rho }\eta _{\nu \sigma }+\eta _{\mu \sigma }\eta _{\nu \rho }\eta _{\mu \nu }\eta _{\rho \sigma })$$ $$\times \left(\frac{i}{(2\pi )^4}\right)d^4k\frac{^2(k^2)}{k^2iϵ}\mathrm{exp}[ik(xx^{})],$$ (87) while the shadow propagator is $$D_{\mu \nu \rho \sigma }^{\mathrm{shad}}(x)=(\eta _{\mu \rho }\eta _{\nu \sigma }+\eta _{\mu \sigma }\eta _{\nu \rho }\eta _{\mu \nu }\eta _{\rho \sigma })$$ $$\times \left(\frac{i}{(2\pi )^4}\right)d^4k\frac{[1^2(k^2)]}{k^2iϵ}\mathrm{exp}[ik(xx^{})].$$ (88) The ghost propagator in momentum space is given by $$D_{\mu \nu }^G(p)=\frac{\eta _{\mu \nu }^2(p^2)}{p^2},$$ (89) while the shadow ghost propagator is $$D_{\mu \nu }^{\mathrm{shad}\mathrm{G}}(p)=\frac{\eta _{\mu \nu }[1^2(p^2)]}{p^2}.$$ (90) In momentum space we have $$\frac{i^2(k^2)}{k^2iϵ}=i_1^{\mathrm{}}\frac{d\tau }{\mathrm{\Lambda }_G^2}\mathrm{exp}\left(\tau \frac{k^2}{\mathrm{\Lambda }_G^2}\right),$$ (91) and $$\frac{i(^2(k^2)1)}{k^2iϵ}=i_0^1\frac{d\tau }{\mathrm{\Lambda }_G^2}\mathrm{exp}\left(\tau \frac{k^2}{\mathrm{\Lambda }_G^2}\right),$$ (92) where $`\mathrm{\Lambda }_G`$ is the gravitational scale parameter. The local propagator is reproduced by subtracting $`D^{\mathrm{shad}}`$ from $`D^{\mathrm{grav}}`$, while the “stripped” vertices are obtained by subtracting the amplitudes containing the shadow propagator $`D^{\mathrm{shad}}`$ from the amplitudes containing the regulator operators (86). We can facilitate the calculations by separating the free and interacting parts of the action $$W_{\mathrm{grav}}(\gamma )=W_{\mathrm{grav}}^F(\gamma )+W_{\mathrm{grav}}^I(\gamma ).$$ (93) The finite regularized gravitational action is given by $$\widehat{W}_{\mathrm{grav}}(\gamma ,s)=W_{\mathrm{grav}}^F(\widehat{\gamma })P_{\mathrm{grav}}(s)+W_{\mathrm{grav}}^I(\gamma +s),$$ (94) where $$\widehat{\gamma }=^1\gamma ,P_{\mathrm{grav}}(s)=d^4x𝒢(\sqrt{s},s_i𝒪_{ij}^1s_j),$$ (95) $`s`$ denotes the graviton shadow field, and $`𝒢`$ denotes the detailed expansion of the contributions formed from the shadow field. The regularized Lagrangian density up to order $`\kappa ^2`$ is invariant under the extended BRST transformations : $$\widehat{\delta }_0\gamma _{\mu \nu }=X_{\mu \nu \lambda }^{(0)}C^\lambda \sigma =(_\nu C_\mu +_\mu C_\nu \eta _{\mu \nu }_\lambda C^\lambda )\sigma ,$$ (96) $$\widehat{\delta }_1\gamma _{\mu \nu }=\kappa ^2X_{\mu \nu \lambda }^{(1)}C^\lambda \sigma =\kappa ^2(2\gamma _{\rho (\mu }^\rho C_{\nu )}_\lambda \gamma _{\mu \nu }C^\lambda \gamma _{\mu \nu }_\lambda C^\lambda ),$$ (97) $$\widehat{\delta }_0\overline{C}^\nu =2_\mu \gamma ^{\mu \nu }\sigma ,$$ (98) $$\widehat{\delta }_1C_\nu =\kappa ^2C^\mu _\mu C_\nu \sigma .$$ (99) The order $`\kappa ^2`$ transformations are $$\widehat{\delta }_2\gamma _{\mu \nu }=\kappa ^2^2[2^\rho C_{(\nu }D_{\mu )\rho \kappa \lambda }^{\mathrm{shad}}(B^{\kappa \lambda }+H^{\kappa \lambda })$$ $$C^\rho D_{\mu \nu \kappa \lambda }^{\mathrm{shad}}(_\rho B^{\kappa \lambda }+_\rho H^{\kappa \lambda })_\rho C^\rho D_{\mu \nu \kappa \lambda }^{\mathrm{shad}}(B^{\kappa \lambda }+H^{\kappa \lambda })$$ $$+2\gamma _{\rho (\mu }D_{\nu )\kappa }^{\mathrm{shad}\mathrm{G}}^\rho H^\kappa _\rho \gamma _{\mu \nu }D^{\mathrm{shad}\mathrm{G}\rho \kappa }H_\kappa \gamma _{\mu \nu }D^{\mathrm{shad}\mathrm{G}\rho \kappa }_\rho H_\kappa ]\sigma ,$$ (100) $$\widehat{\delta }_2C_\nu =\kappa ^2^2(_\mu C_\nu D^{\mathrm{shad}\mathrm{G}\rho \kappa }H_\kappa +C_\mu D_{\nu \kappa }^{\mathrm{shad}\mathrm{G}}^\mu H^\kappa )\sigma .$$ (101) Here, we have $$H^{\alpha \beta }=(^{(\alpha }\overline{C}_\rho ^{\beta )}C^\rho +^\rho \overline{C}^{(\alpha }^{\beta )}C_\rho +^\rho ^{(\beta }\overline{C}^{\alpha )}C_\rho ),$$ (102) $$H^\rho =\gamma _{\lambda \kappa }^\lambda ^\kappa C^\rho +^\kappa \gamma _{\lambda \kappa }^\lambda C^\rho _\kappa _\lambda \gamma ^{\rho \kappa }C^\lambda _\kappa \gamma ^{\rho \kappa }_\lambda C^\lambda ,$$ (103) $$\overline{H}^\rho =^\lambda \overline{C}^\rho ^\kappa \gamma _{\lambda \kappa }+^\lambda ^\kappa \overline{C}^\rho \gamma _{\lambda \kappa }+^\rho \overline{C}^\lambda ^\kappa \gamma _{\lambda \kappa }.$$ (104) Because we have extended the gauge symmetry to nonlinear, nonlocal transformations, we must also supplement the quantization procedure with an invariant measure $$=\mathrm{\Delta }(𝐠,\overline{C},C)D[𝐠_{\mu \nu }]D[\overline{C}_\lambda ]D[C_\sigma ]$$ (105) such that $`\delta =0`$. As we have demonstrated, the quantum gravity perturbation theory is invariant under the FQFT generalized, nonlinear field representation dependent transformations. It is unitary and finite to all orders in a way similar to the non-Abelian gauge theories formulated using FQFT. At the tree graph level all unphysical polarization states are decoupled and nonlocal effects will only occur in graviton and graviton-matter loop graphs. Because the gravitational tree graphs are purely local there is a well-defined classical GR limit. The finite quantum gravity theory is well-defined in four real spacetime dimensions. We quantize by means of the path integral operation $$0\|T^{}(O[𝐠])\|0_{}=[D𝐠]\mu [𝐠](\mathrm{gauge}\mathrm{fixing})O[\widehat{𝐠}]\mathrm{exp}(i\widehat{W}_{\mathrm{grav}}[𝐠]).$$ (106) The quantization is carried out in the functional formalism by finding a measure factor $`\mu [𝐠]`$ to make $`[D𝐠]`$ invariant under the classical symmetry. To ensure a correct gauge fixing scheme, we write $`W_{\mathrm{grav}}[𝐠]`$ in the BRST invariant form with ghost fields; the ghost structure arises from exponentiating the Faddeev-Popov determinant . The algebra of extended gauge symmetries is not expected to close off-shell, so one needs to introduce higher ghost terms (beyond the normal ones) into both the action and the BRST transformation. The BRST action will be regularized directly to ensure that all the corrections to the measure factor are included. ## 8 A Resolution of The Higgs Hierarchy Problem It is time to discuss the Higgs sector hierarchy problem . The gauge hierarchy problem is related to the spin $`0^+`$ scalar field nature of the Higgs particle in the standard model with quadratic mass divergence and no protective extra symmetry at $`m=0`$. In standard point particle, local field theory the fermion masses are logarithmically divergent and there exists a chiral symmetry restoration at $`m=0`$. Writing $`m_H^2=m_{0H}^2+\delta m_H^2`$, where $`m_{0H}`$ is the bare Higgs mass and $`\delta m_H`$ is the Higgs self-energy renormalization constant, we get for the one loop Feynman graph in $`D=4`$ spacetime: $$\delta m_H^2\frac{g}{32\pi ^2}M_c^2,$$ (107) where $`M_c`$ is a cutoff parameter. If we want to understand the nature of the Higgs mass we must require that $$\delta m_H^2O(m_H^2),$$ (108) i.e. the quadratic divergence should be cut off at the mass scale of the order of the physical Higgs mass. Since $`m_H\sqrt{2g}v`$, where $`v=<\varphi >_0`$ is the vacuum expectation value of the scalar field $`\varphi `$ and $`v=246`$ GeV from the electroweak theory, then in order to keep perturbation theory valid, we must demand that $`10\mathrm{GeV}m_H350\mathrm{GeV}`$ and we need $$M_c=M_{\mathrm{Higgs}}1\mathrm{TeV},$$ (109) where the lower bound on $`m_H`$ comes from the avoidance of washing out the spontaneous symmetry breaking of the vacuum. Nothing in the standard model can tell us why (109) should be true, so we must go beyond the local standard model to solve the problem. $`M_c`$ is an arbitrary parameter in point particle field theory with no physical interpretation. Since all particles interact through gravity, then ultimately we should expect to include gravity in the standard model, so we expect that $`M_{\mathrm{Planck}}10^{19}`$ GeV should be the natural cutoff. Then we have using (109) and $`g1`$: $$\frac{\delta m_H^2(M_{\mathrm{Higgs}})}{\delta m_H^2(M_{\mathrm{Planck}})}\frac{M_{\mathrm{Higgs}}^2}{M_{\mathrm{Planck}}^2}10^{34},$$ which represents an intolerable fine-tuning of parameters. This ‘naturalness’ or hierarchy problem is one of the most serious defects of the standard model. There have been two strategies proposed as ways out of the hierarchy problem. The Higgs is taken to be composite at a scale $`M_c1`$ TeV, thereby providing a natural cutoff in the quadratically divergent Higgs loops. One such scenario is the ‘technicolor’ model, but it cannot be reconciled with the accurate standard model data, nor with the smallness of fermion masses and the flavor-changing neutral current interactions. The other strategy is to postulate supersymmetry, so that the opposite signs of the boson and fermion lines cancel by means of the non-renormalization theorem. However, supersymmetry is badly broken at lower energies, so we require that $$\delta m_H^2\frac{g}{32\pi ^2}\|M_{c\mathrm{bosons}}^2M_{c\mathrm{fermions}}^2\|1\mathrm{TeV}^2,$$ or, in effect $$\|m_bm_f\|1\mathrm{TeV}.$$ This physical requirement leads to the prediction that the supersymmetric partners of known particles should have a threshold $`1`$ TeV. A third possible strategy is to introduce the FQFT formalism, and realize a field theory mechanism which will introduce a natural physical scale in the theory $`\mathrm{\Lambda }_H1`$ TeV, which will protect the Higgs mass from becoming large and unstable. Let us consider the regularized scalar field FQFT Lagrangian in Minkowski spacetime $$\widehat{}_S=\frac{1}{2}\widehat{\varphi }(^2m^2)\widehat{\varphi }\frac{1}{2}\rho 𝒪^1\rho +\frac{1}{2}Z^1\delta m^2(\varphi +\rho )^2\frac{1}{24}g_0(\varphi +\rho )^4,$$ (110) where $`\varphi =Z^{1/2}\varphi _R`$ is the bare field, $`\varphi _R`$ is the renormalized field, $`\widehat{\varphi }=^1\varphi `$, $`\rho `$ is the shadow field, $`m_0`$ is the bare mass, $`Z`$ is the field strength renormalization constant, $`\delta m^2`$ is the mass renormalization constant and $`m`$ is the physical mass. The regularizing operator is given by $$_m=\mathrm{exp}\left(\frac{^2m^2}{2\mathrm{\Lambda }_H^2}\right),$$ (111) while the shadow kinetic operator is $$𝒪^1=\frac{^2m^2}{_m^21}.$$ (112) Here, $`\mathrm{\Lambda }_H`$ is the Higgs scalar field energy scale in FQFT, which determines the scale of nonlocalizability of the Higgs particle. The full propagator is $$i\mathrm{\Delta }_R(p^2)=\frac{i_m^2}{p^2+m^2iϵ}=i_1^{\mathrm{}}\frac{d\tau }{\mathrm{\Lambda }_H^2}\mathrm{exp}\left[\tau \left(\frac{p^2+m^2}{\mathrm{\Lambda }_H^2}\right)\right],$$ (113) whereas the shadow propagator is $$i\mathrm{\Delta }_{\mathrm{shadow}}=i\frac{_m^21}{p^2+m^2}=i_0^1\frac{d\tau }{\mathrm{\Lambda }_H^2}\mathrm{exp}\left[\tau \left(\frac{p^2+m^2}{\mathrm{\Lambda }_H^2}\right)\right].$$ (114) Let us define the self-energy $`\mathrm{\Sigma }(p)`$ as a Taylor series expansion around the mass shell $`p^2=m^2`$: $$\mathrm{\Sigma }(p^2)=\mathrm{\Sigma }(m^2)+(p^2+m^2)\frac{\mathrm{\Sigma }}{p^2}(m^2)+\stackrel{~}{\mathrm{\Sigma }}(p^2),$$ (115) where $`\stackrel{~}{\mathrm{\Sigma }}(p^2)`$ is the usual finite part in the point particle limit $`\mathrm{\Lambda }_H\mathrm{}`$. We have $$\stackrel{~}{\mathrm{\Sigma }}(m^2)=0,$$ (116) and $$\frac{\stackrel{~}{\mathrm{\Sigma }}(p^2)}{p^2}(p^2=m^2)=0.$$ (117) The full propagator is related to the self-energy $`\mathrm{\Sigma }(p^2)`$ by $$i\mathrm{\Delta }_R(p^2)=\frac{i_m^2[1+𝒪\mathrm{\Sigma }(p^2)]}{p^2+m^2+\mathrm{\Sigma }(p^2)}=\frac{iZ}{p^2+m^2+\mathrm{\Sigma }_R(p^2)}.$$ (118) Here $`\mathrm{\Sigma }_R(p^2)`$ is the renormalized self-energy which can be written as $$\mathrm{\Sigma }_R(p^2)=(p^2+m^2)\left[\frac{Z}{_m^2(1+𝒪\mathrm{\Sigma })}1\right]+\frac{Z\mathrm{\Sigma }}{_m^2(1+𝒪\mathrm{\Sigma })}.$$ (119) The 1PI two-point function is given by $$i\mathrm{\Gamma }_R^{(2)}(p^2)=i[\mathrm{\Delta }_R(p^2)]^1=\frac{i[p^2+m^2+\mathrm{\Sigma }(p^2)]}{_m^2[1+𝒪\mathrm{\Sigma }(p^2)]}.$$ (120) Since $`_m1`$ and $`𝒪0`$ as $`\mathrm{\Lambda }_H\mathrm{}`$, then in this limit $$i\mathrm{\Gamma }_R^{(2)}(p^2)=i[p^2+m^2+\mathrm{\Sigma }(p^2)],$$ (121) which is the standard point particle result. The mass renormalization is determined by the propagator pole at $`p^2=m^2`$ and we have $$\mathrm{\Sigma }_R(m^2)=0.$$ (122) Also, we have the condition $$\frac{\mathrm{\Sigma }_R(p^2)}{p^2}(p^2=m^2)=0.$$ (123) The renormalized coupling constant is defined by the four-point function $`\mathrm{\Gamma }_R^{(4)}(p_1,p_2,p_3,p_4)`$ at the point $`p_i=0`$: $$\mathrm{\Gamma }_R^{(4)}(0,0,0,0)=g.$$ (124) The bare coupling constant $`g_0`$ is determined by $$Z^2g_0=g+\delta g(g,m^2,\mathrm{\Lambda }_H^2).$$ (125) Moreover, $$Z=1+\delta Z(g,m^2,\mathrm{\Lambda }_H^2),$$ $$Zm_0^2=Zm^2\delta m^2(g,m^2,\mathrm{\Lambda }_H^2).$$ A calculation of the scalar field mass renormalization in D-dimensional space gives : $$\delta m^2=\frac{g}{2^{D+1}\pi ^{D/2}}m^{D2}\mathrm{\Gamma }(1\frac{D}{2},\frac{m^2}{\mathrm{\Lambda }_H^2})+O(g^2),$$ (126) where $`\mathrm{\Gamma }(n,z)`$ is the incomplete gamma function: $$\mathrm{\Gamma }(n,z)=_z^{\mathrm{}}\frac{dt}{t}t^n\mathrm{exp}(t)=(n1)\mathrm{\Gamma }(n1,z)+z^{n1}\mathrm{exp}(z).$$ (127) We have $$\mathrm{\Gamma }(1,z)=E_i(z)+\frac{1}{z}\mathrm{exp}(z),$$ (128) where $`E_i(z)`$ is the exponential integral $$E_i(z)_z^{\mathrm{}}𝑑t\frac{\mathrm{exp}(t)}{t}.$$ For small $`z`$ we obtain the expansion $$E_i(z)=\mathrm{ln}(z)\gamma +z\frac{z^2}{22!}+\frac{z^3}{33!}\mathrm{},$$ (129) where $`\gamma `$ is Euler’s constant. For large positive values of $`z`$, we have the asymptotic expansion $$E_i(z)\mathrm{exp}(z)\left[\frac{1}{z}\frac{1}{z^2}+\frac{2!}{z^3}\mathrm{}\right].$$ (130) Thus, for small $`m/\mathrm{\Lambda }_H`$ we obtain in $`D=4`$ spacetime: $$\delta m^2=\frac{g}{32\pi ^2}\left[\mathrm{\Lambda }_H^2m^2\mathrm{ln}\left(\frac{\mathrm{\Lambda }_H^2}{m^2}\right)m^2(1\gamma )+O\left(\frac{m^2}{\mathrm{\Lambda }_H^2}\right)\right]+O(g^2),$$ (131) which is the standard quadratically divergent self-energy, obtained from a cutoff procedure or a dimensional regularization scheme. We have for $`z\mathrm{}`$: $$\mathrm{\Gamma }(a,z)z^{a1}\mathrm{exp}(z)\left[1+\frac{a1}{z}+O\left(\frac{1}{z^2}\right)\right]$$ (132) so that for $`m\mathrm{\Lambda }_H`$, we get in four-dimensional spacetime $$\delta m^2\frac{g}{32\pi ^2}\left(\frac{\mathrm{\Lambda }_H^4}{m^2}\right)\mathrm{exp}\left(\frac{m^2}{\mathrm{\Lambda }_H^2}\right).$$ (133) Thus, the Higgs self-energy one loop graph falls off exponentially fast for $`m\mathrm{\Lambda }_H`$. We have succeeded in stabilizing the radiative corrections to the Higgs sector, solving the Higgs hierarchy problem for $`\mathrm{\Lambda }_H1`$ TeV. ## 9 Gluon and Gravitational Vacuum Polarization A calculation of the one-loop gluon vacuum polarization in FQFT gives the tensor in D-dimensions $$\mathrm{\Pi }_{ik}^{\mu \nu }(p)=\frac{g^2}{2^Dp^{D/2}}f_{ilm}f_{klm}(p^2\eta ^{\mu \nu }p^\mu p^\nu )\mathrm{\Pi }(p^2),$$ (134) where $`p`$ is the gluon momentum and $$\mathrm{\Pi }(p^2)=2\mathrm{exp}\left(p^2/\mathrm{\Lambda }_{\mathrm{YM}}^2\right)_0^{1/2}𝑑y\mathrm{\Gamma }(2D/2,yp^2/\mathrm{\Lambda }_{\mathrm{YM}}^2)[y(1y)p^2]^{D/22}$$ $$\times [2(D2)y(1y)\frac{1}{2}(D6)].$$ (135) We observe that $`\mathrm{\Pi }_{ik\mu }^\mu (0)=0`$ a result that is required by gauge invariance and the fact that the gluon has zero mass. The dimensionally regulated gluon vacuum polarization result is obtained by the replacement $$\mathrm{\Gamma }(2D/2,yp^2/\mathrm{\Lambda }_{\mathrm{YM}}^2)\mathrm{\Gamma }(2D/2)$$ (136) and choosing $`p^2\mathrm{\Lambda }_{\mathrm{YM}}^2`$. In four-dimensions we get $$\mathrm{\Pi }(p^2)=2\mathrm{exp}(p^2/\mathrm{\Lambda }_{\mathrm{YM}}^2)_0^{1/2}𝑑yE_i(yp^2/\mathrm{\Lambda }_{\mathrm{YM}}^2)[4y(1y)+1],$$ (137) where we have used the relation $$\mathrm{\Gamma }(0,z)E_i(z)=_z^{\mathrm{}}𝑑t\mathrm{exp}(t)t^1.$$ (138) By using the behavior for large Euclidean momentum $`p^2\mathrm{\Lambda }_{\mathrm{YM}}^2`$: $$E_i(yp^2/\mathrm{\Lambda }_{\mathrm{YM}}^2)\frac{\mathrm{\Lambda }_{\mathrm{YM}}^2}{p^2}\mathrm{exp}(yp^2/\mathrm{\Lambda }_{\mathrm{YM}}^2),$$ (139) we find from (137) that $$\mathrm{\Pi }(p^2)\frac{\mathrm{\Lambda }_{\mathrm{YM}}^2}{p^2}\mathrm{exp}(p^2/\mathrm{\Lambda }_{\mathrm{YM}}^2)[\mathrm{exp}(p^2/2\mathrm{\Lambda }_{\mathrm{YM}}^2)+\frac{4\mathrm{\Lambda }_{\mathrm{YM}}^2}{p^2}$$ $$+\frac{\mathrm{\Lambda }_{\mathrm{YM}}^4}{p^4}\frac{16\mathrm{\Lambda }_{\mathrm{YM}}^6}{p^6}].$$ (140) Thus, the gluon vacuum polarization is exponentially damped for $`p^2\mathrm{\Lambda }_{\mathrm{YM}}^2`$. The lowest order contributions to the graviton self-energy in FQFT will include the standard graviton loops, the shadow field graviton loops, the ghost field loop contributions with their shadow field counterparts, and the measure loop contributions. In the regularized perturbative gravity theory the first order vacuum polarization tensor $`\mathrm{\Pi }^{\mu \nu \rho \sigma }`$ must satisfy the Slavnov-Ward identities : $$p_\mu p_\rho D^{\mu \nu \alpha \beta }(p)\mathrm{\Pi }_{\alpha \beta \gamma \delta }(p)D^{\gamma \delta \rho \sigma }(p)=0.$$ (141) By symmetry and Lorentz invariance, the vacuum polarization tensor must have the form $$\mathrm{\Pi }_{\alpha \beta \gamma \delta }(p)=\mathrm{\Pi }_1(p^2)p^4\eta _{\alpha \beta }\eta _{\gamma \delta }+\mathrm{\Pi }_2(p^2)p^4(\eta _{\alpha \gamma }\eta _{\beta \delta }+\eta _{\alpha \delta }\eta _{\beta \gamma })$$ $$+\mathrm{\Pi }_3(p^2)p^2(\eta _{\alpha \beta }p_\gamma p_\delta +\eta _{\gamma \delta }p_\alpha p_\beta )+\mathrm{\Pi }_4(p^2)p^2(\eta _{\alpha \gamma }p_\beta p_\delta +\eta _{\alpha \delta }p_\beta p_\gamma $$ $$+\eta _{\beta \gamma }p_\alpha p_\delta +\eta _{\beta \delta }p_\alpha p_\gamma )+\mathrm{\Pi }_5(p^2)p_\alpha p_\beta p_\gamma p_\delta .$$ (142) The Slavnov-Ward identities impose the restrictions $$\mathrm{\Pi }_2+\mathrm{\Pi }_4=0,4(\mathrm{\Pi }_1+\mathrm{\Pi }_2\mathrm{\Pi }_3)+\mathrm{\Pi }_5=0.$$ (143) The basic lowest order graviton self-energy diagram is determined by : $$\mathrm{\Pi }_{\mu \nu \rho \sigma }^1(p)=\frac{1}{2}\kappa ^2\mathrm{exp}\left(p^2/\mathrm{\Lambda }_G^2\right)d^4q𝒰_{\mu \nu \alpha \beta \gamma \delta }(p,q,qp)D^{\alpha \beta \kappa \lambda }(q)$$ $$\times D^{\gamma \delta \tau \xi }(pq)𝒰_{\kappa \lambda \tau \xi \rho \sigma }(q,pq,p),$$ (144) where $`𝒰`$ is the three-graviton vertex function $$𝒰_{\mu \nu \rho \sigma \delta \tau }(q_1,q_2,q_3)=\frac{1}{2}[q_{2(\mu }q_{3\nu )}(2\eta _{\rho (\delta }\eta _{\tau )\sigma }\frac{2}{D2}\eta _{\mu \nu }\eta _{\delta \tau })$$ $$+q_{1(\rho }q_{3\sigma )}(2\eta _{\mu (\delta }\eta _{\tau )\nu }\frac{2}{D2}\eta _{\mu \nu }\eta _{\delta \tau })+\mathrm{}],$$ (145) and the ellipsis denote similar contributions. To this diagram, we must add the ghost particle diagram contribution $`\mathrm{\Pi }^2`$, the shadow diagram contribution $`\mathrm{\Pi }^3`$ and the measure diagram contribution $`\mathrm{\Pi }^4`$. The dominant finite contribution to the graviton self-energy will be of the form $$\mathrm{\Pi }_{\mu \nu \rho \sigma }(p)\kappa ^2\mathrm{\Lambda }_G^4\mathrm{exp}\left(p^2/\mathrm{\Lambda }_G^2\right)Q_{\mu \nu \rho \sigma }(p^2)$$ $$\frac{\mathrm{\Lambda }_G^4}{M_{\mathrm{PL}}^2}\mathrm{exp}\left(p^2/\mathrm{\Lambda }_G^2\right)Q_{\mu \nu \rho \sigma }(p^2),$$ (146) where $`M_{\mathrm{PL}}`$ is the reduced Planck mass and $`Q(p^2)`$ is a finite remaining part. For renormalizable field theories such as quantum electrodynamics and Yang-Mills theory, we will find that in FQFT the loop contributions are controlled by the incomplete $`\mathrm{\Gamma }`$-function. If we adopt an “effective” quantum gravity theory expansion in the energy , then we would expect to obtain $$\mathrm{\Pi }_{\mu \nu \rho \sigma }(p)\kappa ^2\mathrm{exp}(p^2/\mathrm{\Lambda }_G^2)(\mathrm{\Gamma }(2D/2,p^2/\mathrm{\Lambda }_G^2)Q_{\mu \nu \rho \sigma }(p^2),$$ (147) where $``$ denotes the functional dependence on the incomplete $`\mathrm{\Gamma }`$-function. By making the replacement $$(\mathrm{\Gamma }(2D/2),p^2/\mathrm{\Lambda }_G^2)(\mathrm{\Gamma }(2D/2)),$$ (148) we would then obtain the second order graviton loop calculations using dimensional regularization . The dominant behavior will now be $`\mathrm{ln}(\mathrm{\Lambda }_G^2/q^2)`$ and not $`\mathrm{\Lambda }_G^4`$. However, in a nonrenormalizable theory such as quantum gravity, the dimensional regularization technique may not provide a correct result for the dominant behavior of the loop integral and we expect the result to be of order $`\mathrm{\Lambda }_G^4`$. Indeed, it is well known that dimensional regularization for massless particles removes all contributions from tadpole graphs and $`\delta ^4(0)`$ contact terms. On the other hand, FQFT takes into account all leading order contributions and provides a complete account of all counterterms. Because all the scattering amplitudes are finite, then renormalizability is no longer an issue. The function $$Q_{\mu }^{}{}_{}{}^{\mu \sigma }{}_{\sigma }{}^{}(p^2)p^4$$ (149) as $`p^20`$. Therefore, $`\mathrm{\Pi }_{\mu \nu \rho \sigma }(p)`$ vanishes at $`p^2=0`$ as it should from gauge invariance and for massless gravitons. In Euclidean momentum space, which we can reach by a Wick rotation, we see that for $`p^2\mathrm{\Lambda }_G^2`$ the graviton self-energy (9) is exponentially damped and the quantum gravity loop corrections are negligible for energies greater than $`\mathrm{\Lambda }_G`$. It is often argued in the literature on quantum gravity that the gravitational quantum corrections scale as $`\alpha _G=GE^2`$, so that for sufficiently large values of the energy $`E`$, namely, of order the Planck energy, the gravitational quantum fluctuations become large. We see that in FQFT this will not be the case, because the finite quantum loop corrections become negligible in the high energy limit provided the perturbative approximation is valid. Of course, the contributions of the tree graph exchanges of virtual gravitons can be large in the high energy limit, corresponding to strong classical gravitational fields. It follows that for high enough energies, a classical curved spacetime would be a good approximation, at least until the perturbation calculations break down. In contrast to recent models of branes and strings in which the higher-dimensional compactification scale is lowered to the TeV range , we retain the classical GR gravitation picture and its Newtonian limit. It is perhaps a radical notion to entertain that quantum gravity becomes weaker as the energy scale increases towards the Planck scale $`10^{19}`$ Gev, but there is, of course, no known experimental reason why this should not be the case in nature. However, we do not expect that our weak gravity field expansion is valid at the Planck scale when $`GE^21`$, although the exponential damping of the quantum gravity loop graphs could still persist at the Planck scale. This question remains unresolved until a nonperturbative solution to quantum gravity is found. It is worth noting that in the framework of an effective gravitational field theory , the leading lowest order loop divergence can be “renormalized” by being absorbed into two parameters $`c_1`$ and $`c_2`$. For a non-flat spacetime background metric $`\overline{g}_{\mu \nu }`$, the divergent term at one loop due to graviton and ghost loops is given by : $$_{1\mathrm{l}\mathrm{o}\mathrm{o}\mathrm{p}}^{\mathrm{div}}=\frac{1}{8\pi ^2ϵ}[\frac{1}{120}\overline{R}^2+\frac{7}{20}\overline{R}_{\mu \nu },\overline{R}^{\mu \nu }],$$ (150) where $`ϵ=4D`$ and the effective field theory renormalization parameters are $$c_1^{(r)}=c_1+\frac{1}{960\pi ^2ϵ},c_2^{(r)}=c_2+\frac{7}{160\pi ^2ϵ}.$$ (151) ## 10 A Quantum Gravity Resolution of the Cosmological Constant Problem Zeldovich showed that the zero-point vacuum fluctuations must have a Lorentz invariant form $$T_{\mathrm{vac}\mu \nu }=\lambda _{\mathrm{vac}}g_{\mu \nu },$$ (152) consistent with the equation of state $`\rho _{\mathrm{vac}}=p_{\mathrm{vac}}`$. Thus, the vacuum within the framework of particle quantum physics has properties identical to the cosmological constant. In quantum theory, the second quantization of a classical field of mass $`m`$, treated as an ensemble of oscillators each with a frequency $`\omega (k)`$, leads to a zero-point energy $`E_0=_k\frac{1}{2}\mathrm{}\omega (k)`$. The experimental confirmation of a zero-point vacuum fluctuation was demonstrated by the Casimir effect . A simple evaluation of the vacuum density obtained from a summation of the zero-point energy modes gives $$\rho _{\mathrm{vac}}=\frac{1}{(2\pi )^2}_0^{M_c}𝑑kk^2(k^2+m^2)^{1/2}\frac{M_c^4}{16\pi ^2},$$ (153) where $`M_c`$ is the cutoff. Already at the level of the standard model, we get $`\rho _{\mathrm{vac}}(10^2\mathrm{GeV})^4`$ which is $`55`$ orders of magnitude larger than the bound (2). To agree with the experimental bound (2), we would have to invoke a very finely tuned cancellation of $`\lambda _{\mathrm{vac}}`$ with the “bare ” cosmological constant $`\lambda `$, which is generally conceded to be theoretically unacceptable. We can understand this result by using the language of Feynman graphs. To avoid undue technical issues in FQFT, we shall consider initially the basic lowest order vacuum fluctuation diagram computed from the matrix element in flat Minkowski spacetime $$M_{(2)}^{(0)}g^2d^4pd^4p^{}d^4k\delta (k+pp^{})\delta (k+pp^{})$$ $$\times \frac{1}{k^2+m^2}\mathrm{Tr}\left(\frac{i\gamma ^\sigma p_\sigma m_f}{p^2+m_f^2}\gamma ^\mu \frac{i\gamma ^\sigma p_\sigma ^{}m_f}{p^{}_{}{}^{}2+m_f^2}\gamma _\mu \right)$$ $$\times \mathrm{exp}[\left(\frac{p^2+m_f^2}{\mathrm{\Lambda }_{\mathrm{SM}}^2}\right)\left(\frac{p^2+m_f^2}{\mathrm{\Lambda }_{\mathrm{SM}}^2}\right)\frac{k^2}{\mathrm{\Lambda }_{\mathrm{SM}}^2})],$$ (154) where $`g`$ is a coupling constant associated with the standard model. We have considered a closed loop made of a standard model fermion of mass $`m_f`$, an antifermion of the same mass and an internal standard model boson propagator of mass $`m`$; the scale $`\mathrm{\Lambda }_{\mathrm{SM}}10^210^3`$ GeV. This leads to the result $$M_{(2)}^{(0)}16\pi ^4g^2\delta ^4(a)_0^{\mathrm{}}𝑑pp^3𝑑p^{}p^{}_{}{}^{}3\left[\frac{P^2+p^2+p^{}_{}{}^{}2+4m_f^2}{(P+a)(Pa)}\right]$$ $$\times \frac{1}{(p^2+m_f^2)(p^^2+m_f^2)}\mathrm{exp}\left[\frac{(p^2+p^{}_{}{}^{}2+2m_f^2)}{\mathrm{\Lambda }_{\mathrm{SM}}^2}\frac{P^2}{\mathrm{\Lambda }_{\mathrm{SM}}^2}\right],$$ (155) where $`P=pp^{}`$ and $`a`$ is an infinitesimal constant which formally regularizes the infinite volume factor $`\delta ^4(0)`$. We see that $`\rho _{\mathrm{vac}}M_{(2)}^{(0)}`$ is finite and $`M_{(2)}^{(0)}\mathrm{\Lambda }_{\mathrm{SM}}^4`$. To maintain gauge invariance and unitarity in FQFT, we must add to this result the contributions from the ghost diagram, the shadow diagram and the measure diagram. In flat Minkowski spacetime, the sum of all disconnected vacuum diagrams $`C=_nM_n^{(0)}`$ is a constant factor in the scattering S-matrix $`S^{}=SC`$. Since the S-matrix is unitary $`\|S^{}\|^2=1`$, then we must conclude that $`\|C\|^2=1`$, and all the disconnected vacuum graphs can be ignored. However, due to the equivalence principle gravity couples to all forms of energy, including the vacuum energy density $`\rho _{\mathrm{vac}}`$, so we can no longer ignore these virtual quantum fluctuations in the presence of a non-zero gravitational field. Let us now consider the dominant contributions to the vacuum density arising from the graviton loop corrections. As explained above, we shall perform the calculations by expanding about flat spacetime and trust that the results still hold for an expansion about a curved metric background field, which is strictly required for a non-zero cosmological constant. Since the scales involved in the final answer, including the predicted smallness of the cosmological constant, correspond to a very small curvature of spacetime, we expect that our approximation is justified. We shall adopt a simple model consisting of a massive vector meson $`V_\mu `$, which has the standard model energy scale $`10^210^3`$ GeV. We have for the vector field Lagrangian density $$_V=\frac{1}{4}(𝐠)^{1/2}𝐠^{\mu \nu }𝐠^{\alpha \beta }F_{\mu \alpha }F_{\nu \beta }+m_V^2V_\mu V^\mu ,$$ (156) where $$F_{\mu \nu }=_\nu V_\mu _\mu V_\nu .$$ (157) We include in the Lagrangian density an additional piece $`\frac{1}{2}(_\mu V^\mu )^2`$, and the vector field propagator has the form $$D_{\mu \nu }^\mathrm{V}=\frac{\eta _{\mu \nu }}{p^2+m_V^2iϵ}\mathrm{exp}\left[(p^2+m_V^2)/\mathrm{\Lambda }_{\mathrm{SM}}^2\right]$$ $$=\eta _{\mu \nu }_1^{\mathrm{}}\frac{d\tau }{\mathrm{\Lambda }_{\mathrm{SM}}^2}\mathrm{exp}\left[\tau (p^2+m_V^2)/\mathrm{\Lambda }_{\mathrm{SM}}^2\right],$$ (158) while the shadow propagator is $$D_{\mu \nu }^{\mathrm{shad}\mathrm{V}}=\frac{\eta _{\mu \nu }}{p^2+m_V^2}\left[1\mathrm{exp}\left[(p^2+m_V^2)/\mathrm{\Lambda }_{\mathrm{SM}}^2\right]\right]$$ $$=\eta _{\mu \nu }_0^1\frac{d\tau }{\mathrm{\Lambda }_{\mathrm{SM}}^2}\mathrm{exp}\left[\tau (p^2+m_V^2)/\mathrm{\Lambda }_{\mathrm{SM}}^2\right].$$ (159) The graviton-V-V vertex in momentum space is given by $$𝒱_{\alpha \beta \lambda \sigma }(p,q_1,q_2)=\eta _{\lambda \sigma }q_{1(\alpha }q_{2\beta )}\eta _{\sigma (\beta }q_{1\alpha )}q_\lambda \eta _{\lambda (\alpha }q_{1_\sigma }q_{2\beta )}$$ $$+\eta _{\sigma (\beta }\eta _{\alpha )\lambda }q_1q_2\frac{1}{D2}\eta _{\alpha \beta }(\eta _{\lambda \sigma }q_1q_2q_{1\sigma }q_{2\lambda }),$$ (160) where $`q_1,q_2`$ denote the momenta of the two $`Vs`$ connected to the graviton with momentum $`p`$. We use the notation $`A_{(\alpha }B_{\beta )}=\frac{1}{2}(A_\alpha B_\beta +A_\beta B_\alpha )`$. The lowest order correction to the graviton vacuum loop will have the form $$\mathrm{\Pi }_{\mu \nu \rho \sigma }^V(p)=\kappa ^2\mathrm{exp}\left(p^2/\mathrm{\Lambda }_G^2\right)d^4q𝒱_{\mu \nu \lambda \alpha }(p,q,qp)D^{V,W\lambda \delta }(q)$$ $$\times 𝒱_{\rho \sigma \kappa \delta }(p,pq,q)D^{V,W\alpha \kappa }(qp).$$ (161) We obtain $$\mathrm{\Pi }_{\mu \nu \rho \sigma }^V(p)=\kappa ^2\mathrm{exp}\left(p^2/\mathrm{\Lambda }_G^2\right)\frac{d^4q}{(q^2+m_V^2)[(qp)^2+m_V^2]}K_{\mu \nu \rho \sigma }(p,q)$$ $$\times \mathrm{exp}\left[(q^2+m_V^2)/\mathrm{\Lambda }_{\mathrm{SM}}^2\right]\mathrm{exp}\left\{[(qp)^2+m_V^2]/\mathrm{\Lambda }_{\mathrm{SM}}^2\right\},$$ (162) where in D-dimensions $$K_{\mu \nu \rho \sigma }(p,q)=p_\alpha p_\beta p_\rho p_\sigma +q_\alpha p_\beta p_\rho p_\sigma q_\alpha q_\beta p_\rho p_\sigma +(1D)q_\alpha q_\beta q_\rho p_\sigma $$ $$(1+D)p_\alpha q_\beta q_\rho q_\sigma +(D1)p_\alpha q_\beta p_\rho q_\sigma +Dq_\alpha q_\beta q_\rho q_\sigma .$$ (163) As usual, we must add to (10) the contributions from the fictitious ghost particle diagrams, the shadow field diagrams and the invariant measure diagram. We observe that from power counting of the momenta in the integral (10), we obtain $$\mathrm{\Pi }_{\mu \nu \rho \sigma }^V(p)\kappa ^2\mathrm{\Lambda }_{\mathrm{SM}}^4\mathrm{exp}\left(p^2/\mathrm{\Lambda }_G^2\right)N_{\mu \nu \rho \sigma }(p^2)$$ $$\frac{\mathrm{\Lambda }_{\mathrm{SM}}^4}{M_{\mathrm{PL}}^2}\mathrm{exp}\left(p^2/\mathrm{\Lambda }_G^2\right)N_{\mu \nu \rho \sigma }(p^2),$$ (164) where $`N(p^2)`$ is a finite remaining part of $`\mathrm{\Pi }^V(p)`$. We have as $`p^20`$: $$N_{\mu }^{}{}_{}{}^{\mu \sigma }{}_{\sigma }{}^{}(p^2)p^4.$$ (165) Thus, $`\mathrm{\Pi }_{\mu \nu \rho \sigma }^V(p)`$ vanishes at $`p^2=0`$ as it should because of gauge invariance and the massless graviton. For four-dimensional Euclidean momenta $`p^2\mathrm{\Lambda }_G^2`$, $`\mathrm{\Pi }_{\mu \nu \rho \sigma }^V(p)`$ is exponentially damped. At some value of the external graviton momentum $`p`$, when $`\mathrm{\Lambda }_{\mathrm{SM}}^4`$ could begin to become significant, the exponential damping suppresses this contribution. If we choose $`\mathrm{\Lambda }_G10^4`$ eV, then due to the damping of the gravitational vacuum polarization loop graph in the Euclidean limit $`p^2\mathrm{\Lambda }_G^2`$, the cosmological constant contribution is suppressed sufficiently to satisfy the bound (2), and it is protected from large unstable radiative corrections. Thus, FQFT provides a solution to the cosmological constant problem at the energy level of the standard model and possible higher energy extensions of the standard model. The universal fixed FQFT gravitational scale $`\mathrm{\Lambda }_G`$ corresponds to the fundamental length $`\mathrm{}_G1`$ cm at which virtual gravitational radiative corrections are cut off. We observe that the required suppression of the vacuum diagram loop contribution to the cosmological constant, associated with the vacuum energy momentum tensor at lowest order, demands a low fundamental energy scale $`\mathrm{\Lambda }_G10^4`$ eV, which controls the quantum gravity loop contributions. This is essentially because the external graviton momenta are close to the mass shell, requiring a low energy scale $`\mathrm{\Lambda }_G`$. This seems at first sight a radical suggestion that quantum gravity corrections are weak at energies higher than $`10^4`$ eV, but this is clearly not in contradiction with any known gravitational experiment. Indeed, as has been stressed in recent work on large higher dimensions, there is no experimental knowlege of gravitational forces below 1 mm. In fact , we have no experimental knowledge at present about the strength of graviton radiative corrections. The standard model experimental agreement is achieved for standard model particle states close to the mass shell. However, we expect that the dominant contributions to the vacuum density arise from standard model states far off the mass shell. In our perturbative quantum gravity theory, the tree graphs involving gravitons are identical to the tree graphs in local point graviton perturbation theory, retaining classical, causal GR and Newtonian gravity. In particular, we do not decrease the strength of the classical, large distance gravity force. In order to solve the severe cosmological constant hierarchy problem, we have been led to the surprising conclusion that, in contrast to the conventional folklore, quantum gravity corrections to the classical GR theory are negligible at energies above $`10^4`$ eV, a result that will continue to persist if our perturbative calculations can be extrapolated to near the Planck energy scale $`10^{19}`$ GeV. Since the cosmological constant problem already results in a severe crisis at the energies of the standard model, our quantum gravity resolution based on perturbation theory can resolve the crisis at the standard model energy scale and well beyond this energy scale. ## 11 Conclusions The ultraviolet finiteness of perturbative quantum field theory in four-dimensions is achieved by applying the FQFT formalism. The nonlocal quantum loop interactions reflect the quantum, non-point-like nature of the field theory, although we do not specify the nature of the extended object that describes a particle. Thus, as with string theories, the point-like nature of particles is “fuzzy” in FQFT for energies greater than the scale $`\mathrm{\Lambda }`$. One of the features of superstrings is that they provide a mathematically consistent theory of quantum gravity, which is ultraviolet finite and unitary. FQFT focuses on the basic mechanism behind string theory’s finite ultraviolet behavior by invoking a suppression of bad vertex behavior at high energies, without compromising perturbative unitarity and gauge invariance. FQFT provides a mathematically consistent theory of quantum gravity at the perturbative level. If we choose $`\mathrm{\Lambda }_G10^4`$ eV, then quantum radiative corrections to the classical tree graph gravity theory are perturbatively negligible to all energies greater than $`\mathrm{\Lambda }_G`$, provided that the perturbative regime is valid. The important gauge hierarchy problem, associated with the Higgs sector, is solved by the exponential damping of the Higgs self-energy in the Euclidean $`p^2`$ domain for $`p^2\mathrm{\Lambda }_H^2`$, and for a $`\mathrm{\Lambda }_H`$ scale in the electroweak range $`10^210^3`$ GeV. A damping of the vacuum polarization loop contributions to the vacuum energy density-gravity coupling at lowest order can resolve the cosmological constant hierarchy problem, if the gravity loop scale $`\mathrm{\Lambda }_G10^4`$ eV, by suppressing virtual gravitational radiative corrections above the energy scale $`\mathrm{\Lambda }_G`$. We must still set the physical scale $`\mathrm{\Lambda }_{\mathrm{YM}}`$, which controls the size of radiative loop corrections in the Yang-Mills sector of FQFT. We expect this scale to be much larger than the electroweak scale $`10^210^3`$ Gev, and it could be as large as grand unification theory (GUT) scales $`10^{16}`$ Gev, allowing for possible GUT unification schemes. Recently, new supernovae data have strongly indicated a cosmic acceleration of the present universe . This has brought the status of the cosmological constant back into prominence, since one possible explanation for this acceleration of the expansion of the universe is that the cosmological constant is non-zero but very small. We can, of course, accomodate a small non-zero cosmological constant by choosing carefully the gravity scale $`\mathrm{\Lambda }_G`$. Indeed, this new observational data can be viewed as a means of determining the size of $`\mathrm{\Lambda }_G`$. Our quantum field theory formalism has helped to resolve two critical hierarchy problems in modern physics, given two parameters $`\mathrm{\Lambda }_H10^210^3`$ GeV and $`\mathrm{\Lambda }_G10^4`$ GeV. These parameters will hopefully be explained by a more fundamental non-perturbative theory. Acknowledgments I thank Michael Clayton, John Donoghue, Holger Nielsen, Bob Holdom, George Leibbrandt, Michael Luke, Anton Kapustin, Pierre Savaria, Raman Sundrum and Gerard ’t Hooft for helpful and stimulating discussions. This work was supported by the Natural Sciences and Engineering Research Council of Canada.
warning/0003/astro-ph0003349.html	ar5iv	text	# A far IR study of the CfA Seyfert sample: I. The data 1footnote 11footnote 1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA. ## 1 Introduction The origin of the IR emission from Seyfert galaxies has been a matter of discussion since the early work of Rieke & Low (1972); Stein (1975); Rieke (1978); Neugebauer (1978), and others. These authors found an emission excess between 3 and 5 $`\mu `$m in Seyfert galaxies, that originated a strong controversy between those who supported that the emission was non thermal and those who defended that the excess was emission from dust heated by the nucleus. The use of efficient instruments to observe at 10 $`\mu `$m, and especially the pioneering work at the far IR carried out from the Kuiper Airborn Observatory (KAO) showed the importance of the mid and far IR emission to quantify the bolometric luminosity of Seyfert galaxies (Telesco & Harper 1980; Smith et al. 1983). However, it was the IRAS satellite that provided an extensive set of IR data for a large number of galaxies, from which it was shown that Seyfert galaxies are indeed strong far IR emitters (Rodríguez Espinosa, Rudy & Jones 1987; Edelson, Malkan & Rieke 1987; Spinoglio et al. 1995). The IRAS satellite was key to the understanding of the importance of the IR emission to the total luminosity of active galaxies. However the IRAS data alone are not sufficient to clarify the nature of the IR emission, as there are measurements only at a limited number of wavebands preventing a good definition of the shape of the Spectral Energy Distribution (SED) at mid and far IR wavelength (see, e.g., Telesco 1988; 1990; Bregman 1990). A proper characterization of the mid and far SED is essential to understand the emission mechanisms that produce the high output of Seyfert galaxies in the IR. Recent studies have discussed the origin of the IR emission suggesting that it is of thermal origin. For example Giuricin, Mardirossian & Mezzetti (1995) have studied a complete sample of Seyfert galaxies at 10 $`\mu `$m and propose that the emission is due to thermal reemission by dust. Bonatto & Pastoriza (1997), based on color studies of IRAS data from diverse Seyfert samples, find that the colors obtained can be explained with a combination of dust heated by the nucleus plus cold dust in the host galaxy. Siebenmorgen et al. (1997) show that the SED at IR and milimetric wavelengths can be modeled assuming that the dust heated by a central source dominates the luminosity output of these objects. Maiolino et al. (1998) confirmed this last result finding, in their high resolution IR images of the Circinus galaxy, an unresolved source, with size $`<`$ 1 pc, that is reprocessing the nuclear output via dust reradiation. Rigopoulou et al. (1997) have observed a sample of AGN in the CO mm line suggesting that the far IR emission of Seyfert galaxies is thermal, based on three different evidences: the correlation found between the far IR and the CO emission, the dependence of the far IR emission to hydrogen molecular mass ratio with dust temperature, and the similarity of the profile shape of the CO and HI lines. Another important issue is the understanding of the differences between the two Seyfert types. According to the unified models, Seyfert 2 nuclei are intrinsically similar to Seyfert 1 nuclei, the differences observed being due solely to geometrical effects. In Seyfert 2, neither the broad line region nor the optical, UV and soft X ray continuum can be observed directly because the central region is obscured by intervening material in the line of sight. Some authors argue that this material forms a sort of disc or torus of molecular material. This disc or torus is thought to be reponsible for the collimation of radiation and the observed anisotropies, i.e., biconic structures in emission line images (Simpson et al. 1997; Wilson et al. 1993) or highly collimated jets. It is expected that at sufficiently long wavelengths the optical depth of the torus would decrease and the differences between the two Seyfert types should disappear. Several tests have been made to ascertain the presence of these molecular tori (Heckman 1995; Pier & Krolik 1993). These obscuring tori have been theoretically modeled by Pier & Krolik (1992) and Granato & Danese (1994) among others, predicting that the mid IR optical depth is still considerable, thus it should be expected that Seyfert 1s are more luminous than Seyfert 2s in the mid IR. In this work, we make an attempt to understanding the origin of the mid and far IR emission from Seyfert galaxies through the study of their spectral energy distributions (SED). We present ISO data of the entire CfA Seyfert sample (Huchra & Burg 1992), consisting of 25 Seyfert 1 and 22 Seyfert 2, plus a LINER. Section 2 describes the observations. Section 3 presents the separation of the Spectral Energy Distributions in thermal components by means of the Inverse Planckian Transform. Section 4 describes the thermal emission components obtained from the inversion and in Section 5 we perform a statistical analysis of the parameters obtained and discuss the differences between the two Seyfert types. As the CfA Seyfert sample is a complete sample of Seyfert galaxies we expect that the results obtained here are statistically significant for all Seyfert galaxies and certainly suggestive for other classes of AGNs. ## 2 The data Observations of the CfA Seyfert sample have been carried out with the Infrared Space Observatory (ISO; Kessler et al. 1996) through filters at 16, 25, 60, 90 120, 135, 180 and 200$`\mu `$m. NGC 1068 was also observed at 4.9, 7.3 and 11.5 $`\mu `$m. The filter set was chosen to achieve good coverage in wavelength while producing a good sampling of the SEDs. The ISOPHOT instrument (Lemke et al. 1996) was used in the PHT-P and PHT-C configurations, with the P1, P2, C100 and C200 detectors. P1 and P2 are single element photodiodes. C100 is an array of 3x3 pixels, each one projecting onto 45 arcsec in the sky, while C200 is a 2x2 pixel array, each pixel projecting onto 89.4 arcsec in the sky. Integration times were calculated from the signal/noise ratio estimates produced by the ISOPHOT simulator based on interpolations and extrapolations of the IRAS data. The PHT-P observations were done in chopping mode through a 120 or 180 arcsec aperture depending on the size of the objects. The PHT-C observations were done in staring mode. In this case, to set the background level, an empty area adjacent to the source was observed prior to the actual source and with identical instrumental settings. For NGC 1068 we have done maps with C100, at 60 and 105 $`\mu `$m, moving by half a pixel (23”) between two contiguous map positions. Four additional objects (NGC 3079, NGC 3227, NGC 4051 and NGC 4151) have been mapped at 90 $`\mu `$m with the C100 array. Details of the mapping observations of these four galaxies are given in Pérez García, Rodríguez Espinosa & Fuensalida (2000). The data reduction was done with the help of the PHT-Interactive Analysis (PIA) V.7.0 tool. PIA first deglitchs the data eliminating the cosmic rays and others spurious effects, then corrects for non-linear effects taking into account both the dynamic range of the detectors and the flux level of the observed sources. To determine the signal PIA linearizes the integration ramps, and corrects for detector drifts. The background substraction is achieved by repeating the PIA reduction process for the background measurements files, which are then substracted from the object files. It is also necessary to correct for the PSF fraction that is not seen by the detector (see Lemke et al. 1996). This factor varies with wavelength, hence the actual correction factor for each filter is built-in in a PIA table. The flux calibration is performed using the standard calibration of PIA 7.0. For the PHT-P measurements the process is similar except that the background substraction is performed automatically from the chopped measurements and the flux calibration is done through the use of the internal FCS (Fine Calibration Source; Lemke et al. 1996) since it provides a better match with the actual intensity of the sources. Note that the P detector response is very dependent on the intensity of the source being observed. The photometric uncertainties are small varying between 1% for the objects with higher signal/noise ratio and 15% for the weakest objects. The final uncertainties are however dominated by the uncertainties in the flux calibration. The calibration factors have kept improving as the detector and instrument characteristics have become better known, although there are still some residual effects (non linearity of detectors) that are yet not adequately modelled. As of this writting, the uncertainty of the calibration is equal or better than 30%. Therefore, throughout this paper, we adopt a conservative uncertainty of 30%. The fluxes for the observed objects are given in Table 1 and 2. Fluxes are in Janskys. Some objects could not be observed in some or all of the filters for various reasons. These are also given in Table 1. Figure A far IR study of the CfA Seyfert sample: I. The data <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA. shows the Spectral Energy Distributions (SEDs) of the observed objects. The ISO data have been plotted together with IRAS data (Edelson, Malkan & Rieke 1987) at 12, 25, 60 and 100$`\mu `$m showing that the agreement of the ISO and IRAS data is very good in most cases. ## 3 The inversion of the SEDs It can be seen in fig. 1 that the SEDs describe a well defined energy range with a steady increase from the near IR on to a maximum between 90 and 135$`\mu `$m and the start of a decline toward longer wavelengths. Several results like those mentioned in the introduction point to the mid and far-IR emission of Seyfert galaxies being of thermal origin. Furthermore, we have shown that the SED of a few Seyfert galaxies can be explained as emission produced by two emissivity weighted blackbodies (Rodríguez Espinosa et al. 1996). However, rather than proceeding with a plain fit of a “ad hoc”number of black bodies, we have preferred to use an inversion method to analize the SEDs of the galaxies in the sample. This has the advantage that no assumptions have to be made as to the number or location of the sources responsible for the observed spectrum. In particular, we have used an Inverse Planckian Transform algorithm, that employs an emissivity ($`ϵ\lambda ^{1.5}`$) weighted Planck function kernel to switch from frequency space to the temperature domain, hence revealing the temperature distribution of the sources that originate the observed SEDs. The method applies Bayes theorem of conditional probability and the Richardson-Lucy iteration algorithm which converges quickly to a optimum result. To increase the number of data points used in the inversion algorithm, whenever available, we have added the four IRAS band fluxes at 12, 25, 60 and 100 $`\mu `$m as given by Edelson, Malkan & Rieke (1987). Furthermore, to avoid boundary convergence problems in the inversion algorithm we have used 10 $`\mu `$m ground-based data from several authors (Rieke 1978; Edelson, Malkan & Rieke 1987; Maiolino et al 1995), and 1.3 mm upper limit data from Edelson, Malkan & Rieke (1987). These additional data have been used solely for the purpose of constraining the inversion algorithm at the borders of the wavelength range of interest. Details of this method are given in Salas (1992) and in Pérez García, Rodríguez Espinosa, & Santolaya Rey (1998) and are not repeated here. Figure A far IR study of the CfA Seyfert sample: I. The data <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA. shows the results obtained after application of the Inverse Planckian Transform to the SEDs of the CfA Seyfert galaxies observed. For each object, the upper pannel in Figure A far IR study of the CfA Seyfert sample: I. The data <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA. plots the ISO data (triangles) and the IRAS data (filled squares) and the best fit to the mid and far IR data in heavy black, with the different spectral components contributing to the fit printed in dashed lines; the bottom pannel shows the temperature spectrum that produces these components. It is important to note that the temperature spectra obtained in this way are continuous. Further, the spectral temperatures of the dust grains group themselves in discrete and well defined features that we have called components. These components can be described by their peak temperature. Figure A far IR study of the CfA Seyfert sample: I. The data <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA. shows that 34 out of the 40 galaxies for which we have obtained good temperature spectra show three temperature components: a warm component with central temperature T$``$150 K, a cold component with T$``$40-70 K, and a colder component with T$``$15-25 K. The remaining 6 objects show only two components, although in NGC 1068 when the short wavelength data are added a third component is recovered. Table 3 shows the peak temperatures of these thermal components for each object. This table also gives the spectral and morphological type for each object. These results confirm the analysis made before for a few objects of the same sample (Pérez García, Rodríguez Espinosa, & Santolaya Rey 1998). The fits obtained from the inversion are, in general, very good. In some cases, weak objects or objects with low signal to noise ratios, the uncertainty in the inversion is higher. For some objects, we have rejected some measurements (the flux at 120$`\mu `$m of NGC 4235, and the fluxes at 90$`\mu `$m of Mrk 270, Mrk 461 and 1614+35) because these data points do not follow the shape of the SEDs defined by the rest of the measurements. In all of these cases the rejected measurement is very noisy. Four other objects (Mrk 590, Mrk 573 and Mrk 841) have not been measured at long wavelengths ($`>`$90$`\mu `$m), because they are below the sensitivity limit for detection with ISO. Nevertheless, the method still finds three components for three out of the except for NGC 1068, that will be analized in a separate section. four objects. For Mrk 590, the inversion does not converge. In four of the objects, the warm component peak temperature obtained from the inversion is not high enough to fit the data. These four objects are Mrk 335, NGC 5273, Mrk 1243 and IZw1. Both in Mrk 335 and in NGC 5273 it can be due to the low signal to noise ratio of the data. NGC 5273 shows a discrepancy between the IRAS value at 12 $`\mu `$m and the ISO value at 16 $`\mu `$m (again very likely due to the low signal to noise ratio of the ISO data point) that prevents a good fit to the data. Mrk 1243 has only three data points in the hot part of the spectrum. IZw1 is a special case, since the short wavelength data are of good signal to noise ratio, however it is not detected at 200 $`\mu `$m and the flux at 180 $`\mu `$m is rather low. The fall of the signal at longer wavelengths is fast enough to explain that the object is not detected at 200 $`\mu `$m. Figure A far IR study of the CfA Seyfert sample: I. The data <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA. shows that the SED of IZw1 if made of just two components does not account for all of the observed emission. If we substract this initial two-component fit from the actual SED, the residual seems to correspond to yet one third hotter thermal component. The inversion method does not succeed at identifying completely this third component, due to the temperature of this third component being too high, therefore only a small fraction of the energy within this third component falls within the range of wavelength studied. The shape of the residual is however similar to the other two components, a clear hint of the existance of yet another thermal component with a peak temperature of $``$300 K. This temperature is the highest temperature found among all of the SEDs. The cold and very cold components also have peak temperatures higher than the corresponding components of the SEDs of the rest of the CfA sample. IZw1 is one of the more distant and brighter Seyfert galaxies known. Its different spectroscopic characteristics of line width and line ratios were pointed out by Osterbrock and Pogge (1985) and Goodrich (1989). Recently, several authors have pointed out the different nature of IZw1. For example, the optical and X-ray spectra are steeper in IZw1 than in normal Seyfert 1 and 2. Indeed IZw1 shows fast and large amplitude variability in its X-ray emission (see Halpern & Moran 1998 and references there in), and is possibly better associated with the broad line quasars (Boroson & Meyers 1992; Lawrence et al. 1997). ## 4 The dust components The inversion method assigns to each dust component a range of temperature, i.e., the dust grains have temperatures ranging between a minimum and a maximum temperature. This is indeed an expected and physically meaningful result, as the dust closer to the source or sources of radiation will be warmer than the dust farther away. The existence of temperature components indicates that there are well defined sources with well defined temperature profiles which must be the result of well defined physical scenarios existing in these objects. From the temperature spectra, we can calculate the flux that each dust component contributes to the total mid and far IR emission. For each component, the flux is obtained by integration of the temperature spectrum over the relevant range of temperatures: $$F_i=_{Tmin_i}^{Tmax_i}\mathrm{\Psi }(T)𝑑Ti=1,2,3$$ (1) where T$`_{min_i}`$ and T$`_{max_i}`$ are the extreme values of the temperature range that defines each component, and $`\mathrm{\Psi }(T)`$ is the temperature distribution obtained from the inversion process. Table 4 shows the fluxes obtained in this way for each component and the total flux for each object. The comparison between the far IR IRAS fluxes (FIR), as calculated from the values at 60 and 100 $`\mu `$m (Lonsdale et al. 1984) and the sum of the cold and very cold components calculated here is in good agreement. The luminosities for each of the components are shown in Table 5. We have used $`H_0`$ = 75 Km s<sup>-1</sup> Mpc<sup>-1</sup> throughout this paper. The key point now is wheather these emission components can be physically explained within a sensible scenario. In what follows we review the pieces of evidence that we have for explaining each of components: a) In a previous paper Rodríguez Espinosa & Pérez García (1997) were able to use optical R band images of a subset of low redshift Seyferts of the present sample to separate the fluxes from the central region from those of the galaxy disks. We found a very good correlation between the ratio of the extended to the compact R band fluxes and the ratio of the cold plus very cold component to the warm emission component. The conclusion was that the warm emission component is related with the central regions of the galaxies while the cold and the very cold emission had to originate in the disk of these galaxies. b) Furthermore, a correlation has been found between the flux produced by the warm component and the flux in the high ionization coronal lines fluxes like \[OIV\]$`\lambda `$25.9 and \[NeV\]$`\lambda `$14.3$`\mu `$m (Prieto, Pérez García & Rodríguez Espinosa 2000), indicating that the warm component must be heated by the nucleus of these galaxies. c) In another recent paper, Pérez García et al. (2000) have shown that based on 90$`\mu `$m ISO maps of four nearby Seyferts, the 90$`\mu `$m emission is physically extended up to radii similar or larger than those seen in optical images of these same galaxies. Furthermore, the extension of this 90$`\mu `$m emission has been characterized and its extension, scale length and surface brightness profiles are typical of normal galaxy disks. Based on the above, a scenario arises in which the warm component is associated with dust heated by radiation coming from the nuclear or circumnuclear regions of these galaxies, while the cold and very cold dust must be heated by process occurring in the galaxy disk. Danese et al. (1992) also conclude that the mid IR emission (10-25$`\mu `$m) is dominated by the nucleus or the circumnuclear region. Further support to this scenario comes from the following: * Warm dust. Its characteristic peak temperature is in the range 120-170 K, a range of temperatures warmer tham is normal of dust in typical starforming regions. The nuclear origin of the radiation responsible for the heating of the warm dust was already indicated by Rudy (1984) who found a correlation between the \[OIII\]$`\lambda `$5007 emission line flux and the 10 $`\mu `$m emission for a sample of quasars, Seyfert galaxies and radiogalaxies. This author suggested that the dust responsible for the 10 $`\mu `$m emission is mixed with the ionized gas of the Narrow Line Region (NLR) that produces the \[OIII\]$`\lambda `$5007 emission line. More recently, Giuricin, Mardiossian and Mezzetti (1995) support the idea of the nuclear heating of this warm dust, based on 10 $`\mu `$m small aperture observations of a sample of 100 galaxies. They found that the 10 $`\mu `$m emission correlates very well with the 25 $`\mu `$m IRAS emission, while the correlation is poor with the 60 $`\mu `$m emission. From a different perspective Heckmann et al. (1997) have found in the Seyfert 2 galaxy Mrk 477 a strong starburst with very warm dust in an scale of a few hundred parsecs, that confirms the idea that the warm dust can also be heated by nuclear and circumnuclear starformation regions. * Cold dust. Peak temperatures for this component range between 40 and 70 K, a range of temperatures that is typical of dust in regions of starformation. Note that cold dust is present in mostly all classes of galaxies, including normal and starburst galaxies (Knapp et al. 1996; Chini, Kruger & Kreysa 1992; Klaas et al 1997; Walterbos & Schwering 1987). In all these galaxy types, the heating of the dust is produced by OB stars in star formation regions in the galaxy disks. The higher temperatures corresponding to higher recent starformation rates (Young et al. 1989) * Very cold dust. Peak temperatures for this component range from around 15 to 25 K, temperatures that are normal of dust heated by the general interstellar radiation field. This very cold dust is tipically composed of big grains in thermal equillibrium with the interstellar radiation field and has been observed in normal spiral galaxies (see, e.g., Walterbos & Greenewalt 1996; Walterbos & Schwering 1987; Cox, Kruger & Metzer 1986). For example, Cox, Kruger & Mezger (1986) predict that the dust of the Galaxy is at a temperature of around 15-20 K, and Walterbos & Schwering show that in the disk of M31 there is very cold dust at 21 K. ### 4.1 NGC 1068 NGC 1068 has been observed at a wider spectral range than the rest of the sample, 11 filters between 4 and 200 $`\mu `$m have been employed. NGC 1068 is indeed one of the best studied Seyfert 2 galaxies in all spectral ranges. Morphologically it is clasified as SAb, and was the first galaxy in which broad emission lines in polarized light were found (Miller & Antonucci 1983; Antonucci and Miller 1985). In the IR range, Tresh-Frenberg et al. (1987) have found that the emission between 5 and 20 $`\mu `$m presents an elongated structure in the nucleus, that they explain as thermal emission by dust heated by both the active nucleus and an active starforming region near (1”-2”) the nucleus. The comparison between 10$`\mu `$m images and HST images (Cameron et al. 1993) indicates a correlation between the hot dust emission and the Narrow Line Region. These autors developed a model that replaces the molecular torus by giant molecular clouds of gas and dust that attenuate the broad lines. Braatz et al. (1993) found a correlation between a 12.4 $`\mu `$m image and both optical continuum and \[OIII\] images; they also observed that the IR emission is aligned perpendicular to the torus plane and infered that the extended IR emission is thermal radiation by dust in molecular clouds heated by collimated nuclear emission, without ruling out a possible contribution by dust in the molecular torus. Telesco & Harper (1980) have studied the far IR emission of NGC 1068 with KAO (Kuipern Airborn Observatory) observations between 30 and 300 $`\mu `$m and have found two thermal components with temperatures of $``$36 K and $``$115 K. We have applied the inversion method explained before in the same spectral range that we have used for the rest of the objects in the sample, i.e., between 12 and 200 $`\mu `$m. Figure A far IR study of the CfA Seyfert sample: I. The data <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA. shows the results obtained, i.e., two thermal components, a cold one with peak temperature of T$``$36 K and a warm one with T$``$115 K in the peak. It remains however to extend the fit to the shorter wavelength data which we have for NGC 1068 alone. When we apply the inversion method to the full range observed (4-200 $`\mu `$m) the SED of NGC 1068 appears now as the sum of three thermal components (see Figure A far IR study of the CfA Seyfert sample: I. The data <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA.). Two of them can be identified with the two thermal components previously found before without considering the 4-12 $`\mu `$m range. The peak temperatures of the three components are now 29, 110 and 278 K. The differences between the flux contained in each of the two components common to both inversions is less than 7%. The three components found in NGC 1068 are explained as in the rest of the objects of the CfA sample. The cold component is produced by dust in the host galaxy, the warm component corresponds to dust heated by the active nucleus and/or circumnuclear starforming regions. The dust at higher temperature is also heated by the active nucleus, and should either be located in the inner part of the torus or be heated by active starforming regions placed in a radius $`<`$100 pc from the active nucleus (González Delgado et al. 1998). The flux enclosed under the component with T$``$278 K is F = 9.17 10<sup>-9</sup> erg cm<sup>-2</sup> s<sup>-1</sup>, and the corresponding luminosity L = 2.40 10<sup>45</sup> erg s<sup>-1</sup>, while the total IR luminosity of NGC 1068 is 8.60 10<sup>45</sup> erg s<sup>-1</sup>. Hence this component represents 28% of the total IR luminosity, which agrees with the fraction of the total UV emission that is produced by circumnuclear starforming regions in a sample of Seyfert 2 galaxies (González Delgado et al. 1998). ## 5 Energy balance: The cold emission component One possible test that can confirm the suggestion made before on the relation between the cold emission component and the presence of starforming regions in the galaxy discs consists of testing the ratio of H$`\alpha `$ emission to far IR output. For instance in NGC 3079, the observed H$`\alpha `$ luminosity is 7.93 10<sup>40</sup> erg sec<sup>-1</sup> (Armus, Heckman & Miley 1990), a figure that includes both the nuclear and extended H$`\alpha `$ luminosity. Moreover, with a ratio of extended to nuclear H$`\alpha `$ emission of 0.77 (Armus, Heckman & Miley 1990), the extended H$`\alpha `$ emission amounts to 3.49 10<sup>40</sup> erg sec<sup>-1</sup>, which shows the importance of the H$`\alpha `$ flux produced in star forming regions outside the central region in this object. If we take as an average extinction the value measured by Hawarden et al. (1995) A<sub>V</sub>=7.5 mag, the extinction at $`H_\alpha `$ is A$`_{H_\alpha }`$=6.1 mag, and the extinction corrected extended H$`\alpha `$ luminosity is 0.93 10<sup>43</sup> erg sec<sup>-1</sup>. This is to be compared with the IR output due to dust in starforming regions that competes with the gas for high energy photons from massive stars. In NGC 3079 this is what we have called the cold emission (Pérez García, Rodríguez Espinosa & Santolaya Rey 1998), which amounts to 3.52 10<sup>44</sup> erg sec<sup>-1</sup>, or a ratio of H$`\alpha `$ to cold FIR emission of 0.03. This value is in agreement with the expected value in regions where massive stars are forming (Devereux & Young 1990). Hence for NGC 3079 there is good agreement between the H$`\alpha `$ output from disk HII regions and the cold IR emission, which confirms the results in Pérez García et al. (1998) and in this work indicating that the cold component emission in Seyfert galaxies originates in starforming regions wjthin their disk. In the case of NGC 4051 we consider the ratio between H$`\alpha `$ and far IR luminosity in two limiting cases. First, we consider the total integrated H$`\alpha `$ luminosity, 4.34 10<sup>41</sup> erg sec<sup>-1</sup> (Romanishin 1990), which after correction for extinction (A<sub>V</sub>= 0.24, Ho et al. 1997) becomes 5.22 10<sup>41</sup> erg sec<sup>-1</sup>. Hence, the ratio of total H$`\alpha `$ to cold far IR emission (0.25 10<sup>44</sup> erg sec<sup>-1</sup>) is 0.02. This is an upper limit as the total H$`\alpha `$ emission includes the nuclear emission which to a certain extent can be due to non-thermal processes associated with the AGN. On the other hand, if we consider the H$`\alpha `$ emission coming exclusively from the star-forming regions in the disk of NGC 4051, which according to González Delgado et al. (1997) is 26% of the total emission we get a ratio of extended H$`\alpha `$ to cold far IR emission of 0.006. This is a lower limit as some of the H$`\alpha `$ flux that we have assumed nuclear is very likely due to extended circumnuclear emission that will also contribute to the cold far IR emission. The actual value for the ratio of extended H$`\alpha `$ flux to cold far IR emission ranges therefore between 0.006 and 0.02, a value that is well within the range given by Devereux & Young (1990) as typical of normal HII regions in the disk of spiral galaxies. In the case of NGC 5033 the integrated H$`\alpha `$ flux is 3.98 10<sup>41</sup> erg sec<sup>-1</sup> (Kennicutt 1983), which after correcting for extinction (A<sub>V</sub>= 1.48, Ho et al. 1997) becomes 1.20 10<sup>42</sup> erg sec<sup>-1</sup>. The H$`\alpha `$ emission to cold far IR emission ratio is 0.02, which is again to be taken as an upper limit given that the total H$`\alpha `$ flux includes the nuclear emission which we can not separate because we do not have the necessary data, However this upper limit is similar to that obtained for NGC 4051. As for the rest of the sample a similar exercise with H$`\alpha `$ fluxes obtained from the literature produce similar results. In very few instances have we been able to separate the extended from the nuclear H$`\alpha `$ emission. For these cases the ratio of extended H$`\alpha `$ emission to cold far IR emission ranges between 1. 10<sup>-5</sup> for NGC 5929 and 0.02 for Mrk 270. In most cases this separation between extended and nuclear emission has not been possible and hence ratios between the total H$`\alpha `$ emission and the cold far IR output have been obtained. These range between 1.5 10<sup>-4</sup> (Mrk 766) and 2. 10<sup>-3</sup> (NGC 3516). In all cases it can be seen that the ratios are within those expected if massive stars are responsible for the ionizing photons as well as for the heating of the dust. We therefore conclude that the cold emission component that appeared naturally from the inversion process is related to radiation from dust in starforming regions in the galaxy discs. ## 6 Seyfert type differences The current understanding of Seyfert galaxies assumes that the differences between the two Seyfert types are due to geometric factors in the nuclear region. If this is so we should look for differences in the warm emission component between the Seyfert types. Let us therefore turn our attention to the warm dust component that we have suggested it is directly related to the nuclear and circumnuclear dust emission. If there is a molecular torus obscuring the nuclear region and this is sufficiently thick (as the models predict, see e.g. Pier & Krolik 1992; Granato & Danese 1994) we shall observe differences in the warm component emission produced by each of the two Seyfert types. In order to compare the parameters obtained from the inversion of the SEDs for the different Seyfert types, we establish two different groups of objects within the CfA sample: Seyfert 1s, including galaxies classified as Seyfert 1.5, and Seyfert 2s, including Seyfert 1.8 and 1.9. Therefore, the sample gets divided into 22 Seyfert 1 and 24 Seyfert 2. If we drop those objects not observed or those whose data do not allow obtaining a well defined warm component, the sample with adequate data for this analysis consists of 18 Seyfert 1 galaxies and 22 Seyfert 2 galaxies. If the dust responsible for the warm emission component is located within the molecular torus that hides both the active nucleus and the broad line region, we expect to find differences between some of the parameters that define the warm component of the Seyfert 1 and 2 objects. If however the dust is located outside the molecular torus, then the geometry does not play anymore a key role and we should not see differences between the two types of Seyfert galaxies. ### 6.1 Temperature differences Regarding the temperature distribution, it should be realized that in both Seyfert classes we are sampling dust emission produced under similar physical conditions, probed with the same wavelength range, and heated in either case by the active nucleus or by circumnuclear starforming regions. Therefore as temperature is an intensive quantity as opposed to extensive, we expect not to see any difference between the average temperature of the warm component in the two Seyfert types. Figure A far IR study of the CfA Seyfert sample: I. The data <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA. shows the distribution of warm component temperature for the two groups of Seyfert galaxies. It is seen that the median temperature of the Seyfert 1s is slightly higher than the median of the Seyfert 2s though the differences are small, IZw1 excepted. As we pointed out before, IZw1 is a special object, its warm temperature being way far out of every other object in the sample. We have therefore dropped IZw1 from the statistics. To understand the significance of the differences between the two distributons, we have applied a Kolmogorov-Smirnov (KS) test. The characteristics of the distributions are: $$\begin{array}{ccc}T_1\hfill & =148\hfill & \sigma =12.\hfill \\ T_2\hfill & =144\hfill & \sigma =14.\hfill \end{array}$$ The result of this KS test is that there is a 78% probability that the two distributions are the same and so the differences between them are not significant and both groups (Seyfert 1 and Seyfert 2) enjoy, as expected, the same temperature distribution. ### 6.2 Fluxes and luminosities Recent studies of the mid IR emission from Seyfert galaxies claim that Seyfert 2 galaxies are weaker than Seyfert 1s (Heckman 1995; Maiolino et al. 1995; Giuricin, Mardirossian & Mezzeti 1995; Mulchaey et al. 1994). This result is interpreted within the framework of the unified models as an anisotropy, resulting from the presence of a molecular torus with a given optical thickness in the IR. To test this claim we have compared the fluxes and luminosities of the warm component of the Seyfert SEDs to see if we find significant differences between Seyfert 1 and Seyfert 2s. The median values for the warm luminosities are: $$\begin{array}{ccc}logL_{warm}_1\hfill & =44.7\hfill & \sigma =0.3\hfill \\ logL_{warm}_2\hfill & =44.9\hfill & \sigma =0.8\hfill \end{array}$$ Statistically, the KS test shows that there is a probability 33% that both distributions are the same. Hence there are no significant differences between the Seyfert 1 and 2’s regarding the warm emission. Turning now to the total far IR luminosities, i.e., the sum of the three emission components, the distributiones are: $$\begin{array}{ccc}logL_{IR}_1\hfill & =45.1\hfill & \sigma =0.8\hfill \\ logL_{IR}_2\hfill & =44.9\hfill & \sigma =0.7\hfill \end{array}$$ In this case the distributions are the same with a significance level of 99% therefore it can not be concluded that the Seyfert 1s are brighter than the Seyfert 2s in the mid and far IR. We now turn to search wheather there are differences between Seyfert 1 and Seyfert 2s in the ratio of warm component to total luminosity. First, we refer to the ratios between the warm and total fluxes. The mean and standard deviations of these distributions are: $$\begin{array}{ccc}F_{warm}/F_{IR}_1\hfill & =0.42\hfill & \sigma =0.17\hfill \\ F_{warm}/F_{IR}_2\hfill & =0.28\hfill & \sigma =0.15\hfill \end{array}$$ Figure A far IR study of the CfA Seyfert sample: I. The data <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA. shows the flux ratio distributions. A KS test indicates that the distributions are different at the 99% level of significance. This result suggests that the Seyfert 1s emit fractionally more than the Seyfert 2s in the mid IR (warm component). This is consistent, for example, with the models of circumnuclear tori of Pier & Krolik (1992;1993) and of Granato & Danese (1994), that predict anisotropy by absorption of nuclear emission in the mid IR. However this result should be taken with care, as it could indicate either that the nuclear emission is larger in the Seyfert 1s relative to their total FIR flux or that the contribution of the host galaxy is stronger in the Seyfert 2s. To discriminate between these two possible explanations we must normalize the infrared fluxes with an isotropic property of the galaxies in the sample, i.e., with fluxes emitted at long enough wavelengths that they are not suspect of suffering extinction and thus do not depend on the geometry of the sources. Other authors have used \[OIII\]$`\lambda `$5000 $`\AA `$ fluxes, hard X-ray fluxes or radio fluxes (see, for example, Mulchaey et al. 1994). However, the \[OIII\] $`\lambda `$5000 $`\AA `$ fluxes can be affected by absorption due to dust in the NLR. We prefer to use 20 cm radio emission fluxes to normalize the IR flux, since the radio emission is not affected by selective extinction. We have used integrated radio data from Edelson (1987). These data consists of VLA observations at 1.46 GHz (20 cm) with a bandwidth of 45 MHz. The FWHM beamwidth used is 1.5 arcmin, directly comparable with our ISO data. The distributions of the radio normalized warm IR fluxes show the following characteristics: $$\begin{array}{ccc}log(F_{warm}/F_{20cm})_1\hfill & =6.5\hfill & \sigma =0.3\hfill \\ log(F_{warm}/F_{20cm})_2\hfill & =6.1\hfill & \sigma =0.3\hfill \end{array}$$ KS tests indicate that both distributions are different at a significance level of 99.9%. Therefore, the warm flux is indeed higher in Seyfert 1s than it is in Seyfert 2s, and it can be concluded that the warm emission from Seyfert 2s is affected by dust extinction to a larger extent that in the Seyfert 1 galaxies. This result suggests that at shorter wavelengths (mid IR) the emission is still anisotropic, in agreement with the molecular torus models of Pier & Krolik (1992;1993) and Granato & Danese (1994). These and others authors have proposed different models for the absorbing material. Pier & Krolik (1992;1993), Granato & Danese (1994) and Efstathiou & Rowan-Robinson (1994) have modeled the absorbing structures as axially symmetric tori. The models proposed by Granato & Danese (1994) result in thin and extended tori with optical depths ranging from $`\tau `$ 10 to 300 in the UV band and maximum radii ranging from tens to hundreds of parsecs. On the contrary, the models proposed by Pier & Krolik (1993) show thin and compact accretion disks with very large optical depths with values of $`\tau \stackrel{<}{}`$ 1000 in the UV band, and compact radii with dimension of a few pc. If we consider the Seyfert 1 as canonical unobscured objects, and abscribe the differences found between the two Seyfert types to absorption by the obscuring torus we obtain a mid IR optical depth of $`\tau _{IR}`$ 0.4 or a $`\tau _{UV}`$ 80 (assuming $`\tau _\lambda \mathrm{\hspace{0.33em}1}/\lambda `$) for the Seyfert 2 objects. This value is indeed very mild and within the range predicted for the thin and extended tori of Granato, Danese & Franceschini (1997). The validity of ratioing with the 20 cm radio flux has been however questioned based on the idea that Seyfert 2 galaxies may have more star formation in their disks than Seyfert 1s (Maiolino et al. 1995). This would affect the radio emission, and hence the warm to radio flux ratio. It remains to test wether the Seyfert 2 galaxies are indeed stronger emitters of extended far IR radiation. We consider for this test the ratio of the cold far IR component to 20 cm radio flux. For the two groups the values are: $$\begin{array}{ccc}log(F_{cold}/F_{20cm})_1\hfill & =6.7\hfill & \sigma =0.7\hfill \\ log(F_{cold}/F_{20cm})_2\hfill & =6.5\hfill & \sigma =0.5\hfill \end{array}$$ These distributions are similar (60% probability) according to the KS test. Therefore, Seyfert 2s are not stronger emitters in extended IR radiation than Seyfert 1s, and there are no reasons to suspect that they should be stronger radio emitters based solely on the amount of star formation occurring in Seyfert 2s. It is also worth pointing out that the differences found between the type 1 and 2 Seyferts are restricted to the warm emission, while there are not differences regarding the cold and very cold emission, i.e.,the emission from their respective galaxy disks. ## 7 Summary We have presented far IR photometry with ISO of the CfA Seyfert sample. The data have allowed a detailed study of the far IR SED of these sources using a Bayesian inversion method. It has been shown that the mid and far IR emission of Seyfert galaxies can be explained by the emission of three thermal components, a warm component, associated with dust heated by the nucleus and circumnuclear starformation regions; a cold dust component heated by star forming region in the galaxy disk, and very cold dust component heated by the general interstellar radiation field. The mid to far IR output from Seyfert galaxies does not have a simple origin but different ingredients play an important role in it. The comparison of cold far IR fluxes with H$`\alpha `$ data confirms that the cold emission component that appeared naturally from the inversion process is related to radiation from dust in starforming regions in the galaxy discs. We find that the mid IR emission (warm component) is larger in Seyfert 1 than in 2s, suggesting the presence of obscuring material in Seyfert 2s. The median value obtained for the optical depth is in the range predicted by the thin and extended tori models.
warning/0003/math-ph0003028.html	ar5iv	text	# A Fresh Look at Entropy and the Second Law of Thermodynamics ## The basic question The paradigmatic event that the second law deals with can be described as follows. Take a macroscopic system in an equilibrium state $`X`$ and place it in a room, together with a gorilla equipped with arbitrarily complicated machinery (a metaphor for the rest of the universe), and a weight—and close the door. As in the old advertisement for indestructible luggage, the gorilla can do anything to the system—including tearing it apart. At the end of the day, however, the door is opened and the system is found in some other equilibrium state, $`Y`$, the gorilla and machinery are found in their original state, and the only other thing that has possibly changed is that the weight has been raised or lowered. Let us emphasize that although our focus is on equilibrium states, the processes that take one such state into another can be arbitrarily violent. The gorilla knows no limits. The question that the second law answers is this: What distinguishes those states $`Y`$ that can be reached from $`X`$ in this manner from those that cannot? The answer: There is a function of the equilibrium states, called entropy and denoted by $`S`$, that characterizes the possible pairs of equilibrium states $`X`$ and $`Y`$ by the inequality $`S(X)S(Y)`$. The function can be chosen to be additive (in a sense explained below), and with this requirement it is unique, up to a change of scale. Our main point is that the existence of entropy relies only on a few basic principles, independent of any statistical model—or even of atoms. What is exciting about this apparently innocuous statement is the uniqueness of entropy, for it means that all the different methods for measuring or computing entropy must give the same answer. The usual textbook derivation of entropy as a state function, starting with some version of “the second law”, proceeds by considering certain slow, almost reversible processes (along adiabats and isotherms). It is not at all evident that a function obtained in this way can contain any information about processes that are far from being slow or reversible. The clever physicist might think that with the aid of modern computers, sophisticated feedback mechanisms, unlimited amounts of mechanical energy (represented by the weight) and lots of plain common sense and funding, the system could be made to go from an equilibrium state $`X`$ to $`Y`$ that could not be achieved by the primitive quasistatic processes used to define entropy in the first place. This cannot happen, however, no matter how clever the experimentalist or how far from equilibrium one travels! What logic lies behind this law? Why can’t one gorilla undo what another one has wrought? The atomistic foundation of this logic is not as simple as is often suggested. It not only concerns things like the enormous number of atoms involved ($`10^{23}`$), but also other aspects of statistical mechanics that are beyond our present mathematical abilities. In particular, the interaction of a system with the external world (represented by the gorilla and machinery) cannot be described in any obvious way by Hamiltonian mechanics. Although irreversibility is an important open problem in statistical mechanics, it is fortunate that the logic of thermodynamics itself is independent of atoms and can be understood without knowing its source. The founders of thermodynamics—Rudolf Clausius, Lord Kelvin, Max Planck, Constantin Carathéodory, and so on—clearly had transitions between equilibrium states in mind when they stated the law in sentences such as “No process is possible, the sole result of which is that a body is cooled and work is done” (Kelvin). Later it became tacitly understood that the law implies a continuous increase in some property called entropy, which was supposedly defined for systems out of equilibrium. The ongoing, unsatisfactory debates (see referce \[9, for example) about the definition of this nonequilibrium entropy and whether it increases shows, in fact, that what is supposedly “easily” understood needs clarification. Once again, it is a good idea to try to understand first the meaning of entropy for equilibrium states—the quantity that our textbooks talk about when they draw Carnot cycles. In this article we restrict our attention to just those states; by “state” we always mean “equilibrium state”. Entropy, as the founders of thermodynamics understood the quantity is subtle enough, and it is worthwhile to understand the “second law” in this restricted context. To do so it is not necessary to decide whether Boltzmann or Gibbs had the right view on irreversibility. (Their views are described in Joel L. Lebowitz’s article “Boltzmann’s Entropy and Time’s Arrow”, Physics Today, September 1993, page 32.) ## The basic concepts To begin at the beginning, we suppose we know what is meant by a thermodynamic system and equilibrium states of such a system. Admittedly these are not always easy to define, and there are certainly systems, such as a mixture of hydrogen and oxygen or an interstellar ionized gas, that can behave as if they are in equilibrium even if it is not truly so. The prototypical system is a “simple system”, consisting of a substance in a container with a piston. But a simple system can be much more complicated than that. Besides its volume it can have have other coordinates, which can be changed by mechanical or electrical means—shear in a solid, or magnetization, for example. In any event, a state of a simple system is described by a special coordinate $`U`$, which is its energy, and one or more other coordinates (such as the volume $`V`$) called work coordinates. An essential point is that the concept of energy, which we know about from the moving weight and Newtonian mechanics, can be defined for thermodynamic systems. This fact is the content of the first law of thermodynamics. Another type of system is a “compound system”, which consists of several different or identical independent, simple systems. By means of mixing or chemical reactions, systems can be created or destroyed. Let us briefly discuss some concepts that are relevant for systems and their states, which are denoted by capital letters such as $`X,X^{},Y,\mathrm{}`$. Operationally, the composition, denoted $`(X,X^{})`$, of two states $`X`$ and $`X^{}`$ is obtained simply by putting one system in a state $`X`$ and one in a state $`X^{}`$ side by side on the experimental table and regarding them jointly as a state of a new, compound system. For instance, $`X`$ could be a glass containing 100 g of whiskey at standard pressure and $`20^{}`$ C, and $`X^{}`$ a glass containing 50 g of ice at standard pressure and $`0^{}`$ C. To picture $`(X,X^{})`$ one should think the two glasses standing on a table without touching each other. Another operation is the “scaling” of a state $`X`$ by a factor $`\lambda >0`$, leading to a state denoted $`\lambda X`$. Extensive properties like mass, energy and volume are multiplied by $`\lambda `$, while intensive properties such as pressure stay intact. For the states $`X`$ and $`X^{}`$ as in the example above the example above $`\frac{1}{2}X`$ is 50 g of whiskey at standard pressure and $`20^{}`$ C, and $`\frac{1}{5}X^{}`$ is 10 g of ice at standard pressure and $`0^{}`$ C. Compound systems scale in the same way: $`\frac{1}{5}(X,X^{})`$ is 20 g of whiskey and 10 g of ice in separate glasses with pressure and temperatures as before. A central notion is adiabatic accessibility. If our gorilla can take a system from $`X`$ to $`Y`$, as described above—that is, if the only net effect of the action, besides the state change of the system, is that a weight has possibly been raised or lowered, we say that $`Y`$ is adiabatically accessible from $`X`$ and write $`XY`$ (the symbol $``$ is pronounced “precedes”). It has to be emphasized that for macroscopic systems this relation is an absolute one: If a transition from $`X`$ to $`Y`$ is possible at one time, then it is always possible (that is, it is reproducible), and if it is impossible at one time it never happens. This absolutism is guaranteed by the large powers of 10 involved—the impossibility of a chair’s spontaneous jumping up from the floor is an example. ## The role of entropy Now imagine that we are given a list of all possible pairs of states $`X,Y`$ such that $`XY`$. The foundation on which thermodynamics rests, and the essence of the second law, is that this list can be simply encoded in an entropy function $`S`$ on the set of all states of all systems (including compound systems) so that when $`X`$ and $`Y`$ are related at all, then $$XY\text{if and only if}S(X)S(Y).$$ Moreover, the entropy function can be chosen in such a way that if $`X`$ and $`X^{}`$ are states of two (different or identical) systems, then the entropy of the compound system in this pair of states is given by $$S(X,X^{})=S(X)+S(X^{}).$$ This additivity of entropy is a highly nontrivial assertion. Indeed, it is one of the most far reaching properties of the second law. In compound systems such as the whiskey/ice example above, all states $`(Y,Y^{})`$ such that $`XY`$ and $`X^{}Y^{}`$ are adiabatically accessible from $`(X,X^{})`$. For instance, by letting a falling weight run an electric generator one can stir the whiskey and also melt some ice. But it is important to note that $`(Y,Y^{})`$ can be adiabatically accessible from $`(X,X^{})`$ without $`Y`$ being adiabatically accessible from $`X`$. Bringing the two glasses into contact and separating them again is adiabatic for the compound system but the resulting cooling of the whiskey is not adiabatic for the whiskey alone. The fact that the inequality $`S(X)+S(X^{})S(Y)+S(Y^{})`$ exactly characterizes the possible adiabatic transitions for the compound system, even when $`S(X)S(Y)`$, is quite remarkable. It means that it is sufficient to know the entropy of each part of a compound system in order to decide which transitions due to interactions between these parts (brought about by the gorilla) are possible. Closely related to additivity is extensivity, or scaling of entropy, $$S(\lambda X)=\lambda S(X),$$ which means that the entropy of an arbitrary mass of a substance is determined by the entropy of some standard reference mass, such as 1 kg of the substance. Without this property engineers would have to use different steam tables each time they designed a new engine. In traditional presentations of thermodynamics, based for example on Kelvin’s principle given above, entropy is arrived at in a rather roundabout way which tends to obscure its connection with the relation $``$. The basic message we wish to convey is that existence and uniqueness of entropy are equivalent to certain simple properties of the relation $``$. This equivalence is the concern of . An analogy leaps to mind: When can a vector-field, $`𝐄(x)`$, be encoded in an ordinary function (potential), $`\varphi (x)`$, whose gradient is $`𝐄`$? The well-known answer is that a necessary and sufficient condition is that $`\mathrm{curl}𝐄=0`$. The importance of this encoding does not have to be emphasized to physicists; entropy’s role is similar to the potential’s role and the existence and meaning of entropy are not based on any formula such as $`S=\mathrm{\Sigma }_ip_i\mathrm{ln}p_i`$, involving probabilities $`p_i`$ of “microstates”. Entropy is derived (uniquely, we hope) from the list of pairs $`XY`$; our aim is to figure out what properties of this list (analogous to the curl-free condition) will allow it to be described by an entropy. That entropy will then be endowed with an unambiguous physical meaning independent of anyone’s assumptions about “the arrow of time”, “coarse graining” and so on. Only the list, which is given by physics, is important for us now. The required properties of $``$ do not involve concepts like “heat” or “reversible engines”, not even “hot” and “cold” are needed. Besides the “obvious” conditions “$`XX`$ for all $`X`$” (reflexivity) and “$`XY`$ and $`YZ`$ implies $`XZ`$” (transitivity) one needs to know that the relation behaves reasonably with respect to the composition and scaling of states. By this we mean the following: * Adiabatic accessibility is consistent with the composition of states: $`XY`$ and $`ZW`$ implies $`(X,Z)(Y,W)`$. * Scaling of states does not affect adiabatic accessibility: If $`XY`$, then $`\lambda X\lambda Y`$. * Systems can be cut adiabatically into two parts: If $`0<\lambda <1`$, then $`X((1\lambda )X,\lambda X)`$, and the recombination of the parts is also adiabatic: $`((1\lambda )X,\lambda X)X`$. * Adiabatic accessibility is stable with respect to small perturbations: If $`(X,\epsilon Z)(Y,\epsilon W)`$ for arbitrarily small $`\epsilon >0`$, then $`XY`$. These requirements are all very natural. In fact, in traditional approaches they are usually taken for granted, without mention. They are not quite sufficient, however, to define entropy. A crucial additional ingredient is the comparison hypothesis for the relation $``$. In essence, this is the hypothesis that all equilibrium states, simple or compound, can be grouped into classes, such that if $`X`$ and $`Y`$ are in the same class, then either $`XY`$ or $`YX`$. In nature, a class consists of all states with the same mass and chemical composition—that is, with the same amount of each of the chemical elements. If chemical reactions and mixing processes are excluded, the classes are smaller and may be identified with the “systems” in the usual parlance. But it should be noted that systems can be compound, or consist of two or more vessels of different substances. In any case, the role of the comparison hypothesis is to insure that the list of pairs $`XY`$ is sufficiently long. Indeed, we shall give an example later where the list of pairs satisfies all the other axioms, but which is not describable by an entropy function. ## The construction of entropy Our main conclusion (which we do not claim isobvious, but whose proof can be found in reference ) is that the existence and uniqueness of entropy is a consequence of the comparison hypothesis and the assumptions about adiabatic accessibility stated above. In fact, if $`X_0`$, $`X`$ and $`X_1`$ are three states of a system and $`\lambda `$ is any scaling factor between 0 and 1, then either $`X((1\lambda )X_0,\lambda X_1)`$ or $`((1\lambda )X_0,\lambda X_1)X`$ must be true, by the comparison hypothesis. If both alternatives hold, then the properties of entropy demand that $$S(X)=(1\lambda )S(X_0)+\lambda S(X_1).$$ If $`S(X_0)S(X_1)`$ this equality can hold for at most one $`\lambda `$. With $`X_0`$ and $`X_1`$ as reference states, the entropy is therefore fixed, apart from two free constants, namely the values $`S(X_0)`$ and $`S(X_1)`$. From the properties of the relation $``$ listed above, one can show that there is, indeed, always a $`0\lambda 1`$ with the required properties, provided that $`X_0XX_1`$. It is the largest $`\lambda `$, denoted $`\lambda _{\mathrm{max}}`$, such that $`((1\lambda )X_0,\lambda X_1)X`$. Defining the entropies of the reference states arbitrarily as $`S(X_0)=0`$ and $`S(X_1)=1`$ unit, we obtain the following simple formula for entropy: $$S(X)=\lambda _{\mathrm{max}}\text{units}.$$ The scaling factors $`(1\lambda )`$ and $`\lambda `$ measure the amount of substance in the states $`X_0`$ and $`X_1`$ respectively. The formula for entropy can therefore be stated in the following words: $`S(X)`$ is the maximal fraction of substance in the state $`X_1`$ that can be transformed adiabatically (that is, in the sense of $``$) into the state $`X`$ with the aid of a complementary fraction of substance in the state $`X_0`$. This way of measuring $`S`$ in terms of substance is reminiscent of an old idea, suggested by Pierre Laplace and Antoine Lavoisier, that heat be measured in terms of the amount of ice melted in a process. As a concrete example, let us assume that $`X`$ is a state of liquid water, $`X_0`$ of ice and $`X_1`$ of vapor. Then $`S(X)`$ for a kilogram of liquid, measured with the entropy of a kilogram of water vapor as a unit, is the maximal fraction of a kilogram of vapor that can be transformed adiabatically into liquid in state $`X`$ with the aid of a complementary fraction of a kilogram of ice. In this example the maximal fraction $`\lambda _{\mathrm{max}}`$ cannot be achieved by simply exposing the ice to the vapor, causing the former to melt and the latter to condense. This would be an irreversible process—that is, it would not be possible to reproduce the initial amounts of vapor of ice adiabatically (in the sense of the definition given earlier) from the liquid. By contrast, $`\lambda _{\mathrm{max}}`$ is uniquely determined by the requirement that one can pass adiabatically from $`X`$ to $`((1\lambda _{\mathrm{max}})X_0,\lambda _{\mathrm{max}}X_1)`$ and vice versa. For this transformation it is necessary to extract or add energy in the form of work—for example by running a little reversible Carnot machine that transfers energy between the high-temperature and low-temperature parts of the system. We stress, however, that neither the concept of a “reversible Carnot machine” nor that of “temperature” is needed for the logic behind the formula for entropy given above. We mention these concepts only to relate our definition of entropy to concepts for which the reader may have an intuitive feeling. By interchanging the roles of the three states, the definition of entropy is easily extended to situations where $`XX_0`$ or $`X_1X`$. Moreover, the reference points $`X_0`$ and $`X_1`$, where the entropy is defined to be 0 and 1 unit respectively, can be picked consistently for different systems such that the entropy will satisfy the crucial additivity and extensivity conditions $$S(X,X^{})=S(X)+S(X^{})\mathrm{and}S(\lambda X)=\lambda S(X).$$ It is important to understand that once the existence and uniqueness of entropy has been established one need not rely on the $`\lambda _{\mathrm{max}}`$ formula displayed above to determine it in practice. There are various experimental means to determine entropy that are usually much more practical. The standard method consists of measuring pressures, volumes and temperatures (on some empirical scale), as well as specific and latent heats. The empirical temperatures are converted into absolute temperatures $`T`$ (by means of formulas that follow from the mere existence of entropy but do not involve $`S`$ directly), and the entropy is computed by means of formulas like $`\mathrm{\Delta }S=(dU+PdV)/T`$, with $`P`$ the pressure. The existence and uniqueness of entropy implies that this formula is independent of the path of integration. ## Comparability of states The possibility of defining entropy entirely in terms of the relation $``$ was first clearly stated by Giles . (Giles’s definition is different from ours, albeit similar in spirit.) The importance of the comparison hypothesis had been realized earlier, however . All these authors take the comparison hypothesis as a postulate—that is, they do not attempt to justify it from other simpler premises. However, it is in fact possible to derive comparability for any pair of states of the same system from some natural and directly accessible properties of the relation $``$ . In this derivation of comparison the customary parametrization of states in terms of energy and work coordinates is used, but it has to be stressed that such parametrizations are irrelevant, and therefore not used, for our definition of entropy—once the comparison hypothesis is established. To appreciate the significance of the comparison hypothesis it may be helpful to consider the following example. Imagine a world whose thermodynamical sytems consist exclusively of incompressible solid bodies. Moreover, all adiabatic state changes in this world are supposed to be obtained by means of the following elementary operations: * Mechanical rubbing of the individual systems, increasing their energy. * Thermal equilibration in the conventional sense (by bringing the systems into contact.) The state space of the compound system consisting of two identical bodies, 1 and 2, can be paramertized by their energies, $`U_1`$ and $`U_2`$. If $`X=(U_1,U_2)`$ and $`Y==(U_1^{},U_2^{})`$ are such that $`U_1^{}<U_1<U_2<U_2^{}`$ and $`U_1+U_2<U_1^{}+U_2^{}`$ then one finds that that neither $`XY`$ nor $`YX`$ holds. The comparison hypothesis is therefore violated in this hypothetical example, and it is not possible to characterize adiabatic accessibility by means of an additive entropy function. A major part of our work consists of understanding why such situations do not happen—why the comparison hypothesis appears to be true in the real world. The derivation of the comparison hypothesis is based on an analysis of simple systems, which are the building blocks of thermodynamics. As already mentioned the states of such systems are described by one energy coordinate $`U`$ and at least one work coordinate, like the volume $`V`$. The following concepts play a key role in this analysis: * The possibility of forming “convex combinations” of states of simple systems with respect to the energy $`U`$ and volume $`V`$ (or other work coordinates). This means that given any two states $`X`$ and $`Z`$ of one kilogram of our system one can pick any state $`Y`$ on the line between them in $`U`$, $`V`$ space and, by taking appropriate fractions $`\lambda `$ and $`1\lambda `$ in states $`X`$ and $`Z`$, respectively, there will be an adiabatic process taking this pair of states into the state $`Y`$. This process is usually quite elementary. For example, for gases and liquids one need only remove the barrier that separates the two fractions of the system. The fundamental property of entropy increase will then tell us that $`S(Y)\lambda S(X)+(1\lambda )S(Z)`$. As Gibbs emphasized, this “concavity” is the basis for thermodynamical stability—namely positivity of specific heats and compressibilities. * The existence of at least one irreversible adiabatic state change, starting from any given state. In conjuction with concavity of $`S`$ this seemingly weak requirement excludes the possibility that the entropy is constant in a whole neighborhood of some state. The classical formulations of the second law follow from this. * The concept of thermal equilibrium between simple systems, which means, operationally, that no state changes takes place when the systems are allowed to exchange energy with each other at fixed work coordinates. The zeroth law of thermodynamic says that if two systems are in thermal equilibrium with a third, then they are in thermal equilibrium with one another. This property is essential for the additivity of entropy, because it allows a consistent adjustment of the entropy unit for different systems. It leads to a definition of temperature by the usual formula $`1/T=(S/U)_V`$. Using these notions (and a few others of a more technical nature) the comparison hypothesis can be established for all simple systems and their compounds. It is more difficult to justify the comparability of states if mixing processes or chemical reactions are taken into account. In fact, although a mixture of whiskey and water at $`0^{}`$ C is obviously adiabatically accessible from separate whiskey and ice by pouring whiskey from one glass onto the rocks in the other glass, it is not possible to reverse this process adiabatically. Hence it is not clear that a block of a frozen whiskey/water mixture at $`10^{}`$ C, say, is at all related in the sense of $``$ to a state in which whiskey and water are in separate glasses. Textbooks usually appeal here to gedanken experiments with “semipermeable membrane” that let only water molecules through and withhold the whiskey molecules, but such membranes really exist only in the mind . However, without invoking any such device, it turns out to be possible to shift the entropy scales of the various substances in such a way that $`XY`$ always implies $`S(X)S(Y)`$. The converse assertion, namely, $`S(X)S(Y)`$ implies $`XY`$ provided $`X`$ and $`Y`$ have the same chemical composition, cannot be guaranteed a priori for mixing and chemical reactions, but it is empirically testable and appears to be true in the real world. This aspect of the second law, comparability, is not usually stressed, but it is important; it is challenging to figure out how to turn the frozen whiskey/water block into a glass of whiskey and a glass of water without otherwise changing the universe, except for moving a weight, but such an adiabatic process is possible. ## What has been gained? The line of thought that started more than forty years ago has led to an axiomatic foundation for thermodynamics. It is appropriate to ask what if anything has been gained compared to the usual approaches involving quasi-static processes and Carnot machines on the one hand and statistical mechanics on the other hand. There are several points. One is the elimination of intuitive, but hard-to-define concepts like “hot”, “cold” and “heat”. Another is the recognition of entropy as a codification of possible state changes, $`XY`$, that can be accomplished without changing the rest of the universe in any way except for moving a weight. Temperature is eliminated as an a priori concept and appears in its natural place as a quantity derived from entropy and whose consistent definition really depends on the existence of entropy, rather than the other way around. To define enetropy, there is no need for special machines and processes on the empirical side, and there is no need for assumptions about models on the statistical mechanical side. Just as energy conservation was eventually seen to be a consequence of time translation invariance, in like manner entropy can be seen to be a consequence of some simple properties of the list of state pairs related by adiabatic accessibility. If the second law can be demystified, so much the better. If it can be seen to be a consequence of simple, plausible notions then, as Einstein said, it cannot be overthrown. ## Acknowledgements We are grateful to Shivaji Sondhi and Roderich Moessner for helpful suggestions. Lieb’s work was supported by NSF grant PHY 9820650. Yngvason’s work was supported by the Adalsteinn Kristjánsson Foundation and the University of Iceland. ## References C. Kittel and H. Kroemer, Thermal Physics, p. 57, Freeman, NY (1980). E.H. Lieb and J. Yngvason, Physics Reports 310, 1 (1999). A. Einstein, Autobiographical Notes in Albert Einstein: Philosopher-Scientist P. A. Schilpp (ed.), Library of Living Philosophers, vol VII, p. 33, Cambridge University Press, London, 1970. P.T. Landsberg, Rev. Mod. Phys. 28, 363 (1956). H.A. Buchdahl, The Concepts of Classical Thermodynamics, Cambridge University Press, London (1966). G. Falk and H. Jung Handbuch der Physik, III/2, S. Flügge ed., p. 199 Springer, Berlin (1959). R.Giles, Mathematical Foundations of Thermodynamics, Pergamon, Oxford (1964). D.R.Owen, A First Course in the Mathematical Foundations of Thermodynamics, Springer, Heidelberg (1984). J. Serrin, Arch. Rat. mech. Anal. 70, 355 (1979). M. Silhavý, The Mechanics and Thermodynamics of Continuous Media, Springer, Heidelberg (1997). C. A. truesdell, S. Baharata, The Concepts and Logic of Thermodynamics as a Theory of heat Engines, Heidelberg, (1977). J.L. Lebowitz, I. Prigogine, and D.Ruelle, Physica A 263, 516, 528, 540 (1999). E. Fermi, Thermodynamics, Dover, NY, (1956), page 101.
warning/0003/hep-ph0003315.html	ar5iv	text	# Bounds on “charginos nearly degenerate with the lightest neutralino” mass from precision measurements ## 1 Introduction According to the latest searches performed at LEP II at center-of-mass energies up to $`189\text{ GeV}`$, the present bounds on chargino mass are $`m_{\stackrel{~}{\chi }^\pm }90\text{ GeV}`$ for the higgsino-dominated case (or when the sneutrino is heavy) and $`m_{\stackrel{~}{\chi }^\pm }80\text{ GeV}`$ in the wino-dominated light-sneutrino scenario. However, when the lightest chargino and neutralino (the latter being the LSP) are almost degenerate in mass, the charged decay products of the light chargino are very soft, and the above quoted bounds are no longer valid. A special search for such light charginos has been performed recently by the DELPHI collaboration, and the case of $`\mathrm{\Delta }M^\pm m_{\stackrel{~}{\chi }_1^\pm }m_{\stackrel{~}{\chi }_1^0}100\text{ MeV}`$ is now excluded. In the region $`\mathrm{\Delta }M^\pm 1\text{ GeV}`$ the analysis of the Initial State Radiation (ISR) can be used to put a limit on the chargino mass, but in the case of wino domination with a light sneutrino this technique fails and charginos as light as $`45\text{ GeV}`$ (this bound coming from the measurements of $`Z`$ decays at LEP I and SLC) are still allowed. The case of almost degenerate chargino and neutralino can be naturally realized in SUSY and the possibilities to find such particles are discussed in literature. In this talk we investigate the radiative corrections to the electroweak precision measurements generated by such almost degenerate particles. When their masses are close to $`m_Z/2`$ one-loop contributions are large and they spoil the perfect description of experimental data by the Standard Model. Due to the decoupling property of SUSY models, when $`m_{\stackrel{~}{\chi }^{\pm ,0}}m_Z`$ the radiative corrections are power suppressed. ## 2 Discussion The contributions of new physics to the electroweak precision data through oblique corrections can be conveniently parameterized in terms of the three functions $`V_m`$, $`V_A`$ and $`V_R`$ In the simplest supersymmetric extensions of the Standard Model the chargino-neutralino sector is defined by the numerical values of the four parameters $`M_1`$, $`M_2`$, $`\mu `$ and $`\mathrm{tan}\beta `$, and the case of nearly degenerate lightest chargino and neutralino naturally arise when: * $`M_2\mu `$: in this limit the particles of interest form an $`SU(2)`$ doublet of Dirac fermions, whose wave functions are dominated by higgsinos, and their contribution to the $`V_i`$ functions is: $`\begin{array}{cc}\hfill \delta ^{\stackrel{~}{h}}V_m& ={\displaystyle \frac{16}{9}}[({\displaystyle \frac{1}{2}}s^2+s^4)(1+2\chi )F(\chi )\hfill \\ & ({\displaystyle \frac{1}{2}}s^2)(1+2{\displaystyle \frac{\chi }{c^2}})F\left({\displaystyle \frac{\chi }{c^2}}\right){\displaystyle \frac{s^4}{3}}],\hfill \end{array}`$ (1) $`\delta ^{\stackrel{~}{h}}V_A`$ $`={\displaystyle \frac{16}{9}}\left({\displaystyle \frac{1}{2}}s^2+s^4\right)\left[{\displaystyle \frac{12\chi ^2F(\chi )2\chi 1}{4\chi 1}}\right],`$ (2) $`\delta ^{\stackrel{~}{h}}V_R`$ $`={\displaystyle \frac{16}{9}}c^2s^2\left[\left(1+2\chi \right)F(\chi ){\displaystyle \frac{1}{3}}\right],`$ (3) where $`\chi (m_{\stackrel{~}{\chi }^{\pm ,0}}/m_Z)^2`$, the function $`F`$ is defined in App. B of Ref. , and $`s^2`$ ($`c^2`$) is the sine (cosine) squared of the electroweak mixing angle; * $`\mu M_2`$: in this case we get an $`SU(2)`$ triplet of Majorana fermions, with the wave functions dominated by winos, and the expressions for the corrections to $`V_i`$ are: $`\begin{array}{cc}\hfill \delta ^{\stackrel{~}{w}}V_m& ={\displaystyle \frac{16}{9}}[c^4(1+2\chi )F(\chi )\hfill \\ & (12s^2)(1+2{\displaystyle \frac{\chi }{c^2}})F\left({\displaystyle \frac{\chi }{c^2}}\right){\displaystyle \frac{s^4}{3}}],\hfill \end{array}`$ (4) $`\delta ^{\stackrel{~}{w}}V_A`$ $`={\displaystyle \frac{16}{9}}c^4\left[{\displaystyle \frac{12\chi ^2F(\chi )2\chi 1}{4\chi 1}}\right],`$ (5) $`\delta ^{\stackrel{~}{w}}V_R`$ $`={\displaystyle \frac{16}{9}}c^2s^2\left[\left(1+2\chi \right)F(\chi ){\displaystyle \frac{1}{3}}\right].`$ (6) The $`\chi ^2`$ for the new physics contributions to $`V_i`$ was computed using the computer program LEPTOP and the corresponding Confidence Level (together with the numerical value of the $`\delta V_i`$ functions) is plotted in Fig. 1 against the chargino-neutralino mass $`m_{\stackrel{~}{\chi }^{\pm ,0}}`$. We see that at $`95\%`$ C.L. the bounds $`m_{\stackrel{~}{\chi }^{\pm ,0}}51\text{ GeV}`$ (higgsino-dominated case) and $`m_{\stackrel{~}{\chi }^{\pm ,0}}56\text{ GeV}`$ (wino-dominated case) should be satisfied. Note that the main contribution to $`\chi ^2`$ comes from $`\delta V_A`$, which is singular at $`m_{\stackrel{~}{\chi }^{\pm ,0}}=m_Z/2`$. This singularity is not physical and our formulas are valid only for $`2m_{\stackrel{~}{\chi }^{\pm ,0}}m_Z+\mathrm{\Gamma }_Z`$; however, the existence of $`\chi ^\pm `$ with a mass closer to $`m_Z/2`$ will change Z-boson Breit-Wigner curve, therefore it is also not allowed. ## 3 Conclusions Let us briefly discuss the contributions of other SUSY particles to the $`V_i`$ functions. In the considered limits the remaining charginos and neutralinos are very heavy, so they simply decouple. The contributions of the three generations of sleptons (with masses larger than $`90\text{ GeV}`$) and of the first two generations of squarks (with masses larger than $`200\text{ GeV}`$ to satisfy Tevatron bounds) into $`V_A`$ are smaller than $`0.1`$, so they can safely be neglected. Concerning the contributions of the third generation squarks, although enhanced by the large top-bottom mass difference they are almost universal, so compensating the negative value of $`V_A`$ will generate large positive terms into $`V_R`$ and $`V_m`$ and the overall $`\chi ^2`$ will not be better. Finally, according to Ref. the contribution to radiative corrections of the whole MSSM Higgs sector equals with very good accuracy that of a single SM higgs with the same mass as the lightest MSSM neutral higgs, so it is already accounted for in our analysis. Let us remark that in the case of wino domination with a light sneutrino, which occurs naturally in anomaly-mediated SUSY breaking scenarios, the analysis of the ISR fails and the bound $`m_{\stackrel{~}{\chi }^{\pm ,0}}56\text{ GeV}`$ from precision measurements is presently the strongest constraint which can be imposed on the chargino-neutralino mass. ## Acknowledgments I wish to thank my collaborator M.I. Vysotsky. I am also grateful to A.N. Rozanov for evaluating the C.L. shown in Fig. 1 with the program LEPTOP. This work was supported by DGICYT under grant PB98-0693 and by the TMR network grant ERBFMRX-CT96-0090 of the European Union.
warning/0003/hep-th0003169.html	ar5iv	text	# 1 Introduction ## 1 Introduction The M5-brane plays an important role in studying properties of M-theory , the theory of strings and associated field theories. For instance, many physically important multibrane configurations, realized to be relevant to a brane description of non-Abelian gauge theories and a brane-world scenario , can be considered as a specific compactification of a single D=11 M5-brane down to lower dimensions (with or without its subsequent T-dualization). A direct dimensional reduction of D=11 space-time with an M5-brane down to ten-dimensional space-time produces a so-called NS5-brane of type IIA supergravity which has been intensively studied in relation to six-dimensional gauge theories and “little string theories” <sup>2</sup><sup>2</sup>2The double dimensional reduction of the M5-brane action is well known to result in a type IIA D=10 Dirichlet 4-brane. It has also been shown how by reducing the M5-brane action one may arrive at a duality-symmetric D3-brane action.. The verification of the quantum consistency of M-theory requires, in particular, finding a mechanism of anomaly cancellation in the presence of M5-branes. It has been shown that the anomaly problem has a natural solution in the case of $`D=10`$ NS5-branes , while in the case of the D=11 M5-brane the situation is much more subtle and requires additional study . Mechanisms for the M5-brane anomaly cancellation proposed recently in include as an important feature the reduction of the structure group $`SO(5)`$ of the normal bundle of the M5-brane down to its $`SO(4)`$ subgroup. Such a reduction implies an existence of a covariantly constant vector field and, therefore, looks very much as a dimensional reduction (to be more precise, the dimensional reduction is a particular case of such an ‘M5-brane framing’ ) <sup>3</sup><sup>3</sup>3In contrast to the analysis of ref. is based on the assumption that a full understanding of anomaly cancellation requires keeping the full $`SO(5)`$. We are thankful to Jeff Harvey for clarifying this difference in the approaches.. These facts provide us with a motivation to study in more detail the dynamical and symmetry properties of the NSIIA five-brane by constructing a full worldvolume action describing its dynamics in a type IIA D=10 supergravity background. By now the action for the NS5–brane has been constructed up to a second order in the field strength of a two–rank self–dual worldvolume gauge field of the five–brane and only in a background of the bosonic sector of IIA D=10 supergravity . The aim of this paper is to get a full, nonlinear and $`\kappa `$–symmetric, NS5–brane action in a curved IIA D=10 target superspace by carrying out the direct dimensional reduction of the D=11 M5–brane action , and thus filling in a gap in the list of worldvolume actions for supersymmetric extended objects found in string theory. The fact that the NS5–brane can be regarded as an M5–brane propagating in a dimensionally reduced D=11 supergravity background substantially simplifies the analysis of the NS5–brane model, in particular, allowing one to derive its symmetries and dynamical properties directly from those of the M5–brane. For instance, the physical field content of the IIA D=10 NS5–brane is the same as of the M5–brane. The bosonic sector consists of three degrees of freedom corresponding to the two–rank self–dual worldvolume field and five worldvolume scalars. In the case of the M5–brane the five scalar fields describe its oscillations in a D=11 background in the directions transversal to the M5–brane worldvolume, while in the case of the NS5–brane four scalar fields correspond to transversal oscillations in a D=10 background, and the fifth scalar field (corresponding to the compactified dimension of the D=11 space) ‘decouples’ and becomes a ‘purely’ worldvolume field. This results in the abovementioned reduction of the M5-brane normal bundle structure group $`SO(5)`$ down to $`SO(4)`$. For both five–branes eight fermionic fields can be associated with brane ‘oscillations’ in Grassmann directions of corresponding target superspaces. To get the action describing the dynamics of the physical modes of the NS5–brane as a dimensionally reduced M5–brane action we first briefly remind the structure and properties of the latter. In Sections 2–5 we consider bosonic M5– and NS5–branes and in Section 6 we describe the full target–superspace covariant and $`\kappa `$–invariant NS5–brane action. ## 2 The M5-brane action In the absence of interactions with antisymmetric tensor fields of $`D=11`$ supergravity the action for the bosonic sector of the M5-brane has the following form $$S=d^6\xi \left[\sqrt{det(\widehat{g}_{mn}+i\widehat{H}_{mn}^{})}+\frac{\sqrt{\widehat{g}}}{4\sqrt{\widehat{aa}}}\widehat{H}^{mn}H_{mnr}^ra\right]$$ (1) where $$m,n,\mathrm{}=0,\mathrm{},5;$$ are vector indices of $`d=6`$ worldvolume coordinates $`\xi ^m`$, $$\underset{¯}{\overset{^}{m}},\underset{¯}{\overset{^}{n}},\mathrm{}=0,\mathrm{},10,$$ are vector indices of $`D=11`$ target space coordinates $`\widehat{X}^{\underset{¯}{\overset{^}{m}}}`$ $$\widehat{g}_{mn}=_m\widehat{X}^{\underset{¯}{\overset{^}{m}}}\widehat{g}_{\underset{¯}{\widehat{m}\widehat{n}}}^{(11)}_n\widehat{X}^{\underset{¯}{\overset{^}{n}}},$$ (2) is the worldvolume metric induced by embedding the five–brane into a $`D=11`$ gravity background with a metric $`\widehat{g}_{\underset{¯}{\widehat{m}\widehat{n}}}^{(11)}(X)`$ (we use the ‘almost minus’ Minkowski signature $`(+\mathrm{})`$), $`H_{mnl}(\xi )=3_{[m}b_{nl]}`$ is the field strength of the worldvolume antisymmetric tensor field $`b_{mn}(\xi )`$, $$\widehat{H}_{mn}^{}\frac{1}{\sqrt{\widehat{aa}}}\widehat{H}_{mnr}^{}^ra,\widehat{H}^{mnl}=\frac{1}{3!\sqrt{\widehat{g}}}ϵ^{mnlrsq}H_{rsq},$$ (3) $`a(\xi )`$ is an auxiliary scalar field ensuring the covariance of the model, and $$\widehat{aa}_ma\widehat{g}^{mn}_na$$ (4) denotes the scalar product of the $`d=6`$ vector $`_ma`$ with respect to the metric (2). In what follows the ‘hat’ over quantities indicates that they correspond to or induced by the eleven–dimensional theory. In addition to the usual gauge symmetry of the $`b_2`$ field $$\delta a(\xi )=0,\delta b_{mn}=2_{[m}\phi _{n]}(\xi ),$$ (5) the action (1) is invariant under the following transformations , , $$\delta a(\xi )=0,\delta b_{mn}=2\varphi _{[m}(\xi )_{n]}a(\xi ),$$ (6) $$\delta a=\phi (\xi ),\delta b_{mn}=\frac{\delta a}{\sqrt{\widehat{aa}}}[\widehat{^{}}_{mn}H_{mnp}\widehat{g}^{ps}\frac{_sa}{\sqrt{\widehat{aa}}}],$$ (7) where $$\widehat{^{}}_{mn}=\frac{2}{\sqrt{\widehat{g}}}\frac{\delta _{DBI}}{\delta \widehat{H}^{mn}},_{DBI}\sqrt{det(\widehat{g}_{mn}+i\widehat{H}_{mn}^{})}.$$ (8) Note that at the linearized level, $`\widehat{^{}}_{mn}`$ defined in (8) reduces to $`\widehat{H}_{mn}^{}`$. The symmetries (6) and (7) are characteristic of the covariant approach to the Lagrangian description of duality–symmetric fields. They ensure the $`b_2`$ field equation of motion to reduce to a self–duality condition, as well as the connection with non-covariant formulations Let us briefly describe how one derives the symmetries (6) and (7) and gets the self–duality condition , . To this end note that the second term in the action (1) can be written in terms of differential forms $$d^6\xi _1d^6\xi \sqrt{\widehat{g}}\frac{1}{4\sqrt{\widehat{aa}}}\widehat{H}^{mn}H_{mnr}^ra=_^6\frac{1}{2}\widehat{v}H_3i_{\widehat{v}}H_3,$$ (9) where <sup>4</sup><sup>4</sup>4In our notation $`d\xi ^{m_1}\mathrm{}d\xi ^{m_6}=d^6\xi ϵ^{m_1\mathrm{}m_6}`$. $$\widehat{v}=d\xi ^m\widehat{v}_m,\widehat{v}_k\frac{_ka}{\sqrt{\widehat{aa}}}$$ (10) $$H_3\frac{1}{3!}d\xi ^md\xi ^nd\xi ^lH_{lnm},i_{\widehat{v}}H_3\frac{1}{2}d\xi ^md\xi ^n\widehat{v}_k\widehat{g}^{kl}H_{lnm}.$$ (11) The variation of the first term in (1) with respect to the gauge field and the scalar $`a(\xi )`$ can be written in terms of differential forms as $$d\xi ^6\delta _{DBI}d\xi ^6\delta \sqrt{det(\widehat{g}_{mn}+i\widehat{H}_{mn}^{})}=_^6\widehat{}_2^{}\delta \widehat{H}_2^{},$$ (12) where 2-forms $`\widehat{}_2^{}`$ and $`\widehat{H}_2^{}`$ are constructed respectively from the tensors (8) and (3) $$\widehat{}_2^{}\frac{1}{2}d\xi ^md\xi ^n\widehat{^{}}_{nm},\widehat{H}_2^{}=i_{\widehat{v}}H_3\frac{1}{2}d\xi ^md\xi ^n\widehat{H}_{nm}^{}$$ (13) and $``$ is the Hodge operation in $`d=6`$ dimensions <sup>5</sup><sup>5</sup>5To have $`=I`$ we define $$\mathrm{\Omega }_2=\frac{1}{2!4!}d\xi ^{m_4}\mathrm{}d\xi ^{m_1}\sqrt{g}ϵ_{m_1\mathrm{}m_4n_1n_2}\mathrm{\Omega }^{n_1n_2},$$ $$\mathrm{\Omega }_4=+\frac{1}{2!4!}d\xi ^{m_2}d\xi ^{m_1}\sqrt{g}ϵ_{m_1m_2n_1\mathrm{}n_4}\mathrm{\Omega }^{n_1\mathrm{}n_4}\frac{1}{2!4!}d\xi _{m_2}d\xi _{m_1}\frac{1}{\sqrt{g}}ϵ^{m_1m_2n_1\mathrm{}n_4}\mathrm{\Omega }_{n_1\mathrm{}n_4}$$ . Using the identities $$i_{\widehat{v}}\delta \widehat{v}=0,\mathrm{\Omega }_6i_{\widehat{v}}\mathrm{\Omega }_6\widehat{v},i_{\widehat{v}}H_3=H_3\widehat{v},$$ (14) $$\widehat{v}H_3i_{\widehat{v}}\delta H_3=\widehat{v}i_{\widehat{v}}H_3\delta H_3+H_3\delta H_3,$$ (15) $$\widehat{v}H_3i_{\delta \widehat{v}}H_3=\delta \widehat{v}H_3i_{\widehat{v}}H_3=\delta \widehat{v}\widehat{v}i_{\widehat{v}}H_3i_{\widehat{v}}H_3$$ (16) and $$\widehat{v}\widehat{}_2^{}\widehat{}_2^{}=\widehat{v}\widehat{H}_2^{}\widehat{H}_2^{}ϵ^{abcdef}\widehat{^{}}_{bc}\widehat{^{}}_{de}\widehat{v}_f=ϵ^{abcdef}\widehat{H}_{bc}^{}\widehat{H}_{de}^{}\widehat{v}_f,$$ (17) one can rewrite the variation of the Lagrangian (1) in the form $$𝑑\xi ^6\delta _^6(\frac{1}{2}H_3\delta H_3da_2\delta H_3\frac{1}{2}d\delta ada_2_2),$$ (18) where $$_2\frac{1}{\sqrt{\widehat{aa}}}(\widehat{^{}}i_{\widehat{v}}H_3)=\frac{1}{2}d\xi ^md\xi ^n_{nm},$$ (19) or $$_{mn}\frac{1}{\sqrt{\widehat{aa}}}(\widehat{}_{mn}^{}H_{mnl}\widehat{g}^{lk}\frac{_ka}{\sqrt{\widehat{aa}}}).$$ (20) Since $`H_3=db_2`$, the variation (18) can be written (up to a total derivative) in the following form <sup>6</sup><sup>6</sup>6We use conventions where external derivative acts from the right: $$d\mathrm{\Omega }_q=\frac{1}{q!}d\xi ^{m_q}\mathrm{}d\xi ^{m_1}d\xi ^n_n\mathrm{\Omega }_{m_1\mathrm{}m_q},d(\mathrm{\Omega }_p\mathrm{\Omega }_q)=\mathrm{\Omega }_pd\mathrm{\Omega }_q+()^qd\mathrm{\Omega }_p\mathrm{\Omega }_q.$$ $$𝑑\xi ^6\delta _^6d(da_2)(\delta b_2\delta a_2),$$ (21) from which the invariance of the action under (6) and (7) becomes evident. From (21) it also follows that the equation of motion of $`b_2`$ field is $$d(da_2)=0,$$ (22) and the equation of motion of $`a(x)`$ is a consequence of eq. (22). It can be shown that, using the symmetry (6), the second–order equation (22) reduces to the first–order self–duality condition $$_2\frac{1}{\sqrt{\widehat{aa}}}(\widehat{^{}}i_{\widehat{v}}H_3)=0,$$ (23) or in components $$\widehat{}_{mn}^{}=H_{mnl}\widehat{g}^{lk}\frac{_ka}{\sqrt{\widehat{aa}}}.$$ (24) To prove this note that $`\widehat{H}_2^{}`$ is invariant under the transformations (6) $$\delta b_2=da\varphi _1\sqrt{\widehat{aa}}\widehat{v}\varphi _1$$ $$\delta \widehat{H}_2^{}\delta i_{\widehat{v}}(\widehat{H}_3)=\sqrt{\widehat{aa}}i_{\widehat{v}}((\widehat{v}d\varphi _1))0.$$ Hence, the transformations of the two-form (19) reduce to $$\delta _2=\frac{1}{\sqrt{\widehat{aa}}}\delta i_{\widehat{v}}H_3=i_{\widehat{v}}(\widehat{v}d\varphi _1).$$ (25) Eq. (25) is simplified when one takes into account that $`i_{\widehat{v}}da=\sqrt{\widehat{aa}}i_{\widehat{v}}\widehat{v}=\sqrt{\widehat{aa}}.`$ Then $$\delta _2=d\varphi _1+i_{\widehat{v}}d\varphi _1\widehat{v},$$ and $$\delta (da_2)=dad\varphi _1.$$ (26) We now observe that eq. (26) is similar to the general solution of eq. (22) for $`da_2`$. This means that the general solution of eq. (22) can be gauged to zero with the use of the symmetry (6), and eq. (23) appears just as a result of such gauge fixing. Remember that $`\widehat{}_{mn}^{}`$ is defined in (8) and reduces to $`\widehat{H}_{mn}^{}=H_{mnl}^{}\widehat{g}^{lk}\frac{_ka}{\sqrt{\widehat{aa}}}`$ at the linearized level, the equation (24) becoming the conventional self–duality condition $`\widehat{H}_{lmn}^{}=H_{lmn}`$. Further details on the classical dynamics of the M5–brane the reader may find in , . ## 3 Dimensional reduction of $`D=11`$ gravity and the NS5-brane action The procedure of the direct dimensional reduction assumes a compactification of some of target–space spatial dimensions (one in our case), the worldvolume of the $`p`$–brane being not compactified. A standard (string frame) ansatz for the target–space vielbein under the Kaluza-Klein reduction of one spatial dimension has the following form $$E^{\underset{¯}{\overset{^}{a}}}=(E^{\underset{¯}{a}},E^{10})d\widehat{X}^{\underset{¯}{\overset{^}{m}}}E_{\underset{¯}{\overset{^}{m}}}^{}{}_{}{}^{\underset{¯}{\overset{^}{a}}}(\widehat{X}),\widehat{X}^{\underset{¯}{\overset{^}{m}}}=(X^{\underset{¯}{m}},y),y=\widehat{X}^{10},$$ $$E^{\underset{¯}{a}}=e^{\frac{1}{3}\mathrm{\Phi }}dX^{\underset{¯}{m}}e_{\underset{¯}{m}}^{\underset{¯}{a}}(X),E^{10}=e^{\frac{2}{3}\mathrm{\Phi }}(dydX^{\underset{¯}{m}}A_{\underset{¯}{m}})e^{\frac{2}{3}\mathrm{\Phi }},$$ $$e_{\underset{¯}{\overset{^}{m}}}^{}{}_{}{}^{\underset{¯}{\overset{^}{a}}}=\left(\begin{array}{cc}e^{\frac{1}{3}\mathrm{\Phi }}e_{\underset{¯}{m}}^{}{}_{}{}^{\underset{¯}{a}}& e^{\frac{2}{3}\mathrm{\Phi }}A_{\underset{¯}{m}}\\ 0& e^{\frac{2}{3}\mathrm{\Phi }}\end{array}\right),$$ (27) where $`y`$ is the coordinate compactified into a torus, and the reduction means that the background fields, such as components of (27), do not depend on $`y`$ which is now considered as an intrinsic scalar field in the 5–brane worldvolume. $`\mathrm{\Phi }(X)`$ is the dilaton field and $`A_{\underset{¯}{m}}(X)`$ is the Abelian vector gauge field of $`D=10`$ IIA supergravity. The $`U(1)`$–gauge transformations of $`A_{\underset{¯}{m}}(X)`$ and $`y`$ are $$\delta A_{\underset{¯}{m}}(X)=_{\underset{¯}{m}}\phi ^{(0)}(X),\delta y=\phi ^{(0)}(X).$$ (28) This ansatz leads to the following expression for the $`D=11`$ target space metric in terms of the $`D=10`$ metric $`g_{\underset{¯}{mn}}^{(10)}(X)=e_{\underset{¯}{m}}^{\underset{¯}{a}}e_{\underset{¯}{ma}}`$, $`A_{\underset{¯}{m}}(X)`$ and $`\mathrm{\Phi }(X)`$ $$\widehat{g}_{\underset{¯}{\widehat{m}\widehat{n}}}^{(11)}=\left(\begin{array}{cc}e^{\frac{2}{3}\mathrm{\Phi }}(g_{\underset{¯}{mn}}^{(10)}e^{2\mathrm{\Phi }}A_{\underset{¯}{m}}A_{\underset{¯}{n}})& e^{\frac{4}{3}\mathrm{\Phi }}A_{\underset{¯}{m}}\\ e^{\frac{4}{3}\mathrm{\Phi }}A_{\underset{¯}{n}}& e^{\frac{4}{3}\mathrm{\Phi }}\end{array}\right)$$ (29) and, consequently, to the following form of the six–dimensional induced metric (2) $$\widehat{g}_{mn}=e^{\frac{2}{3}\mathrm{\Phi }}(g_{mn}e^{2\mathrm{\Phi }}_m_n).$$ (30) In (30) $$g_{mn}=_mX^{\underset{¯}{m}}g_{\underset{¯}{mn}}^{(10)}(X)_nX^{\underset{¯}{n}},\underset{¯}{m}=0,\mathrm{},9$$ (31) is the six-dimensional metric induced by embedding the 5-brane worldvolume into the ten–dimensional curved space-time and $$_m=_myA_m,$$ (32) where $`A_m=_mX^{\underset{¯}{m}}A_{\underset{¯}{m}}(X)`$ is the worldvolume pullback of $`A_{\underset{¯}{m}}(X)`$ and $`_m`$ is the pullback of the one–form $``$ introduced in (27). $`_m`$ defined in (32) can be considered as a field strength of the worldvolume scalar field $`y(\xi )`$. It is invariant under the $`U(1)`$ gauge transformations (28). In what follows we will also use an expression for the inverse worldvolume metric $$\widehat{g}^{mn}=e^{\frac{2}{3}\mathrm{\Phi }}(g^{mn}+\frac{e^{2\mathrm{\Phi }}^m^n}{1e^{2\mathrm{\Phi }}^2}).$$ (33) The NS5–brane action follows from the M5-action (1) with the background metric having a particular form (29) and the coordinate $`\widehat{X}^{\underset{¯}{10}}=y(\xi )`$ being considered as an intrinsic worldvolume scalar field. To present the explicit form of the NS5–brane action we should rewrite all its constituents in terms of $`D=10`$ fields, and to ‘rescale’ worldvolume fields and their scalar products with respect to the worldvolume induced metric (31). For instance, the Hodge duality (3) is now redefined with respect to the metric (31) $$\widehat{H}^{mnp}=\sqrt{\frac{g}{\widehat{g}}}H^{mnp},H^{mnl}=\frac{1}{3!\sqrt{g}}ϵ^{mnlrsd}H_{rsd},$$ (34) and the M5–brane field strength $`\widehat{H}^{mn}`$ (3) is related to its NS5–brane counterpart $`H^{mn}`$ as $$\widehat{H}^{mn}=\sqrt{\frac{g}{\widehat{g}}}\sqrt{\frac{(a)^2}{\widehat{aa}}}H^{mn},H^{mn}=\frac{1}{3!\sqrt{g}}ϵ^{kmnpqr}H_{pqr}\frac{_ka}{\sqrt{(a)^2}},$$ (35) where the scalar product (4) has also been correspondingly redefined as $$\widehat{aa}_ka\widehat{g}^{ks}_sa=e^{\frac{2}{3}\mathrm{\Phi }}𝒩^2(a)^2,(a)^2_lag^{lm}_ma,$$ (36) with $`𝒩`$ standing for $$𝒩\sqrt{1+\frac{e^{2\mathrm{\Phi }}(a)^2}{(a)^2(1e^{2\mathrm{\Phi }}^2)}}=e^{\frac{1}{3}\mathrm{\Phi }}\sqrt{\frac{\widehat{aa}}{(a)^2}}.$$ (37) In view of eqs. (30), (35), (36) and (37) the antisymmetric tensor entering the DBI-like part of the M5–brane action is reexpressed in terms of $`H^{lm}`$ as follows $$\widehat{H}_{mn}^{}=\widehat{g}_{ml}\widehat{g}_{nk}\widehat{H}^{lk}=\widehat{g}_{ml}\widehat{g}_{nk}\sqrt{\frac{g}{\widehat{g}}}e^{\frac{1}{3}\mathrm{\Phi }}𝒩^1H^{lk},\widehat{g}_{ml}=e^{\frac{2}{3}\mathrm{\Phi }}(g_{ml}e^{2\mathrm{\Phi }}_m_l).$$ (38) As a result, substituting (30)–(38) into the action (1), we get the action for a bosonic 5–brane coupled to the metric, the dilaton and the gauge vector field of type IIA $`D=10`$ supergravity $$S=d^6\xi e^{2\mathrm{\Phi }}\sqrt{det(g_{mn}e^{2\mathrm{\Phi }}_m_n)}\sqrt{det\left(\delta _{m}^{}{}_{}{}^{n}+i\frac{e^\mathrm{\Phi }(g_{mp}e^{2\mathrm{\Phi }}_m_p)}{𝒩\sqrt{det(\delta _{m}^{}{}_{}{}^{n}e^{2\mathrm{\Phi }}_m^n)}}H^{np}\right)}$$ $$\frac{1}{4}d^6\xi \sqrt{g}\frac{1}{𝒩^2}H^{mn}H_{mnk}\left(g^{kp}+\frac{e^{2\mathrm{\Phi }}^k^p}{1e^{2\mathrm{\Phi }}^2}\right)\frac{_pa}{\sqrt{(a)^2}}.$$ (39) Since the action (39) is nothing but the M5-brane action for a special choice of the $`D=11`$ metric (29), its variation with respect to the gauge field $`b_2(\xi )`$ and the auxiliary scalar $`a(\xi )`$ has the form of eq. (21), and hence (39) is also invariant under the symmetries (6) and (7) which, as we have seen, produce the self-duality condition (24). To rewrite the transformations and the self–duality condition (24) in the form adapted to the NS5–brane propagating in the $`D=10`$ background, let us introduce the NS5 counterpart of the tensor $`\widehat{}_{mn}^{}`$ (8) $$_{mn}^{}=\frac{2}{\sqrt{g}}\frac{\delta _{kin.NS5}}{\delta H^{mn}},$$ (40) where $`_{kin.NS5}`$ denotes the first (DBI-like) term in the action (39), which is just the DBI–like term of the M5-action (1) written in the $`D=10`$ adapted worldvolume frame. Using (35), it is easy to find the relation between $`\widehat{}_{mn}^{}`$ and $`_{mn}^{}`$ $$_{mn}^{}=\widehat{^{}}_{mn}\sqrt{\frac{(a)^2}{\widehat{aa}}}.$$ (41) Taking into account eqs. (24), (33), (36) and (41) we obtain the following form of the local worldvolume symmetries $$\delta a=0,\delta b_{mn}=2_{[m}a\varphi _{n]}(\xi ),$$ (42) $$\delta a=\phi (\xi ),$$ $$\delta b_{mn}=\frac{\delta a}{\sqrt{(a)^2}}[_{mn}^{}\frac{1}{𝒩^2}H_{mnp}(g^{ps}+\frac{e^{2\mathrm{\Phi }}^p^s}{1e^{2\mathrm{\Phi }}^2})\frac{_sa}{\sqrt{(a)^2}}]$$ (43) and the self-duality equation for the NS5-brane gauge field $`b_2`$ $$_{mn}^{}=\frac{1}{𝒩^2}H_{mnp}(g^{ps}+\frac{e^{2\mathrm{\Phi }}^p^s}{1e^{2\mathrm{\Phi }}^2})\frac{_sa}{\sqrt{(a)^2}},$$ (44) with $`𝒩`$ and $`_{mn}^{}`$ defined in (37) and (40). In addition to the worldvolume diffeomorphisms and the symmetries (42) and (43), the action (39) (by construction) has gauge symmetries (5) and (28). Thus, we have obtained the action describing the worldvolume dynamics of the bosonic 5-brane propagating in the ‘Kaluza–Klein’ part (29) of the IIA $`D=10`$ supergravity background. In the next section we extend this action to describe coupling of the NS5–brane to antisymmetric gauge fields of IIA $`D=10`$ supergravity. ## 4 Coupling to the background gauge fields When the M5-brane couples to the 3-form background field $`\widehat{C}^{(3)}`$ of $`D=11`$ supergravity the field strength $`H_3`$ gets extended by the worldvolume pullback of $`\widehat{C}^{(3)}`$ $$H^{(3)}\widehat{H}^{(3)}=db^{(2)}\widehat{C}^{(3)}.$$ (45) As a result, up to a total derivative, the variation (18) of the action (1) with respect to $`b_{mn}(\xi )`$ and $`a(\xi )`$ acquires an additional term in comparison with eq. (21) $$𝑑\xi ^6\delta =\left[d(da_2)(\delta b_2\delta a_2)+\frac{1}{2}d\widehat{C}_3\delta b_2\right],$$ (46) The symmetries (6) and (7) spoiled by the last term of (46) are restored if to the action (1) one adds the Wess-Zumino term $$S_{WZ}=_^6(\widehat{C}^{(6)}+\frac{1}{2}db^{(2)}\widehat{C}^{(3)}),$$ (47) As it was shown in , the symmetries (6) and (7) uniquely fix the relative factor between $`S_{WZ}`$ and the action (1). In (47) $`\widehat{C}^{(6)}`$ is the pullback of a six–form gauge potential whose field strength is $`D=11`$ Hodge–dual to the field strength of $`\widehat{C}^{(3)}`$ $$d\widehat{C}^{(6)}+\frac{1}{2}\widehat{C}^{(3)}d\widehat{C}^{(3)}={}_{}{}^{}d\widehat{C}^{(3)}$$ (48) In addition to the symmetries (6) and (7) with $`H^{(3)}`$ generalized as in (45), the M5-brane action (1) extended by the Wess–Zumino term (47) is invariant under the following transformations of the antisymmetric gauge fields $$\delta \widehat{C}^{(6)}=d\widehat{\phi }^{(5)}\frac{1}{2}\delta \widehat{C}^{(3)}\widehat{C}^{(3)},\delta \widehat{C}^{(3)}=d\widehat{\phi }^{(2)},$$ (49) $$\delta b^{(2)}=\widehat{\phi }^{(2)}(\widehat{X}(\xi )).$$ (50) To get the form of the coupling of the NS5–brane to the antisymmetric gauge fields of type IIA D=10 supergravity we should dimensionally reduce $`\widehat{C}^{(3)}`$, $`\widehat{C}^{(6)}`$ and the Wess–Zumino term (47) of the M5–brane. The dimensional reduction of $`\widehat{C}^{(3)}`$ produces a ten–dimensional R–R three–form $`C^{(3)}`$ and an NS–NS two–form $`B^{(2)}`$ $$\widehat{C}^{(3)}=\frac{1}{3!}d\widehat{X}^{\underset{¯}{\overset{^}{l}}}d\widehat{X}^{\underset{¯}{\overset{^}{n}}}d\widehat{X}^{\underset{¯}{\overset{^}{m}}}\widehat{C}_{\underset{¯}{\overset{^}{m}}\underset{¯}{\overset{^}{n}}\underset{¯}{\overset{^}{l}}}(\widehat{X})=$$ (51) $$=\frac{1}{3!}dX^{\underset{¯}{l}}dX^{\underset{¯}{n}}dX^{\underset{¯}{m}}C_{\underset{¯}{m}\underset{¯}{n}\underset{¯}{l}}(X)+\frac{1}{2}dX^{\underset{¯}{n}}dX^{\underset{¯}{m}}B_{\underset{¯}{m}\underset{¯}{n}}(X)\left(dydX^{\underset{¯}{l}}A_{\underset{¯}{l}}\right)$$ $$C^{(3)}+B^{(2)},$$ and the dimensional reduction of $`\widehat{C}^{(6)}`$ produces a ten–dimensional five–form $`C^{(5)}`$ and a six–form $`B^{(6)}`$ which are dual to $`C^{(3)}`$ and $`B^{(2)}`$, respectively, $$\widehat{C}^{(6)}=B^{(6)}+C^{(5)},$$ (52) the duality relations can be easely derived by the dimensional reduction of eq. (48). Thus, the field strength of the self-dual gauge field of the NS5–brane coupled to the $`D=10`$ background gauge fields is extended as follows $$H^{(3)}=db^{(2)}C^{(3)}B^{(2)},$$ (53) and the NS5–brane action (39) is enlarged with the following Wess-Zumino term $$S_{WZ}=__6\left[B^{(6)}+C^{(5)}+\frac{1}{2}db^{(2)}C^{(3)}+\frac{1}{2}db^{(2)}B^{(2)}\right],$$ (54) where $`=d\xi ^m(_myA_m)`$ is now the worldvolume pullback of the one–form (51). We have now obtained the action for the NS5–brane propagating in a background of the bosonic sector of type IIA $`D=10`$ supergravity $$S=d^6\xi e^{2\mathrm{\Phi }}\sqrt{det(g_{mn}e^{2\mathrm{\Phi }}_m_n)}\sqrt{det\left(\delta _{m}^{}{}_{}{}^{n}+i\frac{e^\mathrm{\Phi }(g_{mp}e^{2\mathrm{\Phi }}_m_p)}{𝒩\sqrt{det(\delta _{m}^{}{}_{}{}^{n}e^{2\mathrm{\Phi }}_m^n)}}H^{np}\right)}$$ $$\frac{1}{4}d^6\xi \sqrt{g}\frac{1}{𝒩^2}H^{mn}H_{mnk}\left(g^{kp}+\frac{e^{2\mathrm{\Phi }}^k^p}{1e^{2\mathrm{\Phi }}^2}\right)\frac{_pa}{\sqrt{(a)^2}}$$ $$+__6\left(B^{(6)}+C^{(5)}+\frac{1}{2}db^{(2)}C^{(3)}+\frac{1}{2}db^{(2)}B^{(2)}\right).$$ (55) This action is invariant under the worldvolume gauge transformations (5), (42), (43) and (28), with $`H^{(3)}`$ now having the form (53), and under target–space gauge transformations $$\delta C^{(3)}=d\phi ^{(2)}+d\phi ^{(1)}A,\delta B^{(2)}=d\phi ^{(1)},\delta A=d\phi ^{(0)},$$ $$\delta b^{(2)}=\phi ^{(2)}\phi ^{(1)}dy,\delta y=\phi ^{(0)},$$ (56) under which $`\delta H^{(3)}=0`$, and $$\delta B^{(6)}=d\phi ^{(5)}+d\phi ^{(4)}A\frac{1}{2}d\phi ^{(2)}C^{(3)}\frac{1}{2}d\phi ^{(1)}AC^{(3)},$$ $$\delta C^{(5)}=d\phi ^{(4)}\frac{1}{2}d\phi ^{(2)}B^{(2)}+\frac{1}{2}d\phi ^{(1)}C^{(3)}\frac{1}{2}d\phi ^{(1)}AB^{(2)}.$$ (57) Before proceeding with the consideration of the full super–NS5–brane action let us demonstrate how the action of ref. is obtained from eq. (55). ## 5 NS5–brane action in the second order approximation. The action of is a second–order approximation in powers of $`H_{mnk}`$ of the NS5–brane action, with the self-duality condition being regarded as an extra (actually on–shell) constraint. To get the second–order action we should expand (55) in series of $`H^{(3)}`$ and truncate it down to the second order in $`H^{(3)}`$ assuming the worldvolume gauge field to be weak. Since the Wess–Zumino term is already linear and quadratic in $`H`$, we shall write down only the “kinetic” part of the action <sup>7</sup><sup>7</sup>7We should note that our choice of the dimensionally reduced $`C_3`$ and $`C_5`$ differs from that in , so the Wess-Zumino term in Eq. (55) is related to the one of Refs. by the following field redefinition: $$yc^{(0)},b^{(2)}a^{(2)},$$ $$B^{(2)}B^{(2)},C^{(3)}B^{(2)}AC^{(3)},$$ $$C^{(5)}C^{(5)}\frac{1}{2}C^{(3)}B^{(2)},B^{(6)}C^{(5)}A\stackrel{~}{B}^{(6)}.$$ The WZ term of also contains the curl of an auxiliary worldvolume 5-form field which ensures the exact gauge invariance of the WZ term. . To carry out such a truncation the simplest way is to first truncate the M5–brane action (1) and then perform its dimensional reduction. Up to the second order in $`H^{(3)}`$ the M5–brane action has the form $$S=d^6\xi \sqrt{\widehat{g}}\left[1\frac{1}{4}\widehat{H}_{mn}^{}\widehat{H}^{mn}+\frac{1}{4}\widehat{H}^{mn}H_{mnp}\widehat{v}^p+\mathrm{}\right]$$ (58) with the self-duality condition (24) reducing to $$H_{mnl}\widehat{H}_{mnl}^{}=0.$$ (59) Taking into account the expressions $$\widehat{H}_{mnk}^{}=\frac{1}{3!\sqrt{\widehat{g}}}ϵ_{mnkqrs}H^{qrs}$$ and $$ϵ_{mnlpqr}ϵ^{mnlstv}=(3!)^2\delta _p^{[s}\delta _q^t\delta _r^{v]},$$ after some algebra one can rewrite (58) in the following form $$S=d^6\xi \sqrt{\widehat{g}}\left[1\frac{1}{24}H_{mnl}H^{mnl}\frac{1}{8\widehat{aa}}_ma(H^{mnl}\widehat{H}^{mnl})(H_{nlp}\widehat{H}_{nlp}^{})^pa+\mathrm{}\right].$$ (60) Discarding in (60) the term containing the auxiliary field $`a(\xi )`$ and the anti-selfdual tensor $`H\widehat{H}^{}`$ (which is zero on the mass shell (59)), and carrying out the direct dimensional reduction of (60) we recover the NS5–brane action of $$S=d^6\xi e^{2\mathrm{\Phi }}\sqrt{det(g_{mn}e^{2\mathrm{\Phi }}_m_n)}[1\frac{1}{24}(e^{2\mathrm{\Phi }}H_{mnk}H^{mnk}$$ $$+3\frac{e^{4\mathrm{\Phi }}}{1e^{2\mathrm{\Phi }}^2}_mH^{mnk}H_{nkp}^p)+\mathrm{}].$$ (61) Alternatively, the action (61) can be obtained directly by truncating the NS5–brane action (55), and discarding terms containing the auxiliary field and the linearized NS5–brane self–duality condition (44) $$H_{mnl}^{}H_{mnp}\left(\delta _l^p+\frac{e^{2\mathrm{\Phi }}^p_l}{1e^{2\mathrm{\Phi }}^2}\right)=0.$$ ## 6 The $`\kappa `$–symmetric super–NS5–brane action To generalize the results of previous sections to describe the propagation of an NS5–brane in a curved IIA $`D=10`$ target superspace parametrized by ten bosonic coordinates $`X^{\underset{¯}{m}}`$ and 32–component Majorana-spinor fermionic coordinates $`\mathrm{\Theta }^{\underset{¯}{\alpha }}`$ forming a IIA, $`D=10`$ superspace coordinate system $$Z^{\underset{¯}{M}}=(X^{\underset{¯}{m}},\mathrm{\Theta }^{\underset{¯}{\alpha }}),$$ (62) we again start with an M5–brane propagating in a generic D=11 supergravity background parametrized by eleven bosonic coordinates $`\widehat{X}^{\underset{¯}{\overset{^}{m}}}`$ and 32–component Majorana-spinor fermionic coordinates $`\mathrm{\Theta }^{\underset{¯}{\alpha }}`$ forming a $`D=11`$ superspace coordinate system $$\widehat{Z}^{\underset{¯}{\overset{^}{M}}}=(\widehat{X}^{\underset{¯}{\overset{^}{m}}},\mathrm{\Theta }^{\underset{¯}{\alpha }})=(Z^{\underset{¯}{M}},y),$$ (63) where we have separated the eleventh coordinate $`y=X^{10}`$ keeping in mind the dimensional reduction of $`D=11`$ superspace down to type IIA $`D=10`$ superspace. $`D=11`$ superspace geometry is described by a supervielbein $$\widehat{E}^{\underset{¯}{\overset{^}{A}}}=d\widehat{Z}^{\underset{¯}{\overset{^}{M}}}\widehat{E}_{\underset{¯}{\overset{^}{M}}}^{\underset{¯}{\overset{^}{A}}}(\widehat{Z})=(\widehat{E}^{\underset{¯}{\overset{^}{a}}},\widehat{E}^{\underset{¯}{\alpha }}),$$ (64) where $`\underset{¯}{\overset{^}{A}}=(\underset{¯}{\overset{^}{a}},\underset{¯}{\alpha })`$ are locally flat tangent superspace indices, by a superconnection $$\widehat{w}_{\underset{¯}{\overset{^}{A}}}^{\underset{¯}{\overset{^}{B}}}=d\widehat{Z}^{\underset{¯}{\overset{^}{M}}}\widehat{w}_{\underset{¯}{\overset{^}{M}}}{}_{\underset{¯}{\overset{^}{A}}}{}^{\underset{¯}{\overset{^}{B}}}(\widehat{Z}),$$ (65) and by a three–superform generalization of the bosonic gauge field (51) $$\widehat{C}^{(3)}=\frac{1}{3!}d\widehat{Z}^{\underset{¯}{\overset{^}{N}}}d\widehat{Z}^{\underset{¯}{\overset{^}{M}}}d\widehat{Z}^{\underset{¯}{\overset{^}{L}}}C_{\underset{¯}{\widehat{L}\widehat{M}\widehat{N}}}(\widehat{Z}).$$ (66) The supervielbein, the superconnection and the gauge superfield are subject to supergravity constraints which put the superfield formulation of eleven–dimensional supergravity on the mass shell. An explicit form of the $`D=11`$ supergravity constraints relevant to the description of M5–brane dynamics the reader may find in . The super–M5–brane action has the similar form as the bosonic action (1) enlarged with the WZ term (47), where the worldvolume induced metric is now $$\widehat{g}_{mn}=_m\widehat{Z}^{\underset{¯}{\overset{^}{M}}}_n\widehat{Z}^{\underset{¯}{\overset{^}{N}}}\widehat{E}_{\underset{¯}{\overset{^}{N}}}^{\underset{¯}{\overset{^}{a}}}(\widehat{Z})\widehat{E}_{\underset{¯}{\overset{^}{M}}\underset{¯}{\overset{^}{a}}}(\widehat{Z}),$$ (67) and $`\widehat{C}^{(3)}`$ and $`\widehat{C}^{(6)}`$ are worldvolume pullbacks of the three–superform (66) and its six–superform dual . In addition to all symmetries discussed above and target–space superdiffeomorphisms the super–M5–brane action is invariant under the following fermionic $`\kappa `$–symmetry transformations $$i_\kappa \widehat{E}^{\underset{¯}{\overset{^}{a}}}\delta _\kappa \widehat{Z}^{\underset{¯}{\overset{^}{M}}}\widehat{E}_{\underset{¯}{\overset{^}{M}}}^{\underset{¯}{\overset{^}{a}}}(\widehat{Z})=0,i_\kappa \widehat{E}^{\underset{¯}{\overset{^}{\alpha }}}\delta _\kappa \widehat{Z}^{\underset{¯}{\overset{^}{M}}}\widehat{E}_{\underset{¯}{\overset{^}{M}}}^{\underset{¯}{\overset{^}{\alpha }}}(\widehat{Z})=(I\overline{\mathrm{\Gamma }})^{\underset{¯}{\overset{^}{\alpha }}\underset{¯}{\overset{^}{\beta }}}\kappa _{\underset{¯}{\overset{^}{\beta }}},$$ (68) $$\delta _\kappa b_2=i_\kappa \widehat{C}_3\frac{1}{2}d\widehat{Z}^{\underset{¯}{\overset{^}{M}}_3}d\widehat{Z}^{\underset{¯}{\overset{^}{M}}_2}\delta _\kappa \widehat{Z}^{\underset{¯}{\overset{^}{M}}_1}\widehat{C}_{\underset{¯}{\overset{^}{M}}_1\underset{¯}{\overset{^}{M}}_2\underset{¯}{\overset{^}{M}}_3}(\widehat{Z}),\delta _\kappa a=0.$$ where the spinor matrix $`\overline{\mathrm{\Gamma }}`$ has the following expansion in products of $`D=11`$ Dirac matrices $$\overline{\mathrm{\Gamma }}=\frac{\sqrt{\widehat{g}}}{\sqrt{det(\widehat{g}+i\widehat{H}^{})}}\left(\widehat{\mathrm{\Gamma }}^{(6)}+\frac{i}{2}\widehat{H}_{mn}^{}\widehat{v}_l(\widehat{\mathrm{\Gamma }}^{mnl})+\widehat{t}_m\widehat{v}_n(\widehat{\mathrm{\Gamma }}^{mn})\right)$$ (69) $$\widehat{\mathrm{\Gamma }}^{(6)}=\frac{1}{6!}\epsilon ^{m_1\mathrm{}m_6}\widehat{\mathrm{\Gamma }}_{m_1}\mathrm{}\widehat{\mathrm{\Gamma }}_{m_6},\widehat{\mathrm{\Gamma }}_m_m\widehat{Z}^{\underset{¯}{\overset{^}{M}}}\widehat{E}_{\underset{¯}{\overset{^}{M}}}^{\underset{¯}{\overset{^}{a}}}(\widehat{Z})\widehat{\mathrm{\Gamma }}_{\underset{¯}{a}},$$ $$\widehat{t}^m=\frac{1}{8}\epsilon ^{mnkplq}\widehat{H}_{nk}^{}\widehat{H}_{pl}^{}\widehat{v}_q\frac{1}{8}\epsilon ^{mnkplq}\widehat{}_{nk}^{}\widehat{}_{pl}^{}\widehat{v}_q.$$ As is characteristic of all superbranes, for the M5–brane action to be $`\kappa `$–symmetric the superbackground must satisfy the supergravity constraints . When they are taken into account, from (68) we get $$\delta _\kappa \widehat{H}_3=i_\kappa \widehat{F}_4=\widehat{E}^{\underset{¯}{\overset{^}{a}}}\widehat{E}^{\underset{¯}{\overset{^}{b}}}\widehat{E}^{\underset{¯}{\alpha }}\left(\mathrm{\Gamma }_{\underset{¯}{\overset{^}{a}}\underset{¯}{\overset{^}{b}}}(I\overline{\mathrm{\Gamma }})\right)_{\underset{¯}{\alpha }\underset{¯}{\beta }}\kappa ^{\underset{¯}{\beta }},$$ (70) $$\delta _\kappa \widehat{g}_{mn}=4i\widehat{E}_{(m}^{\underset{¯}{\alpha }}\left(\mathrm{\Gamma }_{n)}(I\overline{\mathrm{\Gamma }})\right)_{\underset{¯}{\alpha }\underset{¯}{\beta }}\kappa ^{\underset{¯}{\beta }}.$$ We now turn to the consideration of the super–NS5–brane action. It can be obtained from the super–M5–brane action by the direct dimensional reduction of the $`D=11`$ supergravity superfields. A consistent ansatz for the dimensionally reduced supervielbein (64) was proposed in . This is the following superfield generalization of eq. (27) $$\widehat{E}^{\underset{¯}{a}}=e^{\frac{1}{3}\mathrm{\Phi }(Z)}E^{\underset{¯}{a}},\widehat{E}^{10}=e^{\frac{2}{3}\mathrm{\Phi }(Z)}(dydZ^{\underset{¯}{M}}A_{\underset{¯}{M}}(Z))e^{\frac{2}{3}\mathrm{\Phi }(Z)},$$ (71) $$\widehat{E}^{\underset{¯}{\alpha }}=e^{\frac{1}{6}\mathrm{\Phi }(Z)}E^{\underset{¯}{\alpha }}(Z)+\chi ^{\underset{¯}{\alpha }}(Z),$$ (72) where $`E^{\underset{¯}{A}}(Z)=dZ^{\underset{¯}{M}}E_{\underset{¯}{M}}^{\underset{¯}{A}}=(E^{\underset{¯}{a}},E^{\underset{¯}{\alpha }})`$ are supervielbeins of type IIA $`D=10`$ supergravity, $`\mathrm{\Phi }(Z)`$ is the dilaton superfield, $`A_{\underset{¯}{M}}(Z)`$ are components of the one–form gauge superfield $`A=dZ^{\underset{¯}{M}}A_{\underset{¯}{M}}(Z)`$, and $`\chi ^{\underset{¯}{\alpha }}(Z)`$ is a Grassmann–odd Majorana spinor superfield, which is actually the Grassmann derivative of the dilaton superfield $`\mathrm{\Phi }(Z)`$. The superfields which describe IIA $`D=10`$ supergravity are subject to the constraints which are obtained from the $`D=11`$ supergravity constraints using the ansatz (71),(72) and solving for Bianchi identities. Different forms of these constraints have been considered in , . We do not write the super–NS5–brane action explicitly since it has exactly the same form as Eq. (55) where now the worldvolume induced metric is $$g_{mn}=_mZ^{\underset{¯}{M}}_nZ^{\underset{¯}{N}}E_{\underset{¯}{N}}^{\underset{¯}{a}}(Z)E_{\underset{¯}{M}\underset{¯}{a}}(Z),$$ (73) and all the bosonic background fields are replaced with corresponding superfields, in particular, $`B^{(6)}`$, $`C^{(5)}`$, $`C^{(3)}`$ and $`B^{(2)}`$ are the worldvolume pullbacks of the type IIA $`D=10`$ superforms $$C^{(n)}(Z)=\frac{1}{n!}E^{\underset{¯}{A}_n}\mathrm{}E^{\underset{¯}{A}_1}C_{\underset{¯}{A}_1\mathrm{}\underset{¯}{A}_n}(Z).$$ (74) Note that the spinor superfield $`\chi ^{\underset{¯}{\alpha }}`$ does not appear in the action (55). The super–NS5–brane action is invariant under $`\kappa `$–symmetry transformations obtained from eqs. (68) by substituting into the latter the ansatz (71) and (72) $$i_\kappa E^{\underset{¯}{\overset{^}{\alpha }}}\delta _\kappa Z^{\underset{¯}{M}}E_{\underset{¯}{M}}^{\underset{¯}{\alpha }}(Z)=(I\overline{\mathrm{\Gamma }})^{\underset{¯}{\alpha }\underset{¯}{\beta }}\kappa _{\underset{¯}{\beta }},$$ (75) $$i_\kappa E^{\underset{¯}{a}}=0,i_\kappa =0\delta _\kappa y=i_\kappa E^{\underset{¯}{\alpha }}A_{\underset{¯}{\alpha }}(Z),$$ (76) $$\delta _\kappa b_2=i_\kappa C_3+i_\kappa B_2,\delta _\kappa a=0.$$ where $`A_{\underset{¯}{\alpha }}(Z)`$ is a fermionic component of the Kaluza–Klein connection form $$AdZ^{\underset{¯}{M}}A_{\underset{¯}{M}}=E^{\underset{¯}{\alpha }}A_{\underset{¯}{\alpha }}+E^{\underset{¯}{a}}A_{\underset{¯}{a}}.$$ ## 7 Conclusion and Discussion To summarize, we have obtained the covariant $`\kappa `$-symmetric action for the super–NS5–brane in a IIA $`D=10`$ supergravity background by the direct dimensional reduction of the M-theory super-five-brane action. In addition to worldvolume diffeomorphisms, gauge symmetry, $`\kappa `$–symmetry and background supergravity symmetries the super–NS5–brane action possesses special local symmetries ensuring the covariance of actions with self-dual gauge fields and serving for deriving the self-duality condition directly from the action as a consequence of the equation of motion of the gauge field. An interesting problem for future study is to construct the Lagrangian description of the consistent coupling of a type IIA supergravity action to an NS5-brane source. The latter requires the construction of a duality–symmetric version of type IIA supergravity by the dimensional reduction of the duality-symmetric $`D=11`$ supergravity . The truncation of such a IIA supergravity action shall produce the duality–symmetric version of the $`N=1`$, $`D=10`$ supergravity, which should naturally couple to a heterotic five-brane . Note that recent investigations of interacting brane actions may provide one with a possibility of making this coupling supersymmetric. Another problem for further studying is to perform the T-duality transformation of the complete NS5–brane action and to arrive at a non–linear and supersymmetric action for a type IIB D=10 Kaluza-Klein (KK) monopole. A quadratic approximation for the bosonic part of this action has been constructed in . One of possible ways of deriving appropriate T–duality transformation rules for the antisymmetric gauge fields is to T–dualize the duality–symmetric version of type IIA supergravity to the duality–symmetric version of type IIB supergravity . As it was noted in and proved in the second order approximation in , the type IIB D=10 KK monopole is expected to be a self–dual object under the S–duality symmetry of type IIB supergravity. The construction of the complete action for the type IIB KK monopole should allow one to explicitly verify this statement. Acknowledgements. The authors are grateful to Kurt Lechner, Paolo Pasti and Mario Tonin for interest to this work and valuable discussions and to Christopher Hull, Jeffrey Harvey, Bernard Julia and Kellog Stelle for useful comments. I.B. and A.N. also acknowledge kind hospitality extended to them at the Abdus Salam International Centre for Theoretical Physics where part of this work was done. This work has been partially supported by the Ukrainian GKNT Grant 2.5.1/52 and INTAS Grant No 96-0308.
warning/0003/hep-th0003160.html	ar5iv	text	# Untitled Document hep-th/0003160 Noncommutative Solitons Rajesh Gopakumar, Shiraz Minwalla and Andrew Strominger Jefferson Physical Laboratory, Harvard University Cambridge, MA 02138, USA Abstract We find classically stable solitons (instantons) in odd (even) dimensional scalar noncommutative field theories whose scalar potential, $`V(\varphi )`$, has at least two minima. These solutions are bubbles of the false vacuum whose size is set by the scale of noncommutativity. Our construction uses the correspondence between non-commutative fields and operators on a single particle Hilbert space. In the case of noncommutative gauge theories we note that expanding around a simple solution shifts away the kinetic term and results in a purely quartic action with linearly realised gauge symmetries. Mar. 2000 1. Introduction Quantum field theory on a noncommutative space is of interest for a variety of reasons. It appears to be a self-consistent deformation of the highly constrained structure of local quantum field theory. Noncommutative field theories are nonlocal; unraveling the consequences of the breakdown of locality at short distances may help understanding non-locality in quantum gravity. The discovery of noncommutative quantum field theory in a limit of string theory provides new inroads to the subject. Perturbative aspects of noncommutative field theories have been analyzed in \[2--27\]. This study has thrown up some evidence for the renormalizability of a class of noncommutative field theories, and has revealed an intriguing mixing of the UV and IR in these theories. In this paper we will construct localized classical solutions in some simple noncommutative field theories. We expect these objects to play a role in the quantum dynamics of the theory. We first consider a scalar field with a polynomial potential. A scaling argument due to Derrick shows that, in the commutative case, solitonic solutions do not exist in more than $`1+1`$ dimensions, as the energy of any field configuration can always be lowered by shrinking. Perhaps surprisingly, for sufficiently large noncommutativity parameter $`\theta `$, we will find classically stable solitons in any theory with a scalar potential with more than one $`\mathrm{𝚕𝚘𝚌𝚊𝚕}`$ minimum. These solitons are asymptotic to the true vacuum, and reach a second (possibly false) vacuum in their core. They cannot decay simply by shrinking to zero size because sharply peaked field configurations have high energies in noncommutative field theories. These solitons are metastable in the quantum theory, but by adjusting parameters in the scalar potential, their lifetime can be made arbitrarily long while their mass is kept fixed. Solutions are found corresponding to solitons in $`2l+1`$ dimensions or instantons in $`2l`$ dimensions for any $`l`$. Our construction of these solutions exploits the connection between non-commutative fields and operators in single particle quantum mechanics. Under this correspondence, the $``$ product maps onto usual operator multiplication, and the equation of motion translates into algebraic operator equations. The noncommutative scalar action can be rewritten as the trace over operators (which can be regarded as $`\mathrm{}\times \mathrm{}`$ matrices). This leads to a connection between noncommutative field theories, and zero dimensional matrix models. Next we consider noncommutative $`U(N)`$ Yang-Mills theory. When expanded around a simple solution of the equations of motion, the action takes the simple quartic form (up to constants and topological terms) $$S_{YM}=\frac{1}{4g_{YM}^2}d^{2l}x\delta ^{\mu \lambda }\delta ^{\nu \rho }\mathrm{Tr}\left([\mathrm{\Phi }_\mu ,\mathrm{\Phi }_\nu ][\mathrm{\Phi }_\lambda ,\mathrm{\Phi }_\rho ]\right),$$ where $`\mathrm{\Phi }_\mu `$ are $`N\times N`$ hermitian matrices and all commutators are constructed from the $``$ product. Note that the kinetic term has been shifted away! The usual space-time gauge symmetries act linearly as unitary transformations on the fields $`\mathrm{\Phi }_\mu `$, and the $`\mathrm{\Phi }_\mu =0`$ vacuum leaves even local gauge symmetries unbroken. This construction is similar to that of , in which the kinetic term of Witten’s string field theory action is shifted away. Indeed, our search for such a construction in noncommutative field theory was motivated by the tantalizing analogy, noted in , between noncommutative field theories and string field theories. The existence of the formulation (1.1) of noncommutative gauge theories strengthens the analogy. We also reproduce, as an illustration, the $`U(1)`$ instanton solutions of . Rewriting noncommutative fields as the large $`N`$ limit of matrices, (1.1) is closely related to the IKKT matrix theory \[32\]. Indeed, our construction is essentially equivalent to that presented by Aoki et. al. in this context. Related observations are also made in \[37,,1,,33--36\]. This paper is organized as follows. In section 2 we describe the action for noncommutative scalar field theory. In section 3 we consider the limit, $`\theta \mathrm{}`$, in which the equations simplify considerably. The general solution can be found exactly and is given in terms of quantum mechanical projection operators. In Section 4 we show that there are stable solitons in this limit, as long as the potential has at least two local minima. In section 5 we argue that there are stable solitons at large but finite $`\theta `$ which can be constructed perturbatively in $`\theta ^1`$. In section 6 we turn to the noncommutative gauge theory where the purely quartic action is constructed. The $`U(1)`$ instanton solution of is also reproduced. In an Appendix we give an explicit construction, of the leading $`\frac{1}{\theta }`$ correction to the simplest stable soliton of the scalar field theory. 2. The Noncommutative Scalar Action Consider first a noncommutative field theory of a single scalar $`\varphi `$ in $`(2+1)`$ dimensions with non-commutativity purely in the spatial directions. The spatial $`R^2`$ is parametrized by complex coordinates $`z,\overline{z}`$. The energy functional $$E=\frac{1}{g^2}d^2z\left(_z\varphi _{\overline{z}}\varphi +V(\varphi )\right),$$ where $`d^2z=dxdy`$. (We will comment on the generalization to arbitrary dimensions in the appropriate places.) Fields in this non-local action are multiplied using the Moyal star product, $$\left(AB\right)(z,\overline{z})=e^{\frac{\theta }{2}\left(_z_{\overline{z}^{}}_z^{}_{\overline{z}}\right)}A(z,\overline{z})B(z^{},\overline{z}^{})\|_{z_=z^{}}.$$ Note that in the quadratic part of the action, the star product reduces to the usual product. We seek finite energy (localized) solitons of (2.1). These can also be interpreted as finite action instantons in the two-dimensional euclidean theory. We will, however, refer to the solutions as solitons in the following. Since no solutions exist in the commutative limit $`\theta =0`$ , we begin our search in the limit of large noncommutativity, $`\theta \mathrm{}`$. It is useful to non-dimensionalize the coordinates $`zz\sqrt{\theta }`$, $`\overline{z}\overline{z}\sqrt{\theta }`$. As a result, the $``$ product will henceforth have no $`\theta `$; i.e. it will be given by (2.1) with $`\theta =1`$. Written in rescaled coordinates, the dependence on $`\theta `$ in the energy is entirely in front of the potential term: $$E=\frac{1}{g^2}d^2z\left(\frac{1}{2}(\varphi )^2+\theta V(\varphi )\right)$$ In the limit $`\theta \mathrm{}`$, with $`V`$ held fixed, the kinetic term in (2.1) is negligible in comparison to $`V(\varphi )`$, at least for field configurations varying over sizes of order one in our new coordinates. Our considerations apply to generic potentials $`V(\varphi )`$, but we will, for definiteness, mostly discuss those of polynomial form $$V(\varphi )=\frac{1}{2}m^2\varphi ^2+\underset{j=3}{\overset{r}{}}\frac{b_j}{j}\varphi ^j.$$ We have, of course, abbreviated $$\varphi ^j=\varphi \varphi \mathrm{}\varphi .$$ 3. Scalar Solitons in the $`\theta =\mathrm{}`$ Limit After neglecting the kinetic term, the energy $$E=\frac{\theta }{g^2}d^2zV(\varphi ),$$ is extremised by solving the equation $$\frac{V}{\varphi }=0.$$ For instance, (3.1) is $$m^2\varphi +b_3\varphi \varphi =0$$ for a cubic potential and $$m^2\varphi +b_3\varphi \varphi +b_4\varphi \varphi \varphi =0$$ for a quartic potential. If $`V(\varphi )`$ were the potential in a commutative scalar field theory, the only solutions to (3.1) would be the constant configurations $$\varphi =\lambda _i,$$ where $`\lambda _i\{\lambda _1,\lambda _2,\mathrm{},\lambda _k\}`$ are the various real extrema of the function $`V(x)`$<sup>1</sup> For $`V(\varphi )`$ as in (2.1), $`\lambda _i`$ are the real roots of the equation $`m^2x+_{j=3}^rb_jx^{j1}=0.`$. As we shall see below, the derivatives in the definition of the star product allow for more interesting solutions of (3.1). 3.1. A Simple Nontrivial Solution A non-trivial solution to (3.1) can easily be constructed. Given a function $`\varphi _0(x)`$ that obeys $$\left(\varphi _0\varphi _0\right)(x)=\varphi _0(x),$$ it follows by iteration that $`\varphi _0^n(x)=\varphi _0(x)`$,<sup>2</sup> This equation and its solution has also appeared in earlier work involving the Moyal Product. See \[38,,39\]. and that $`f\left(a\varphi _0(x)\right)=f(a)\varphi _0(x)`$ (fields in $`f`$ are multiplied using the star product). In particular, $`\lambda _i\varphi _0(x)`$ solves (3.1) when $`\lambda _i`$ is an extremum of $`V(x)`$. Thus, in order to find a solution of (3.1), it is sufficient to find a function that squares to itself under the star product. We proceed to construct such a function below. If we take the ordinary product of a smooth function of width $`\mathrm{\Delta }`$ with itself, the spatial size of the function shrinks to a fraction of $`\mathrm{\Delta }`$, which is why non-constant functions never square to themselves! The non-locality of the star product, however, introduces an additional effect, adding roughly<sup>3</sup> The added width is actually $`K`$, the typical momentum in the Fourier transform of the function. For a function of size $`\mathrm{\Delta }`$ with no oscillations, $`K\frac{1}{\mathrm{\Delta }}`$. For a function of size $`\mathrm{\Delta }`$ with $`n`$ oscillations, $`K\frac{n}{\mathrm{\Delta }}`$. $`\frac{1}{\mathrm{\Delta }}`$ to the width of the product. This makes it possible for a lump of approximately unit size to square to itself under the star product. Consider a gaussian packet of the form $$\psi _\mathrm{\Delta }(r)=\frac{1}{\pi \mathrm{\Delta }^2}e^{\frac{r^2}{\mathrm{\Delta }^2}},$$ with radial width $`\mathrm{\Delta }`$ (here $`r^2=x^2+y^2`$). The star product of $`\psi _\mathrm{\Delta }`$ with itself is easily computed by passing to momentum space, $$\stackrel{~}{\psi }_\mathrm{\Delta }(k)=e^{ikx}\psi _\mathrm{\Delta }(x)d^2x=e^{\frac{k^2\mathrm{\Delta }^2}{4}},$$ $$\begin{array}{cc}\hfill \left(\stackrel{~}{\psi }_\mathrm{\Delta }\stackrel{~}{\psi }_\mathrm{\Delta }\right)(p)& =\frac{1}{(2\pi )^2}d^2k\stackrel{~}{\psi }_\mathrm{\Delta }(k)\stackrel{~}{\psi }_\mathrm{\Delta }(pk)e^{\frac{i}{2}ϵ_{\mu \nu }k^\mu (pk)^\nu }\hfill \\ & =\frac{1}{2\pi \mathrm{\Delta }^2}e^{\frac{p^2}{8}\left(\mathrm{\Delta }^2+\frac{1}{\mathrm{\Delta }^2}\right)}.\hfill \end{array}$$ Therefore $$\left(\psi _\mathrm{\Delta }\psi _\mathrm{\Delta }\right)(x)=\frac{1}{\pi ^2\mathrm{\Delta }^2(\mathrm{\Delta }^2+\frac{1}{\mathrm{\Delta }^2})}\mathrm{exp}\left[\frac{2r^2}{\mathrm{\Delta }^2+\frac{1}{\mathrm{\Delta }^2}}\right].$$ In particular<sup>4</sup> We note in passing that in the limit $`\mathrm{\Delta }0`$, (3.1) reduces to $`\delta ^2(x)\delta ^2(x)=\frac{1}{(2\pi )^2}`$., when $`\mathrm{\Delta }^2=1`$, the gaussian squares to itself (up to a factor of $`2\pi `$). That is, $$\varphi _0(x)=2\pi \psi _1(x)=2e^{r^2}$$ solves (3.1) and $`\lambda _i\varphi _0(x)`$ solves (3.1). 3.2. The General Solution In order to find all solutions of (3.1) we will exploit the connection between Moyal products and quantization. Given a $`C^{\mathrm{}}`$ function $`f(q,p)`$ on $`R^2`$ (thought of as the phase space of a one-dimensional particle), there is a prescription which uniquely assigns to it an operator $`O_f(\widehat{q},\widehat{p})`$, acting on the corresponding single particle quantum mechanical Hilbert space, $``$. It is convenient for our purposes to choose the Weyl or symmetric ordering prescription $$O_f(\widehat{q},\widehat{p})=\frac{1}{(2\pi )^2}d^2k\stackrel{~}{f}(k)e^{i\left(k_q\widehat{q}+k_p\widehat{p}\right)},$$ where $$\stackrel{~}{f}(k)=d^2xe^{i(k_qq+k_pp)}f(q,p),$$ and $$[\widehat{q},\widehat{p}]=i.$$ With this prescription, it may be verified that $$\frac{1}{2\pi }𝑑p𝑑qf(q,p)=\mathrm{Tr}_{}O_f,$$ and that the Moyal product of functions is isomorphic to ordinary operator multiplication $$O_fO_g=O_{fg}.$$ In order to solve any algebraic equation involving the star product, it is thus sufficient to determine all operator solutions to the equation in $``$. The functions on phase space corresponding to each of these operators may then be read off from (3.1). We will now employ this procedure to find all solutions of (3.1). As noted above, any solution to (3.1) may be rescaled into a solution of (3.1). Particular solutions of (3.1) may thus be obtained by constructing operators in $``$ that obey (3.1), i.e. $`O_\varphi ^2=O_\varphi `$. This equation is solved by any projection operator in $``$. $``$ possesses an infinite number of projection operators, which can be classified by the dimension of the subspace they project onto. Each class contains a large continuous infinity of operators, each of which, upon rescaling, yields a solution to (3.1). The most general solution to (3.1) hence takes the form $$O=\underset{j}{}a_jP_j$$ where $`\{P_j\}`$ are mutually orthogonal projection operators onto one dimensional subspaces, with $`a_j`$ taking values in the set $`\{\lambda _i\}`$ of real extrema of $`V(x)`$. In order to obtain the functions in space corresponding to the solutions (3.1), it is convenient to choose a particular basis in $``$. Let $`\|n`$ represent the energy eigenstates of the one dimensional harmonic oscillator whose creation and annihilation operators are defined by $$a=\frac{\widehat{q}+i\widehat{p}}{\sqrt{2}};a^{}=\frac{\widehat{q}i\widehat{p}}{\sqrt{2}}.$$ Note that $`a\|n=\sqrt{n}\|n1`$ and $`a^{}\|n=\sqrt{n+1}\|n+1`$. Any operator may be written as a linear combination of the basis operators $`\|mn\|`$’s, which, in turn, may be expressed in terms of $`a`$ and $`a^{}`$ as $$\|mn\|=:\frac{a^m}{\sqrt{m}!}e^{a^{}a}\frac{a^n}{\sqrt{n!}}:$$ where double dots denote normal ordering. We will first describe operators of the form (3.1) that correspond to radially symmetric functions in space. As $`a^{}a\frac{r^2}{2}`$, operators corresponding to radially symmetric wavefunctions are functions of $`a^{}a`$. From (3.1), the only such operators are linear combinations of the diagonal projection operators $`\|nn\|=\frac{1}{n!}:a^ne^{a^{}a}a^n:`$. Hence all radially symmetric solutions of (3.1) correspond to operators of the form $`O=_na_n\|nn\|`$, where the numbers $`a_n`$ can take any values in the set $`\{\lambda _i\}`$. We now translate these operator solutions back to field space. From the Baker-Campbell-Hausdorff formula $$e^{i\left(k_q\widehat{q}+k_p\widehat{p}\right)}=e^{i\left(k_{\overline{z}}a+k_za^{}\right)}=e^{\frac{k^2}{4}}:e^{i\left(k_{\overline{z}}a+k_za^{}\right)}:,$$ where $$k_z=\frac{k_x+ik_y}{\sqrt{2}},k_{\overline{z}}=\frac{k_xik_y}{\sqrt{2}},k^2=2k_zk_{\overline{z}}.$$ Any operator $`O`$ expressed as a normal ordered function of $`a`$ and $`a^{}`$, $`f_N(a,a^{})`$, can be rewritten in Weyl ordered form as follows. By definition, $$O=:f_N(a,a^{}):=\frac{1}{(2\pi )^2}d^2k\stackrel{~}{f}_N(k):e^{i\left(k_{\overline{z}}a+k_za^{}\right)}:.$$ Using (3.1), (3.1) may be rewritten as $$O=\frac{1}{(2\pi )^2}d^2k\stackrel{~}{f}_N(k)e^{\frac{k^2}{4}}e^{i\left(k_{\overline{z}}a+k_za^{}\right)}.$$ Thus, the momentum space function $`\stackrel{~}{f}`$ associated with the operator $`O`$, according to the rule (3.1) is $$\stackrel{~}{f}(k)=e^{\frac{k^2}{4}}\stackrel{~}{f}_N(k).$$ For the operator $`O_n=\|nn\|`$ we find, using (3.1) and (3.1), that the corresponding normal ordered function $`\stackrel{~}{\varphi }_N^{(n)}(k)=2\pi e^{\frac{k^2}{2}}L_n(\frac{k^2}{2})`$. (3.1) then becomes $$\|nn\|=\frac{1}{(2\pi )}d^2ke^{\frac{k^2}{4}}L_n(\frac{k^2}{2})e^{i\left(k_{\overline{z}}a+k_za^{}\right)}$$ where $`L_n(x)`$ is the $`n^{th}`$ Laguerre polynomial. The field $`\varphi _n(x,y)`$ that corresponds to the operator $`O_n=\|nn\|`$ is, therefore, $$\varphi _n(r^2=x^2+y^2)=\frac{1}{(2\pi )}d^2ke^{\frac{k^2}{4}}L_n(\frac{k^2}{2})e^{ik.x}=2(1)^ne^{r^2}L_n(2r^2).$$ Note that $`\varphi _0(r^2)`$ is precisely the gaussian solution found in Sec. 3.1. In summary, (3.1) has an infinite number of real radial solutions, given by $$\underset{n=0}{\overset{\mathrm{}}{}}a_n\varphi _n(r^2)$$ where $`\varphi _n(r^2)`$ is given by (3.1) and each $`a_n`$ takes values in $`\{\lambda _i\}`$. In order to generate all non radially symmetric solutions to (3.1), we rewrite (3.1) in operator language, using (3.1) as $$E=\frac{2\pi \theta }{g^2}\mathrm{Tr}V(O_\varphi ).$$ (3.1) is manifestly invariant under unitary transformations of $`O_\varphi `$ and so has a $`U(\mathrm{})`$ global symmetry. In other words, if $`O`$ is a solution to the equation of motion, so is $`UOU^{}`$, where $`U`$ is any unitary operator acting on $``$. A general Hermitian operator (one that corresponds to a real field $`\varphi `$) may be obtained by acting on a diagonal operator (i.e. an operator that corresponds to a radially symmetric field configuration) by an element of the $`U(\mathrm{})`$ symmetry group (since any hermitian operator is unitarily diagonalizable). Thus every solution to (3.1) may be obtained from a radially symmetric solution by means of $`U(\mathrm{})`$ symmetry transformations. Therefore solutions to (3.1) consist of disjoint infinite dimensional manifolds labelled by the set of eigenvalues of the corresponding operator. Points on the same manifold can be mapped into each other by $`U(\mathrm{})`$ transformations. Each manifold includes several<sup>5</sup> Distinct diagonal operators having the same eigenvalues lie on the same manifold, being related by the “Weyl” subgroup of $`U(\mathrm{})`$ that permutes eigenvalues. diagonal operators (radially symmetric solutions). We will have more to say about the moduli space of these solutions in the next section. As all solutions are related to radially symmetric solutions by a symmetry transformation, we will mostly discuss only radially symmetric solutions. 3.3. UV/IR Mixing $`\varphi _0(r^2)`$, the Gaussian solution worked out in subsec 3.1, is a lump of unit size centred at the origin, as shown in Fig. 1. Fig. 1: A plot of $`\varphi _0(r)`$ versus $`r`$. The solution is a blob centred at the origin. $`\varphi _n(r)`$, at large $`n`$, looks quite different (see Fig. 2.). It is a solution of size $`\sqrt{n}`$ that undergoes $`n`$ oscillations<sup>6</sup> Using asymptotic formulae for Laguerre polynomials we find $$\varphi _n(r)=\{\begin{array}{cc}2(1)^n\hfill & r\sqrt{\frac{1}{8n}}\hfill \\ \frac{2(1)^n}{(2\pi ^2r^2)^{\frac{1}{4}}}\mathrm{cos}(\sqrt{2n}2r\frac{\pi }{4})\hfill & \sqrt{\frac{1}{8n}}r\sqrt{2n})\hfill \\ \frac{2(2r^2)^n}{n!}e^{r^2}\hfill & r\sqrt{2n}\hfill \\ .\hfill & \end{array}$$ in that interval, with oscillation period $`\frac{1}{\sqrt{n}}`$. $`\varphi _n(r^2)`$ thus receives significant contributions from momenta up to $`\sqrt{n}`$ in momentum space. These solutions exemplify the UV-IR mixing pointed out in ; oscillations with frequency $`\sqrt{n}`$ produce an object of size $`\sqrt{n}`$ (instead of $`\frac{1}{\sqrt{n}}`$) in a noncommutative theory. Fig. 2: A plot of $`\varphi _{30}(r)`$ versus $`r`$. 3.4. Generalization to Higher Dimensions All considerations of the preceding subsections may easily be generalized to higher dimensions. Consider a scalar field theory in $`2l+1`$ dimensions with non-commutativity only in the spatial directions. By a choice of axes, the $`2l\times 2l`$ dimensional noncommutativity matrix $`\mathrm{\Theta }`$ may always be brought into block diagonal form. In other words, it is possible to choose spatial coordinates $`z_i,\overline{z}_{\overline{j}}`$ $`(i,\overline{j}=1\mathrm{}l)`$, in terms of which the non-commutativity matrix $`\mathrm{\Theta }_{i\overline{j}}=\theta _i\delta _{i\overline{j}}`$, $`\mathrm{\Theta }_{ij}=\mathrm{\Theta }_{\overline{i},\overline{j}}=0`$., As before we consider the limit where $`\theta _i`$ are uniformly taken to $`\mathrm{}`$ and non-dimensionalize $`z_iz_i\sqrt{\theta _i}`$. As in the previous subsections, the kinetic term in the action may be dropped in this limit. Solutions to the equations of motion (3.1) are once again in correspondence with operator solutions to the same equations; the operators in question now acting on $`\times \times \mathrm{}\times `$, $`l`$ copies of the Hilbert space of the previous subsection. The general solution to (3.1) once again takes the form (3.1) in terms of projection operators on this space. As in the previous subsection, the general solution may be obtained from diagonal solutions via $`U(\mathrm{})`$ rotations. Diagonal solutions to (3.1) are given by $$O=\underset{\stackrel{}{n}}{}a_\stackrel{}{n}\|\stackrel{}{n}\stackrel{}{n}\|\underset{\stackrel{}{n}}{}a_\stackrel{}{n}\underset{i}{}\varphi _{n_i}(\|z_i\|^2),$$ where $`\stackrel{}{n}`$ is shorthand for the set of quantum numbers $`\{n_i\}`$ for the $`l`$ dimensional oscillator and $`\varphi _{n_i}`$ are defined in (3.1). As in (3.1), the coefficients $`a_\stackrel{}{n}`$ take values in $`\{\lambda _i\}`$. A subset of the solutions (3.1) are actually invariant under $`SO(2l)`$ rotations and can be written in terms of associated Laguerre polynomials. These are displayed in Sec.A.3 of the Appendix. In summary, in the limit of maximal noncommutativity, the construction of solitons in two spatial dimensions generalizes almost trivially to every even spatial dimension. 4. Stability and Moduli Space at $`\theta =\mathrm{}`$ In this section we study the stability of the solitons constructed in the previous section. We will also describe the moduli space of stable solitons. 4.1. Stability at $`\theta =\mathrm{}`$ We wish to examine the stability of the radial solution $$\varphi (r^2)=\underset{n=0}{\overset{\mathrm{}}{}}\lambda _{a_n}\varphi _n(r^2)$$ to small fluctuations. Since any $`U(\mathrm{})`$ rotation does not change the energy of our solution (4.1), it is sufficient to study the stability of (4.1) to radially symmetric fluctuations. These are most conveniently parameterized as deformations of the eigenvalues. The energy for an arbitrary radially symmetric state $`\varphi (r^2)=_{n=0}^{\mathrm{}}c_n\varphi _n(r^2)`$ is $$E=\frac{2\pi \theta }{g^2}\underset{n=0}{\overset{\mathrm{}}{}}V(c_n).$$ The solution $`c_n=\lambda _{a_n}`$ is manifestly an extremum of $`S`$, as, by definition, $`\lambda _{a_i}`$ are extrema of the function $`V(x)`$. Clearly (4.1) is a local minimum of the energy (and so a stable solution) if, and only if, $`\lambda _{a_n}`$ is a local minimum of $`V(x)`$ for all $`0n\mathrm{}`$. As an example consider the cubic potential of Fig. 3. with a maximum at $`\lambda =1`$. In this case, all $`\lambda _{a_n}`$ in (4.1) are either zero or -1. The only stable solution is that for which all $`\lambda _{a_n}=0`$, i.e. the vacuum. The solution $`\varphi _0(r^2)`$, for instance, is unstable, as the energy of this field configuration is decreased by scaling this solution by a constant near unity. This instability shows up as a negative eigenvalue of the quadratic form for fluctuations about $`\varphi _0(x)`$; the corresponding eigenmode $`\delta \varphi _0`$ is $`\varphi _0`$. Fig. 3: The $`\varphi ^3`$ theory with an unstable extremum On the other hand the field theory with $`V(\varphi )`$ (say, for a quartic potential) graphed in Fig. 4 has stable solitons; these are solutions of the form $`\varphi (r^2)=_{n=0}^{\mathrm{}}\lambda _{c_n}\varphi _n(r^2)`$ with $`\lambda _{c_n}`$ taking the values of the minima – $`0`$ or $`\lambda 1.4`$ for all $`n`$. In particular $`\lambda \varphi _0(r^2)`$ is a stable solution, manifestly stable to rescalings. Again, one may check that the quadratic form for fluctuations about $`\lambda \varphi _0(r^2)`$ is positive. In particular, $`\delta \varphi \varphi _0`$ is an eigenmode of this quadratic form with positive eigenvalue. Fig. 4: A $`\varphi ^4`$ potential with two minima. The stability of $`\varphi (r^2)=\lambda \varphi _0(r^2)`$ in the previous example may qualitatively be understood as follows. $`\varphi _0`$ is a Gaussian of height $`2\lambda `$. Far away from the origin, $`\varphi _0(x)=0`$, but near $`x=0`$, $`\varphi _0(x)`$ is in the vicinity of the second vacuum. In other words, the static solution corresponds to a bubble of the “false” vacuum. The area of the bubble is of order one (or $`\theta `$ in our original coordinates), the non-commutativity scale. In a commutative theory such a bubble would decay by shrinking to zero size. Noncommutativity prevents the bubble from shrinking to a spatial size smaller than $`\sqrt{\theta }`$. In order to decay, $`\varphi _0`$ actually has to scale to zero - but that process involves going over the hump in the potential and so is classically forbidden. Fig. 5: Profile of the Gaussian soliton with a false vacuum region (above the horizontal bar) of radius 1. The energy of this soliton is proportional to the vacuum energy density $`\frac{V(\lambda )}{g^2}`$ at the ‘false’ vacuum times the volume of the soliton $`\theta `$. It is remarkable that the energy of the soliton is completely insensitive to the value of the scalar potential at any point except $`\varphi =\lambda `$. Thus the mass of the soliton is unchanged if the height of the barrier in $`V(\varphi )`$ (between $`\varphi =\lambda `$ and $`\varphi =0`$, see Fig. 4.) is taken to infinity while $`V(\lambda )`$ is kept fixed. This is true even though $`\varphi _0(r)`$, the solitonic field configuration corresponding to $`\lambda \|00\|`$, decreases continuously from $`\varphi =2\lambda `$ at $`r=0`$ to $`\varphi =0`$ at $`r=\mathrm{}`$! Consider a 2+1 dimensional scalar theory, noncommutative only in space, at infinite $`\theta `$. Using the correspondence between functions and operators (matrices) described in the previous section, the noncommutative scalar field theory is equivalent to the matrix quantum mechanics of an $`N\times N`$ hermitian matrix $`H`$, at infinite $`N`$, with the usual relativistic kinetic term $`\mathrm{Tr}\left(_tH\right)^2`$, and a potential $`\mathrm{Tr}\left(V(H)\right)`$. The amplitude for an eigenvalue of $`H`$ to tunnel from $`\lambda `$ to $`0`$ is exponentially suppressed by the area under the potential barrier in Fig. 4., and goes to zero as this area is taken to infinity. Thus the finite mass soliton $`\lambda \|00\|`$ is stable, even quantum mechanically, in this limit. The $`U(\mathrm{})`$ symmetry of (3.1) is spontaneously broken by every nonzero solution, $`\varphi (x)`$, of (3.1). As a consequence, every solution has a number of exact zero modes (Goldstone modes) corresponding to small displacements about $`\varphi (x)`$ on the manifold of solutions. As $`R_{nm}=\|nm\|+\|mn\|`$ and $`S_{nm}=i(\|nm\|\|mn\|)`$ are the generators of $`U(\mathrm{})`$, these zero modes are given by the nonzero elements of $`\delta \varphi [R_{nm},\varphi ],[S_{nm},\varphi ]`$. The $`U(\mathrm{})`$ group of symmetry transformations that generates these zero modes is certainly not manifest (at least to the untrained eye) in the energy written in coordinate space in the form (3.1). In addition to the two translations, (3.1) possesses three manifest local symmetries, corresponding to a linear change of the coordinates $`x,y`$ by an SL(2,R) matrix. The remaining $`U(\mathrm{})`$ transformations act non-locally on $`\varphi (x,y)`$, according to $`\varphi ^{}(x,y)=\left(U\varphi U^{}\right)(x,y)`$ where $`U(x,y)`$ is any function that obeys $`UU^{}=1`$ (such functions correspond to $`U(\mathrm{})`$ operators under the map (3.1)). All arguments in this subsection may be applied (after straightforward generalizations) to higher dimensional solitons. 4.2. Multi Solitons In this subsection we will qualitatively describe a part of the moduli space of stable solitons (at $`\theta =\mathrm{}`$) in the simple case of the potential graphed in Fig. 4 with a single non-zero minimum at $`\varphi =\lambda `$. The stable solitons can be characterized by their ‘level’ (number of $`\lambda `$ eigenvalues). All stable level one solitons correspond to operators of the form $$\lambda U\|00\|U^{}$$ where $`U`$ is a unitary operator. As mentioned above, the set of level one solitons span an infinite dimensional manifold parameterized by $`U(N)/U(N1)`$ (for $`N=\mathrm{}`$). The soliton looks very different at different points on the manifold. $`U=I`$ in (4.1) corresponds to the gaussian blob of Fig. 1. If $`U`$ happens to be a unitary transformation that maps $`\|0`$ to $`\|m`$, for large $`m`$, the corresponding wave function is qualitatively similar to that in Fig. 2. When $`U=e^{a^{}za\overline{z}}`$ is the generator of translations, the operator in (4.1), $`\lambda \|zz\|`$, is proportional to the projection operator onto a gaussian centred around $`z=\frac{1}{\sqrt{2}}(x+iy)`$. (Here $`\|z=e^{\frac{\|z\|^2}{2}}e^{a^{}z}\|0`$ is the usual coherent state.) Again, if $`U`$ corresponds to one of the $`SL(2,R)`$ operators, we obtain squeezed states; gaussians elongated in the $`y`$ direction and shrunk in the $`x`$ direction. And so on. Turn now to solitons at arbitrary level $`n`$. All such solitons may be obtained by acting on $$\lambda (\varphi _0+\varphi _1+\mathrm{}+\varphi _{n1})$$ by arbitrary unitary transformations. The manifold of solutions thus generated is parameterized by $`\frac{U(N)}{U(Nn)}`$ (and has dimension $`d_n2nN`$) where $`N\mathrm{}`$. Notice that $`d_nnd_1`$. This fact has a nice explanation; in a particular limit the manifold of level $`n`$ solutions reduces to $`n`$ copies of level $`1`$ solitons very far from each other. This conclusion follows from the observation that the operator that represents $`n`$ widely separated level one solitons (with centres $`z_j`$), for instance $$M=\lambda \underset{j}{}\|z_jz_j\|$$ is approximately a level $`n`$ soliton (and exponentially close to a true level $`n`$ soliton) when $`\|z_iz_j\|\mathrm{}`$ for all $`i,j`$. We demonstrate this explicitly below for the case $`n=2`$. Using $`z\|z=e^{2\|z\|^2}`$, it is easy to check that the kets $$\|z_\pm =\frac{\|z\pm \|z}{\sqrt{2(1\pm e^{2\|z\|^2})}}$$ are orthogonal. From (3.1) we conclude that the projector $$O_z=\lambda \left(\|z_+z_+\|+\|z_{}z_{}\|\right)=\lambda \frac{\|zz\|+\|zz\|+e^{2\|z\|^2}\left(\|zz\|+\|zz\|\right)}{(1e^{4\|z\|^2})}$$ corresponds to a level 2 solution. Up to corrections of order $`e^{2\|z\|^2}`$, $`O_z`$ is equal to $`\|zz\|+\|zz\|`$, the superposition of field configurations corresponding to two widely separated level one solitons <sup>7</sup> It is curious that the kinetic energy of this field configuration is independent of $`z`$ indicating that there is no force between the two solitons even to next leading order in $`\frac{1}{\theta }`$. . We conclude that a part of the level $`n`$ moduli space describes $`n`$ widely separated level one solitons. We have, so far, worked in the strict limit $`\theta =\mathrm{}`$. The picture developed in this limit is qualitatively modified at large but finite $`\theta `$, as we will describe in the next section. 5. Scalar Solitons at Large but Finite $`\theta `$. We have argued that, under certain conditions on $`V(\varphi )`$, (3.1) has an infinite number of stable solutions. Each solution has an infinite number of exact zero modes, the Goldstone modes of the spontaneously broken $`U(\mathrm{})`$ symmetry of (3.1). At finite $`\theta `$, the kinetic term in (2.1) explicitly breaks this $`U(\mathrm{})`$ symmetry down to the Euclidean group in 2 dimensions. Finite $`\theta `$ effects may thus be expected to 1. Lift the $`\theta =\mathrm{}`$ manifold of solutions to a discrete set of solutions. 2. Give (positive or negative) masses to the $`U(\mathrm{})`$ Goldstone bosons about these discrete solutions. In Appendix A we will argue that, at large enough $`\theta `$, corresponding to every radially symmetric solution $`s`$ of (3.1), there is a radially symmetric saddle point of (2.1), that reduces to $`s`$ as $`\theta \mathrm{}`$. It is likely that these are the only saddle points of (2.1). Not all these radially symmetric solutions are stable, however. In fact, it might seem likely that some of the infinite number of zero modes, at $`\theta =\mathrm{}`$, about each solution $`s`$, might become tachyonic at finite $`\theta `$. If this were true, (2.1) would have no classically stable extremum at any finite $`\theta `$, no matter how large. We will find that is not the case. In subsection 5.1 below we will argue that any small perturbation of (3.1) must preserve the existence of at least one classically stable level one soliton. In subsection 5.2 we will identify this soliton to be the one near the gaussian $`\lambda \varphi _0(r^2)`$. 5.1. Existence of A Stable Soliton For definiteness, through the rest of this section we assume that the potential $`V(\varphi )`$ has the shape shown in Fig. 4. In particular, it is positive definite. Let the stable extremum of $`V`$ occur at $`\varphi =\lambda `$ and the unstable extremum at $`\varphi =\beta `$, ($`\lambda <\beta <0`$). Consider, first, (3.1) i.e. the energy functional in the limit where we neglect the kinetic term. We will show that any path in field space leading from the soliton $`\lambda \varphi _0(r^2)`$ to the the vacuum passes through a point whose energy is larger than $`\frac{2\pi \theta }{g^2}V(\beta )`$. Since the energy of the stable soliton is $`\frac{2\pi \theta }{g^2}V(\lambda )<\frac{2\pi \theta }{g^2}V(\beta )`$, every path from the soliton to the vacuum must pass over a barrier of height $`𝒪(\frac{\theta }{g^2})`$. The energy evaluated on an operator $`A`$ is $$E=\frac{2\pi \theta }{g^2}\mathrm{Tr}(V(A))=\frac{2\pi \theta }{g^2}\underset{n=1}{\overset{\mathrm{}}{}}V(c_n)$$ where $`c_n`$ are the eigenvalues of $`A`$. Since $`V`$ is positive definite, $$E\frac{2\pi \theta }{g^2}V(b),$$ where $`b`$ is the smallest eigenvalue of $`A`$. Consider a path in field, or operator space, leading from $`\lambda \varphi _0`$ to the vacuum. At the beginning of this path $`b=\lambda `$. At its end $`b=0`$. Since $`\lambda <\beta <0`$, any smooth path with these endpoints must have a point at which $`b=\beta `$. At that point $`E>\frac{2\pi \theta }{g^2}V(\beta )`$, as was to be shown. Now include the the kinetic term in (3.1). Barring singular behaviour, this changes the energies of all field configurations by terms of $`𝒪(\frac{1}{g^2})`$. For large enough $`\theta `$, the arguments of the previous paragraph imply that the field configuration that describe the level one soliton at $`\theta =\mathrm{}`$ cannot decay to the vacuum. Hence there must exist at least one stable soliton near one of the unperturbed level one solutions. In fact, as we will show in the next subsection, there is a stable soliton near the gaussian $`\lambda \varphi _0(r^2)`$. In the Appendix we will present an approximate construction of this solution at large but finite $`\theta `$. A similar argument demonstrates the existence of at least one stable solution at level $`n`$. 5.2. Approximate Description of the Stable Soliton All level one solutions to (3.1) take the form $`\lambda U\|00\|U^{}`$ where $`U`$ is a unitary operator. We wish to determine the contribution of the kinetic term to the energy of such an operator. The kinetic term in (2.1) for an operator $`A`$ is $$K=\frac{2\pi }{g^2}\mathrm{Tr}[a,A][A,a^{}].$$ Setting $`A=\lambda U\|00\|U^{}`$ we find $$\frac{g^2K(U)}{2\pi \lambda ^2}=1+\underset{k}{}2k\|U_{k,0}\|^22\|\underset{k}{}\sqrt{k+1}U_{k,0}U_{k+1,0}^{}\|^2.$$ We expand (5.1) to quadratic order in deviations from $`U=I`$. Choose $`U_i=U_{i,0}`$ for $`i1`$ as the coordinates for this expansion ($`\|U_{00}\|`$ is determined in terms of $`U_i`$ as $`U`$ is unitary). To quadratic order in $`U_i`$ $$\frac{g^2K(U)}{2\pi \lambda ^2}=1+2\underset{k=2}{\overset{\mathrm{}}{}}k\|U_k\|^2.$$ As $`U_1`$ and $`\overline{U}_1`$ do not appear in (5.1), they parameterize flat directions of $`K(U)`$ (to quadratic order). This was to be expected. Any localized extremum of (2.1) has two exact translational zero modes. Infinitesimally, $`U_{01}`$ and its complex conjugates act as derivatives on $`\varphi _0(r^2)`$, generating these zero modes. Modulo these zero modes, the fluctuation matrix about $`U=1`$ is positive definite. While $`K(U)`$ has several critical points other than $`U=I`$, it has no further local minima. For example, $`U=U^{(m)}`$, the unitary transformation that rotates $`\|00\|`$ to $`\|mm\|`$, is an unstable critical point of $`K(U)`$ for all $`m`$. In fact $`U=U^{(m)}`$ is unstable to decay into $`U=I`$. This may be demonstrated by considering the path in field space $`\|\alpha \alpha \|`$ where $`\|\alpha =\mathrm{cos}\alpha \|0+\mathrm{sin}\alpha \|m`$. (5.1) evaluated on such a path is equal to $`1+2m\mathrm{sin}^2\alpha `$ (for $`m>1`$; $`1+2\mathrm{sin}^4\alpha `$ for $`m=1`$) indicating that the state $`\|mm\|`$ can decay to $`\|00\|`$. We will now argue that, at large enough $`\theta `$, the finite $`\theta `$ saddle point $`\varphi (x,y)`$ of (2.1) that reduces to $`\lambda \|00\|`$ as $`\theta \mathrm{}`$ is classically stable. Consider the mass matrix for fluctuations about $`\varphi (x,y)`$. Since any operator may be written as $`UDU^{}`$ where $`D`$ is diagonal and $`U`$ unitary, small fluctuations may be decomposed into radial (fluctuations of $`D`$) and angular ones (fluctuations of $`U`$). The mass matrix for purely radial fluctuations is $`𝒪(\theta )`$ to leading order, and has been shown to be positive definite in sec 3.4. The mass matrix for purely angular fluctuations is $`𝒪(1)`$ to leading order, and has been shown to be positive definite, modulo the two zero modes. Since angular modes completely disappear from the potential, mixing between radial and angular fluctuations occurs only through the kinetic term, and are also $`𝒪(1)`$. These cross terms result in corrections to the eigenvalues of the mass matrix only at $`𝒪(\frac{1}{\theta })`$. Hence, to leading order in $`\frac{1}{\theta }`$, the mass matrix is positive. The two zero modes of the angular mass matrix cannot be driven negative by $`\frac{1}{\theta }`$ corrections as they are exact. A similar argument demonstrates the instability of all other radially symmetric level one solitons (those that reduce to $`\lambda \|nn\|`$ at $`\theta =\mathrm{}`$) at large enough $`\theta `$. The considerations of this subsection may easily be generalized to solitons in $`2l`$ spatial dimensions, using the higher dimensional analogue of (5.1): $$\frac{g^2K(U)}{(2\pi )^l\lambda ^2}=1+2\underset{j,\stackrel{}{k}}{}k_j\|U_{\stackrel{}{k},\stackrel{}{0}}\|^22\underset{j}{}\|\underset{\stackrel{}{k}}{}\sqrt{k_j+1}U_{\stackrel{}{k},\stackrel{}{0}}U_{\stackrel{}{k}+\stackrel{}{i},\stackrel{}{0}}^{}\|^2$$ and (5.1) $$\frac{g^2K(U)}{(2\pi )^l\lambda ^2}=1+2\underset{\stackrel{}{k}}{}\left(\underset{j=1}{\overset{l}{}}k_j\underset{j=1}{\overset{l}{}}\delta _{\stackrel{}{k},\stackrel{}{i}}\right)\|U_{\stackrel{}{k},\stackrel{}{0}}\|^2.$$ We use the notation of (3.1); $`\stackrel{}{k}`$ is an $`l`$ dimensional vector, $`j`$ runs from $`1`$ to $`l`$ and $`\stackrel{}{i}`$ is the basis unit vector in the $`i^{th}`$ direction; in components $`i_n=\delta _{i,n}`$. Notice that $`K(U)`$ in (5.1) is independent of $`U_{\stackrel{}{i},0}`$ for all $`i`$, a consequence of the exact translational invariance in all $`2l`$ spatial directions. 6. Noncommutative Yang-Mills 6.1. Quartic Action for the $`U(1)`$ Theory in Two Dimensions Consider the action $$S=\frac{1}{4g_{YM}^2}d^2z[\overline{\mathrm{\Phi }},\mathrm{\Phi }][\overline{\mathrm{\Phi }},\mathrm{\Phi }],$$ where $`\mathrm{\Phi }`$ is a complex field and $$[\mathrm{\Phi },\overline{\mathrm{\Phi }}]\mathrm{\Phi }\overline{\mathrm{\Phi }}\overline{\mathrm{\Phi }}\mathrm{\Phi }.$$ The equation of motion following from (6.1) is $$[\overline{\mathrm{\Phi }},[\mathrm{\Phi },\overline{\mathrm{\Phi }}]]=0.$$ $`\mathrm{\Phi }`$ can also be viewed as a quantum mechanical operator and $`\overline{\mathrm{\Phi }}`$ as it’s hermitian conjugate. The commutators in (6.1)-(6.1) are then ordinary operator commutators, and the integral is the trace over the Hilbert space. In the operator representation a simple solution of the equation of motion (6.1) is $$\mathrm{\Phi }=a,\overline{\mathrm{\Phi }}=a^{}.$$ Let us expand around this solution by defining $$\mathrm{\Phi }=a+iA_{\overline{z}},\overline{\mathrm{\Phi }}=a^{}iA_z.$$ One then finds, translating back to functions (with $`\sqrt{2}z=q+ip`$, $`[a,]=_{\overline{z}}`$ and $`[a^{},]=_z`$, that $$[\mathrm{\Phi },\overline{\mathrm{\Phi }}]=1+i_zA_{\overline{z}}i_{\overline{z}}A_z[A_z,A_{\overline{z}}]=1+iF_{z\overline{z}}.$$ The operator representation of (6.1) has the manifest $`U(N=\mathrm{})`$ symmetry under which $`\mathrm{\Phi }\mathrm{\Phi }^{}=U^{}\mathrm{\Phi }U`$ just as in the scalar field theory. Infinitesimally, $$\delta \mathrm{\Phi }=i[\mathrm{\Phi },\mathrm{\Lambda }],$$ where $`U=\mathrm{exp}i\mathrm{\Lambda }`$. When gauged, this is just the usual $`U(1)`$ gauge symmetry of the non-commutative theory, $$\delta A=d\mathrm{\Lambda }+i[A,\mathrm{\Lambda }].$$ The equation of motion (6.1) is $$D_{\overline{z}}F_{z\overline{z}}=0.$$ The action (6.1) is then $$S=\frac{1}{4g_{YM}^2}d^2z\left(F_{z\overline{z}}i\right)^2,$$ the standard two dimensional non-commutative $`U(1)`$ Yang-Mills action up to constants and topological terms. 6.2. The U(N) Theory in $`2l`$ Dimensions (6.1) can be generalized to $$S=\frac{1}{4g_{YM}^2}d^{2l}x\delta ^{\mu \lambda }\delta ^{\nu \rho }\mathrm{Tr}\left([\mathrm{\Phi }_\mu ,\mathrm{\Phi }_\nu ][\mathrm{\Phi }_\lambda ,\mathrm{\Phi }_\rho ]\right),$$ where $`\mu ,\nu =1,,,2l`$ and $`\mathrm{\Phi }_\mu `$ are real $`N\times N`$ matrices. Though we have restricted ourselves to a flat euclidean metric, one can generalise the argument below to the Minkowski metric as well. The equation of motion is $$\delta ^{\mu \nu }[\mathrm{\Phi }_\mu ,[\mathrm{\Phi }_\nu ,\mathrm{\Phi }_\lambda ]]=0.$$ We choose complex coordinates such that $`\mathrm{\Theta }_{a\overline{b}}=i\delta _{a\overline{b}},`$ with $`a,b=1\mathrm{}l`$. (6.1) has the solution $$\mathrm{\Phi }_b=a_b,\mathrm{\Phi }_{\overline{b}}=a_{\overline{b}}^{},$$ where $`[a_b,a_{\overline{c}}^{}]=\delta _{b\overline{c}}.`$ Expanding around this solution with $$\mathrm{\Phi }_b=a_b+iA_{\overline{b}}$$ one finds $$S=\frac{1}{4g_{YM}^2}d^{2l}z\left(F_{a\overline{b}}\mathrm{\Theta }_{a\overline{b}}^1\right)^2.$$ As before the manifest $`U(\mathrm{})U(N)`$ symmetry corresponds to the non-commutative $`U(N)`$ gauge symmetry. 6.3. The U(1) Instanton The four dimensional non-commutative gauge theory has instanton solutions which are deformed versions of the usual non-abelian instantons. In particular, the $`U(1)`$ non-commutative theory also has non-singular finite action saddle points . We exhibit the operators $`\mathrm{\Phi }_a`$ corresponding to the simplest such $`U(1)`$ instanton. The operators $`\mathrm{\Phi }_a`$ corresponding to an anti self dual field strength $`\delta ^{a\overline{b}}F_{a\overline{b}}=0`$ $`(a,\overline{b}=1,2)`$, obey $$[\mathrm{\Phi }_b,\mathrm{\Phi }_c]=0,\delta ^{a\overline{b}}[\mathrm{\Phi }_a,\mathrm{\Phi }_{\overline{b}}]=2.$$ In four dimensions, the operators $`\mathrm{\Phi }_a(a=1,2)`$ live in a Hilbert space generated by the creation and annihilation operators of a two-dimensional harmonic oscillator (See Sec. 3.3). Rather than work in the conventional number basis $`\|n_1,n_2`$, it is convenient to work in Schwinger’s angular momentum basis, $$\|j,m\frac{(a_1^{})^{j+m}}{\sqrt{(j+m)!}}\frac{(a_2^{})^{jm}}{\sqrt{(jm)!}}\|0,0,$$ with $`0j<\mathrm{},\|m\|j`$. The operators $$J_+=a_1^{}a_2,J_{}=a_2^{}a_1,J_z=\frac{1}{2}(a_1^{}a_1a_2^{}a_2)$$ obey the usual angular momentum algebra. We will find a solution to (6.1) of the form $$\begin{array}{cc}\hfill \mathrm{\Phi }_b& =a_b\underset{j,m}{}(1+c_j)\|j,mj,m\|=a_b+a_b\underset{j,m}{}c_j\|j,mj,m\|,\hfill \\ \hfill \mathrm{\Phi }_{\overline{b}}& =a_{\overline{b}}^{},\hfill \end{array}$$ and put it into Hermitian form via a complexified gauge transformation $`W`$. The ansatz (6.1) satisfies the holomorphic part of (6.1) for any $`c_j`$. For a real $`c_j`$, the only condition comes from the equation $`F_{1\overline{1}}=F_{2\overline{2}}`$. Using $$\begin{array}{cc}\hfill a_{1,2}^{}\|j,m& =\sqrt{j\pm m+1}\|j+\frac{1}{2},m\pm \frac{1}{2};\hfill \\ \hfill a_{1,2}\|j,m& =\sqrt{j\pm m}\|j\frac{1}{2},m\frac{1}{2},\hfill \end{array}$$ yields the equation $`jc_j=(j+1)c_{j+\frac{1}{2}}`$. Which has the solution $$c_j=\frac{c}{j(2j+1)},(j>0).$$ The complexified gauge transformation $$W=W^{}=\underset{j,m}{}\sqrt{\frac{j}{j+1}}\|j,m$$ puts the solution (6.1) into Hermitian form for $`c=1`$. The field strength then takes the compact form $$[\mathrm{\Phi }_{\overline{b}},\mathrm{\Phi }_c]=\delta _{\overline{b}c}iF_{\overline{b}c}=\delta _{\overline{b}c}(\stackrel{}{J}\stackrel{}{\sigma })_{\overline{b}c}\underset{j,m}{}\frac{1}{j(j+1)(2j+1)}\|j,mj,m\|.$$ Here $`\stackrel{}{J}`$ are the angular momentum generators defined in (6.1) and $`\stackrel{}{\sigma }`$, the usual Pauli matrices. This solution is exactly the same as the simplest charge one $`U(1)`$ instanton in . It may be checked that $`\frac{1}{2}TrF_{a\overline{b}}^2=1`$. Acknowledgements We are grateful to S. Coleman, J. Harvey, T. Kinoshita, J. Maldacena, B. Pioline, N. Seiberg, I. Singer, C. Vafa, M. Van Raamsdonk, S. Sinha, A. Vishwanath and S.-T. Yau for useful discussions. This work was supported in part by DOE grant DE-FG02-91ER40654. Appendix A. Solutions at Finite $`\theta `$ In this appendix we will examine radially symmetric saddle points of (2.1) at finite $`\theta `$. In subsection A.1 we study the equation of motion resulting from (2.1) at finite $`\theta `$, and examine the existence of radially symmetric solutions to these equations. In A.2 we concentrate on a particular solution; the one that reduces to the stable soliton $`\lambda \|00\|`$ as $`\theta `$ is taken to infinity. We present an approximate construction of this soliton at large $`\theta `$. In A.3 we briefly comment on the generalization of these results to solitons in higher dimensions. A.1. The Perturbation Expansion and a Recursion Relation The full equation of motion derived from (2.1) may be written in momentum space as $$\stackrel{~}{\varphi }(k^2)+\underset{j=3}{\overset{r}{}}\frac{b_j}{m^2}\stackrel{~}{\varphi }^{j1}(k^2)=\frac{k^2}{m^2\theta }\stackrel{~}{\varphi }(k^2)$$ While the LHS of (A.1) is independent of $`\theta `$, the RHS is of order $`\frac{1}{\theta }`$, and so is a small parameter at large $`\theta `$. For notational convenience, we set $`\frac{b_j}{m^2}=d_j`$ and $`\frac{1}{m^2\theta }=ϵ`$. Let $$\underset{n=0}{\overset{\mathrm{}}{}}c_n\stackrel{~}{\varphi }_n(k^2)$$ be a solution to (A.1). Substituting (A.1) into (A.1), using the recurrence relation for Laguerre polynomials, and equating coefficients of $`\stackrel{~}{\varphi }_n(k^2)`$, we arrive at the difference equations $$c_n+\underset{j=3}{\overset{r}{}}d_jc_n^{j1}=2ϵ[nc_{n1}(2n+1)c_n+(n+1)c_{n+1}].$$ We are interested in finite energy solutions to (3.1), i.e. solutions to (A.1) for which $$\underset{n}{}V(c_n)<\mathrm{}.$$ Since $`V(0)=0`$, (A.1) will be satisfied if the $`c_n`$s approach zero sufficiently fast as $`n`$ approaches infinity. For such a solution, all nonlinear terms in (A.1) may be neglected at large enough $`n`$. At sufficiently large $`n`$, $`n`$ may also be replaced by a continuous variable $`u`$, and (A.1) turns into the second order differential equation $$c(u)=2ϵu\frac{d^2c(u)}{du^2}.$$ (A.1) is the Schroedinger equation for a zero energy state of a particle in a $`\frac{1}{u}`$ potential. $`\sqrt{ϵ}`$ plays the role of Planck’s constant, and at small $`ϵ`$ (A.1) is easily solved in the WKB approximation, yielding $$c(u)=A_{}u^{\frac{1}{4}}e^{\sqrt{\frac{2u}{ϵ}}}+A_+u^{\frac{1}{4}}e^{+\sqrt{\frac{2u}{ϵ}}}$$ where $`A_\pm `$ are arbitrary constants. In order that $`c_n`$ tend to zero at large $`n`$, $`A_+=0`$. Thus, for large<sup>8</sup> (A.1) is a good approximation when $`\|c_n\|1`$ (so that dropping nonlinear terms in (A.1) is justified) and $`\frac{c_nc_{n1}}{c_n}1`$, i.e. $`nϵ1`$ (so that the transition from (A.1) to (A.1) is justified). $`n`$, $$c_nAn^{\frac{1}{4}}e^{\sqrt{\frac{2n}{ϵ}}}.$$ (A.1) has an undetermined parameter $`A`$, the scale of the solution at large $`n`$. As (A.1) is a nonlinear equation, $`A`$ is not an arbitrary parameter, but is determined to be one of a discrete set of values. Given $`c_p`$ and $`c_{p+1}`$, the $`(p+1)`$ equations (A.1) with $`n=0\mathrm{}p`$ overdetermine the $`p`$ unknowns $`c_n`$ for $`n<p`$. The extra equation constrains the scale $`A`$, as we’ll see in the next subsection. A.2. The Gaussian Soliton Corrected In this section we present an approximate construction of the stable soliton that reduces to the gaussian at infinite $`\theta `$. Our construction approximates the true solution to arbitrary accuracy at small enough $`ϵ`$. We wish to find a solution of (A.1) such that $$\underset{ϵ0}{lim}c_0=\lambda $$ and $$\underset{ϵ0}{lim}c_m=0$$ uniformly in $`m`$, for $`m1`$. (A.1) ensures that, on such a solution, (A.1) for $`n1`$ reduces to $$c_n=2ϵ[nc_{n1}(2n+1)c_n+(n+1)c_{n+1}]$$ for small enough $`ϵ`$. It is easy to find an explicit solution to (A.1) that obeys (A.1), (A.1). Consider a function $`\varphi (x,y)`$ that obeys the differential equation $$(ϵ^2+1)\varphi =b\varphi _0.$$ Expanding $`\varphi `$ in the form $$\varphi =\underset{n=0}{\overset{\mathrm{}}{}}c_n\varphi _n$$ and imitating the manipulations of section 4.3, we find that $`c_ns`$ obey (A.1) for $`n1`$, but obey $$c_0=2ϵ\left[c_1c_0\right]+b$$ instead of (A.1) (with $`n=0`$). This relation will fix the free parameter $`b`$. (A.1) is easily solved in momentum space $$\stackrel{~}{\varphi }(k)=b\frac{\stackrel{~}{\varphi }_0(k)}{1+2ϵk^2}.$$ Using the explicit forms for $`\stackrel{~}{\varphi }_n(k)`$ and orthogonality of the Laguerre polynomials we find $$c_n=b_0^{\mathrm{}}𝑑x\frac{e^xL_n(x)}{1+2ϵx}.$$ In particular $$c_0=b𝑑x\frac{e^x}{1+2ϵx}=bF(ϵ)\mathrm{where}F(ϵ)=1ϵ+𝒪(ϵ^2).$$ Using (A.1) we conclude that (A.1) and (A.1) (at $`n=0`$) are identical on $`\{c_n\}`$ if $`b`$ is chosen such that $$bF(ϵ)+\underset{j=3}{\overset{r}{}}d_j\left(bF(ϵ)\right)^{j1}=b\left(F(ϵ)1\right).$$ We wish to find a solution to (A.1) that obeys (A.1), i.e. (from (A.1)) one for which $`lim_{ϵ0}b=\lambda `$. As $`\lambda +_{j=3}^rd_j\lambda ^{j1}=0`$, such a solution exists, and takes the form $$b(ϵ)=\lambda (1+Kϵ+𝒪(ϵ^2))$$ at small $`ϵ`$ where $`K`$ is a number that may easily be determined. In summary, $`\{c_n\}`$ given by (A.1) with $`b`$ given by (A.1), (A.1), solve (A.1) for $`n1`$ and (A.1) (with $`n=0`$). $`\{c_n\}`$ therefore also approximately satisfy the true difference equations (A.1) for all $`n`$ as long as $`\|c_n\|1`$ for all $`n1`$. But it is easy to verify that all $`\|c_n\|`$ for all $`n1`$ are arbitrarily small at small enough $`ϵ`$. Using the completeness of the Laguerre polynomials, $$\underset{n=0}{\overset{\mathrm{}}{}}c_n^2=b^2_0^{\mathrm{}}\frac{e^x}{(1+2ϵx)^2}<b^2.$$ But $$c_0=b_0^{\mathrm{}}𝑑x\frac{e^x}{1+2ϵx}>b_0^{\mathrm{}}𝑑xe^x(12ϵx)=b(12ϵ).$$ Combining (A.1) and (A.1) $$\underset{n=1}{\overset{\mathrm{}}{}}c_n^2<4ϵb^2$$ establishing (A.1) uniformly in $`n`$ on our solution. Thus $`\{c_n\}`$ provides an approximate solution to the full nonlinear difference equations (A.1) for all $`n`$ at small enough $`ϵ`$. Furthermore, from (A.1), this solution has finite energy. As $`\{c_n\}`$ obey the linearized recursion relation (A.1) and are small at small $`ϵ`$, we can conclude, from the previous subsection, that $`d_n`$ takes the form (A.1) for $`nϵ1`$. In order to estimate the behaviour of $`c_n(ϵ)`$ for $`n\frac{1}{ϵ}`$ we formally expand the denominator in (A.1) in a power series in $`ϵx`$ and integrate term by term, arriving at the asymptotic expansion $$c_n=\underset{m=n}{\overset{\mathrm{}}{}}(1)^{m+n}(2ϵ)^m\frac{m!^2}{n!(mn)!}.$$ This expansion is useful only when the first few terms in the series in (A.1) are successfully smaller, i.e. for $`nϵ1`$. A.3. Generalization to Higher Dimensions In this subsection we will outline the generalization of the arguments of A.1 and the construction of A.2, for the case of the maximally isotropic noncommutativity in $`2l`$ dimensions, i.e. a theory with noncommutativity matrix $`\mathrm{\Theta }`$, all of whose eigenvalues are $`\pm i\theta `$. It is likely that these arguments can be further extended to generic $`\mathrm{\Theta }`$. We first note that a subset of the diagonal $`\theta =\mathrm{}`$ solutions (3.1) are (in non-dimensionalized coordinates) invariant under $`SO(2l)`$ rotations. These solutions take the form $$\underset{\stackrel{}{n}}{}c_J\frac{1}{\sqrt{D_J}}\delta _{(J,_in_i)}\|\stackrel{}{n}\stackrel{}{n}\|\frac{1}{\sqrt{D_J}}\underset{J}{}c_J\varphi _J^{(l)}(r^2).$$ Here $$\varphi _J^{(l)}(r^2=\underset{i}{}\|z_i\|^2)=2^l(1)^JL_J^{(l1)}(r^2).$$ where $`L_J^{(l1)}(r^2)`$ is an associated Laguerre polynomial. (A.1) is obtained from (3.1) by repeated use of the identity $$\underset{m=0}{\overset{n}{}}L_{nm}^\alpha (x)L_m^\beta (y)=L_n^{\alpha +\beta +1}(x+y).$$ $`D_J=\left(\begin{array}{c}J+l1\\ J\end{array}\right)`$ is a convenient normalization factor. When the noncommutativity matrix is maximally isotropic, the kinetic term in (2.1) is invariant under $`SO(2l)`$ rotations of rescaled coordinates. Thus the corrections to an $`SO(2l)`$ invariant $`\theta =\mathrm{}`$ solution, of the form (A.1), are also $`SO(2l)`$ invariant. Restricting to $`SO(2l)`$ invariant functions, the arguments of section A.1 are easily generalized. Any $`SO(2l)`$ invariant function takes the form $$\stackrel{~}{\varphi }(k^2)=\underset{n=0}{\overset{\mathrm{}}{}}c_J\stackrel{~}{\varphi }_J(k^2);\stackrel{~}{\varphi }_J(k^2)=\frac{1}{\sqrt{D_J}}(2\pi )^lL_J^{(l1)}(\frac{k^2}{2})e^{\frac{k^2}{4}}.$$ The equation of motion implies that $`c_J`$ obey the following generalization of (A.1) $$c_J+\underset{j=3}{\overset{r}{}}d_jc_J^{j1}=2ϵ[(J+l1)c_{J1}(2J+l)c_J+(J+1)c_{J+1}].$$ For large $`J`$ (A.1) and (A.1) are identical, hence all conclusions of section A.1 carry over to this case. 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warning/0003/cond-mat0003334.html	ar5iv	text	# 𝜇⁺-Knight Shift Measurements in U0.965Th0.035Be13 Single Crystals ## Abstract Muon spin rotation ($`\mu `$SR) measurements of the temperature dependence of the $`\mu ^+`$-Knight shift in single crystals of U<sub>0.965</sub>Th<sub>0.035</sub>Be<sub>13</sub> have been used to study the static spin susceptibility $`\chi _s`$ below the transition temperatures $`T_{c1}`$ and $`T_{c2}`$. While an abrupt reduction of $`\chi _s`$ with decreasing temperature is observed below $`T_{c1}`$, $`\chi _s`$ does not change below $`T_{c2}`$ and remains at a value below the normal-state susceptibility $`\chi _n`$. In the normal state we find an anomalous anisotropic temperature dependence of the transferred hyperfine coupling between the $`\mu ^+`$-spin and the U $`5f`$-electrons. An intriguing feature of the heavy-fermion compound U<sub>1-x</sub>Th<sub>x</sub>Be<sub>13</sub> is that for $`0.019x0.045`$ a second phase transition $`T_{c2}(x)`$ appears at a temperature below the superconducting (SC) transition $`T_{c1}(x)`$ . The nature of the lower transition $`T_{c2}`$ is still a matter of considerable debate. Initially $`T_{c2}`$ was identified as a second distinct SC transition from measurements of a specfic heat peak , the pressure dependence of $`T_c(x)`$ and the increased slope of $`H_{c1}`$ vs. $`T`$ . The observation of a $`T^3`$ dependence of the <sup>9</sup>Be NMR spin-lattice relaxation rate below $`T_{c2}`$ suggested a SC state characterized by line nodes in the energy gap . Later zero-field $`\mu `$SR measurements clearly revealed the onset of small moment magnetism ($`10^3`$ $`\mu _B`$/U) below $`T_{c2}`$. The appearance of a small internal field could arise from a SC state that breaks time-reversal symmetry . On the other hand, it could originate from a spin-density wave instability or the formation of long-range AFM correlations within a single SC phase. However, these latter interpretations fail to explain the large specific heat jump at $`T_{c2}`$. Early $`\mu `$SR measurements on polycrystalline samples of U<sub>0.967</sub>Th<sub>0.033</sub>Be<sub>13</sub> showed a constant or perhaps weakly increasing $`\mu ^+`$-Knight shift upon cooling below $`T_{c1}`$ . In the SC state the temperature dependence of the Knight shift $`K`$ reflects the change in the static spin susceptibility $`\chi _s`$ due to the formation of Cooper pairs. For the case of orbital $`s`$-wave ($`L=0`$) spin singlet ($`S=0`$) pairing, Yosida calculated from the BCS theory that $`\chi _s(T)`$ vanishes as $`T0`$ K. Modifications to this temperature dependence are expected for spin-orbit scattering by impurities and unconventional pairing states. In this Letter we report on the temperature dependence of the $`\mu ^+`$-Knight shift in single crystals of U<sub>0.965</sub>Th<sub>0.035</sub>Be<sub>13</sub>. These measurements differ from earlier studies on polycrystalline samples in that there are two magnetically inequivalent $`\mu ^+`$-sites which facilitate a determination of $`\chi _s`$ in the SC state. We find that upon cooling through $`T_{c1}`$, $`\chi _s`$ rapidly decreases, but remains independent of temperature below $`T_{c2}`$. Our study also reveals a temperature dependence in the normal state of the transferred hyperfine coupling at one of the two $`\mu ^+`$-sites, roughly coinciding with features observed in resistivity and specific heat data for pure UBe<sub>13</sub>. The single crystals of U<sub>0.965</sub>Th<sub>0.035</sub>Be<sub>13</sub> were grown from an Al flux as described in Ref. . From zero-field specific heat measurements the upper and lower transitions occur at $`T_{c1}=0.47(5)`$ K and $`T_{c2}=0.35(2)`$ K, respectively. The $`\mu `$SR measurements were carried out using a top loading dilution refrigerator on the M15 beam line at the TRI-University Meson Facility (TRIUMF), Canada and using a <sup>4</sup>He gas-flow cryostat on the $`\pi `$M3 beam line at the Paul Scherrer Institute (PSI), Switzerland. The crystals were mounted on a Ag plate attached to a cold finger. The magnetic field $`𝐇`$ was applied parallel to the crystallographic $`\widehat{𝐜}`$-axis and transverse to the initial $`\mu ^+`$-spin polarization direction. As a local spin-$`1/2`$ probe, the muon is sensitive only to magnetic interactions and precesses about the local magnetic field $`B_\mu `$ with a Larmor frequency $`\omega =\gamma _\mu B_\mu `$, where $`\gamma _\mu /2\pi =13.55342`$ MHz/kOe. The applied field results in a uniform polarization of the localized U $`5f`$-moments, which reside at the corners of a cubic lattice. The Fourier transform of the $`\mu ^+`$-spin precession signal in U<sub>0.965</sub>Th<sub>0.035</sub>Be<sub>13</sub> shows two distinct symmetric lines with an amplitude ratio of 1:2. In the time domain, each signal was best fit by a Gaussian relaxation function $`G(t)=\mathrm{exp}(\sigma ^2t^2/2)`$, where $`\sigma `$ is the $`\mu ^+`$-spin depolarization rate. From the amplitude ratio and the frequencies of these two signals, we have determined that the $`\mu ^+`$ stops at the (0, 0, 1/4) site, half way between nearest-neighbor U atoms. Muons stopping between U atoms adjoined along the $`\widehat{𝐜}`$-axis direction experience a net dipolar field from the $`5f`$-moments which is parallel to $`𝐇`$, and thus precess at a frequency $`\omega _{}`$ that is greater than those stopping in Ag (which provide a zero-shift reference frequency). On the other hand, twice as many $`\mu ^+`$ stop between U atoms adjoined along the $`\widehat{𝐚}`$\- and $`\widehat{𝐛}`$-axis directions, where the net dipolar field is antiparallel to $`𝐇`$. These muons precess at a frequency $`\omega _{}`$ that is lower than those stopping in Ag. The Knight shift at the two magnetically inequivalent $`\mu ^+`$-sites is given by $$K_,=\left(\omega _,\omega _{\mathrm{Ag}}\right)/\omega _{\mathrm{Ag}}.$$ (1) Figure 1 shows measurements of the temperature dependence of $`K_{}`$ and $`K_{}`$ below 30 K at $`H=10`$ kOe and above 2 K at 6 kOe (insets). The reduction of $`K_{}`$ above $`T50`$ K is attributed to crystal electric field (CEF) excitations, which have been inferred from specific heat and NMR spin-lattice relaxation studies in pure UBe<sub>13</sub>. The effect on the hyperfine coupling is observable for both $`\mu ^+`$-sites from plots of $`K`$ vs. $`\chi _{\mathrm{mol}}`$ in the normal state (see Fig. 2), where $`\chi _{\mathrm{mol}}`$ is the isotropic bulk molar susceptibility. The plots are essentially linear between 5 and 50 K (where $`K`$ follows a Curie-Weiss behavior) and at temperatures above 63 K, with a change of slope between the two regions. The temperature dependence of $`\chi _{\mathrm{mol}}^1`$ is shown in the inset of Fig. 2 compared with that for proposed CEF spittings of U<sup>4+</sup> $`J=4`$ and U<sup>3+</sup> $`J=9/2`$ manifolds in cubic symmetry. The CEF models have been corrected by adding a molecular-field constant of 57 emu/mol, compared to 52 emu/mol in CeCu<sub>2</sub>Si<sub>2</sub> . Although the data are consistent with $`J=9/2`$, the $`J=4`$ energy scheme results in similar behavior when the hybridization proposed in the quadrupolar Kondo model is included in the calculation of $`\chi _{\mathrm{mol}}^1(T)`$ — as was shown for the case of pure UBe<sub>13</sub> . A linear fit to $`\chi _{\mathrm{mol}}^1(T)`$ above 100 K yields an effective moment of 3.62(1) $`\mu _B`$/U. Equation (1) can be expressed in terms of the individual contributions to $`K`$, so that for the axial symmetry of the $`\mu ^+`$-site $$K_{}=(A_\mathrm{c}^{}+A_{\mathrm{dip}}^{zz})\chi _{5f}+K_{\mathrm{dem},\mathrm{L}}+K_0+K_{\mathrm{dia}}$$ (2) and $$K_{}=(A_\mathrm{c}^{}\frac{1}{2}A_{\mathrm{dip}}^{zz})\chi _{5f}+K_{\mathrm{dem},\mathrm{L}}+K_0+K_{\mathrm{dia}},$$ (3) where $`A_\mathrm{c}`$ and $`A_{\mathrm{dip}}^{zz}`$ are the contact hyperfine and dipolar coupling constants pertaining to the interaction of the $`\mu ^+`$ with the $`5f`$-electrons (i.e. $`AH_{hf}/N_A\mu _B`$, where $`H_{hf}`$ is the hyperfine field, $`N_\mathrm{A}`$ is Avogadro’s number and $`\mu _B`$ is the Bohr magneton), $`\chi _{5f}`$ is the isotropic molar $`5f`$-electron susceptibility, $`K_{\mathrm{dem},\mathrm{L}}=4\pi (1/3N)\rho _{\mathrm{mol}}\chi _{\mathrm{mol}}`$ is the correction for the demagnetization and Lorentz fields (where $`N1`$ is the demagnetization factor and $`\rho _{\mathrm{mol}}=0.01227`$ mol/cm<sup>3</sup> is the molar density), $`K_0`$ is the isotropic $`T`$-independent contribution from the non-$`5f`$ conduction electrons, and $`K_{\mathrm{dia}}`$ is due to flux expulsion in the SC state. The total normal-state susceptibility is given by $`\chi _{\mathrm{mol}}=\chi _{5f}+\chi _0`$, where $`\chi _0`$ is the $`T`$-independent non-$`5f`$ contribution. From the normal-state plot of $`K_{}K_{}`$ vs. $`\chi _{\mathrm{mol}}`$ at 10 kOe (see Fig. 3), $`\chi _0=0.0039(2)`$ emu/mol was obtained from the intercept of the extrapolated linear line, where $`\chi _{5f}1/T0`$ and $`K_{}=K_{}K_0=1846(90)`$ ppm. In general, $`A_\mathrm{c}^{}=A_\mathrm{c}^{}`$, in which case the slope of the solid line (3/2)$`A_{\mathrm{dip}}^{zz}`$ gives $`A_{\mathrm{dip}}^{zz}=2066(22)`$ Oe/$`\mu _B`$. This value agrees with the result $`A_{\mathrm{dip}}^{zz}=2062`$ Oe/$`\mu _B`$ obtained from a simple dipolar-field calculation for U moments sitting on the corners of a cubic lattice of edge-length 5.134 Å. We note that the value of $`A_{\mathrm{dip}}^{zz}`$ obtained from the 6 kOe data is only 1693(28) Oe/$`\mu _B`$. Although this may imply that $`A_\mathrm{c}`$ is anisotropic, the time spectra recorded at 6 kOe had a larger time resolution and far fewer muon-decay events than the spectra taken at 10 kOe. Thus, it is likely that there is a systematic uncertainty in the temperature dependence of the $`\mu ^+`$-Knight shift at 6 kOe. $`A_\mathrm{c}`$ represents the transferred hyperfine coupling between the $`\mu ^+`$-spin and the U $`5f`$-electrons via the conduction $`s`$-electrons. Substituting the value of $`A_{\mathrm{dip}}^{zz}`$ into Eqs. (2) and (3) gives the temperature dependence of $`A_\mathrm{c}^{}`$ and $`A_\mathrm{c}^{}`$ shown in Fig. 4. The offset of the 6 kOe data stems from the discussion in the previous paragraph. The decrease above 50 K is likely due to the mixing of the wave functions associated with the different CEF levels. The strong reduction of $`A_\mathrm{c}^{}`$ and lack of change of $`A_\mathrm{c}^{}`$ for $`T_{c1}T4`$ K is the source of the nonlinearity above $`\chi _{\mathrm{mol}}0.014`$ emu/mol in Fig. 3. This is not a muon induced effect, since similar departures from linearity have been observed in $`K`$-$`\chi _{\mathrm{mol}}`$ plots for the <sup>9</sup>Be NMR Knight shift in UBe<sub>13</sub> and the <sup>63</sup>Cu and <sup>29</sup>Si NMR Knight shifts in CeCu<sub>2</sub>Si<sub>2</sub> . A decrease of $`A_\mathrm{c}^{}`$ over nearly the same temperature range is also observed in pure UBe<sub>13</sub> . This anomaly roughly coincides with the peak in the resistivity and specific heat at $`2.5`$ K in UBe<sub>13</sub> , which is completely suppressed when 3.55 % Th is added. The decrease of $`A_\mathrm{c}^{}`$ in both the pure and doped systems is not necessarily inconsistent with this latter behavior, because most of the $`\mu ^+`$ stopping in U<sub>0.965</sub>Th<sub>0.035</sub>Be<sub>13</sub> do not reside near a Th impurity. In the SC state the flux expulsion term $`K_{\mathrm{dia}}`$ in Eqs. (2) and (3) is related to the value of the magnetic penetration depth $`\lambda `$ and the coherence length $`\xi _0`$. To our knowledge there have been no measurements of the absolute value of $`\lambda `$ in U<sub>0.965</sub>Th<sub>0.035</sub>Be<sub>13</sub>. However, the lack of any increase in the $`\mu ^+`$-spin depolarization rate $`\sigma `$ below $`T_{c1}`$ is consistent with a value $`\lambda (0)12100`$ Å, as reported in pure UBe<sub>13</sub> . Using the simple theoretical model developed by Hao et al. for the reversible magnetization of a type-II superconductor and the value $`H_{c2}(0)55`$ kOe , we calculate that $`\|K_{\mathrm{dia}}\|72`$ ppm. Since we observed no field dependence for $`K_,`$ below $`T_{c1}`$ in the range 5 kOe $`H15`$ kOe, we conclude that the internal field is essentially uniform and the diamagnetic shift $`K_{\mathrm{dia}}`$ is negligible. Because $`A_\mathrm{c}^{}`$ is temperature independent in the normal state, we make the reasonable assumption that it remains so below $`T_{c1}`$, allowing $`\chi _s`$ (i.e., $`\chi _{5f}`$ in the SC state) to be determined from Eq. (2). As shown in Fig. 5, $`\chi _s(T)`$ exhibits two different behaviors (in agreement with the raw Knight shift data in Fig. 1) which coincide with the two phase transitions in the specific heat. The decrease of $`\chi _s(T)`$ between $`T_{c1}`$ and $`T_{c2}`$ is consistent with a phase in which the Cooper pairs have a substate of opposite spin projection (i.e., $`S_z=0`$). However, the data cannot distinguish between even and odd parity spin states possessing this substate, because Fermi-liquid corrections and spin-orbit (SO) scattering by impurities may be significant. For the case of an even parity SC phase we can estimate the importance of SO scattering from the relation $`\chi _s(T_{c2})/\chi _n=12l_{\mathrm{SO}}/\pi \xi _0`$ , where $`\chi _n`$ is the normal-state spin susceptibility at $`T_{c1}`$, $`\chi _s(T_{c2})/\chi _n=0.61`$ and $`\xi _077`$ Å from $`H_{c2}(0)`$ . This gives a SO scattering mean free path of $`l_{\mathrm{SO}}47`$ Å. The average distance between Th atoms $`15`$ Å, represents a lower limit for the mean free path $`l`$ between collisions of the electrons with the Th impurities. Since $`l_{\mathrm{SO}}`$ is of the same order of $`l`$, modification of $`\chi _s(T)`$ due to SO scattering cannot be ruled out. The lack of a temperature dependence for $`\chi _s`$ below $`T_{c2}`$ is characteristic of a spin-triplet ($`S=1`$) odd-parity ($`L=1`$) superconductor with parallel spin pairing, except that $`\chi _s<\chi _n`$. This unusual behavior suggests that the component of the order parameter corresponding to the phase $`T_{c2}<T<T_{c1}`$ stops or slows down its growth at $`T_{c2}`$, where a second component develops. In terms of the $`𝐝`$-vector of the triplet order parameter $`\widehat{\mathrm{\Delta }}(𝐤)=i(𝐝\sigma )\sigma _y`$, a possibile scenario is that (i) one component corresponds to $`𝐝𝐇`$, so that $`\chi _s`$ decreases below $`T_{c1}`$ and (ii) the second component corresponds to $`𝐝𝐇`$, in which case $`\chi _s`$ is unchanged below $`T_{c2}`$. The idea of a two-component $`𝐝`$-vector is similar to the weak spin-orbit coupling model recently developed for UPt<sub>3</sub> from detailed <sup>195</sup>Pt NMR Knight shift measurements . Finally, substituting $`\chi _s(T)`$ for $`\chi _{5f}`$ in Eq. (3) we find that the magnitude of $`A_\mathrm{c}^{}`$ rapidly increases to a constant value below $`T_{c2}`$ (see Fig 4). In conclusion, our study of U<sub>0.965</sub>Th<sub>0.035</sub>Be<sub>13</sub> has identified different behavior for the temperature dependence of the spin susceptibility $`\chi _s`$ below the two transitions observed in the specific heat. A possible explanation for the absence of a change below $`T_{c2}`$ is that U<sub>0.965</sub>Th<sub>0.035</sub>Be<sub>13</sub> is an odd parity spin-triplet superconductor. However, we stress that this may not be the only interpretation of our measurements. A definitive identification of the pairing state will require further measurements as a function of magnetic-field direction to unambiguously determine the relative orientation of the $`𝐝`$-vector and $`𝐇`$. We are extremely grateful to K. Machida, D.L. Cox, M.J. Graf, A.V. Balatsky and A. Schenck for informative discussions. Work at Los Alamos was performed under the auspices of the U.S. DOE. Other support was provided by U.S. NSF, Grant DMR-9731361 (U.C. Riverside), DMR-9820631 (California State) and the Dutch Foundations FOM and NWO (Leiden). FIGURE CAPTIONS Figure 1. Temperature dependence of $`K_{}`$ and $`K_{}`$ measured at TRIUMF in an applied field $`H=10`$ kOe. Inset: Data taken above $`T_{c1}`$ at PSI, where the maximum available field was $`H=6`$ kOe. Figure 2. Plot of the normal-state $`\mu ^+`$-Knight shift at $`H=6`$ kOe vs. the bulk molar susceptibility. Inset: Temperature dependence of the inverse susceptibility. Dashed and solid lines are calculations using Eq. (3) of Ref. for the CEF schemes described in the text. Figure 3. Plot of $`K_{}K_{}`$ vs. the bulk molar susceptibility at $`T>T_{c1}`$ and $`H=10`$ kOe. The solid line is a linear fit to the data above 5 K. Figure 4. Temperature dependence of $`A_\mathrm{c}^{}`$ (open symbols) and $`A_\mathrm{c}^{}`$ (solid symbols) at $`H=6`$ kOe and 10 kOe. Note: We have assumed that $`A_\mathrm{c}^{}`$ is unchanged below $`T_{c1}`$. The data for $`A_\mathrm{c}^{}`$ below $`T_{c1}`$ was obtained under this assumption. Figure 5. Temperature dependence of the specific heat (open circles) and magnetic susceptibility (solid circles).
warning/0003/math0003176.html	ar5iv	text	# Bordism-finiteness and semi-simple group actions ## 1 Introduction A direct consequence of the Atiyah-Bott-Segal-Singer fixed point theorem in equivariant index theory is the fact that all the Pontrjagin numbers of an oriented closed manifold with smooth fixed point free $`S^1`$-action vanish. Since the oriented bordism ring is determined by Pontrjagin and Stiefel-Whitney numbers this result may be interpreted as a bordism-finiteness theorem: In a given dimension the class of oriented closed manifolds which support a fixed point free $`S^1`$-action contains only finitely many oriented bordism types. In this paper we give bordism-finiteness results for smooth actions on connected closed oriented manifolds with possible non-empty fixed point set. We also look at the related question of diffeomorphism-finiteness in a given homotopy type. In Riemannian geometry finiteness theorems (involving bounds for curvature, diameter, volume etc.) go back to the work of Cheeger (cf. \[Ch\]) and are the topic of active research since then. It is natural to look what happens if the geometrical bounds arise from symmetries, i.e. group actions. Concerning the questions on bordism- and diffeomorphism-finiteness it is desirable to obtain positive results for groups as “small” as possible. Our results involve smooth actions by the semi-simple group $`S^3`$. Before we state these we take a brief look at actions by “smaller” groups. For finite groups one cannot expect any general finiteness results, even if one restricts to free actions or to a fixed homotopy type. Dovermann and Masuda used equivariant surgery to construct infinitely many homotopy $`P^3`$’s with effective smooth $`_p`$-action for every prime number $`p`$ (cf. \[DoMa\] and the references therein). Löffler and Raußen (cf. \[LoRa\]) used Sullivan’s theory of minimal models to construct free $`_p`$-actions on certain high connected manifolds for any large enough prime number $`p`$ and conjectured that non-trivial $`_p`$-actions exist on any simply-connected manifold (cf. also \[We\]; for the question of bordism-finiteness cf. \[CoFl\] and \[Di\]). On the other hand, by a result of Schultz, one knows that in dimension $`4`$ any oriented bordism class contains infinitely many manifolds with infinite fundamental group which admit no effective finite group action (cf. \[Sc\]). Next we discuss the question of bordism-finiteness for circle actions. By taking products or equivariant connected sum along invariant submanifolds one can construct new $`S^1`$-manifolds from given pieces. Using these methods it is easy to show that for any given oriented manifold a suitable multiple (disjoint union) is bordant to a connected oriented manifold with non-trivial $`S^1`$-action. On the other hand the Lefschetz fixed point formula in equivariant index theory (cf. \[AtSe\], \[AtSiIII\], \[AtBo\], \[Bo\]) implies that the bordism class of a manifold with $`S^1`$-action is modulo torsion completely determined by the local geometry of the action near the fixed point set (see Cor. 3.2). In order to get interesting bordism-finiteness results one needs additional conditions which limit the constructions mentioned above but are not too restrictive on the local geometry of the action. This leads us to consider $`S^1`$-actions with isolated fixed points which satisfy a prescribed upper bound for the number of fixed points (by the classical Lefschetz fixed point formula the number of $`S^1`$-fixed points is just the Euler characteristic). In this way we exclude products with trivial $`S^1`$-manifolds and limit the use of the connected sum construction but still allow various local geometries although the non-equivariant part of these is trivial. For these $`S^1`$-actions bordism-finiteness still fails (see the examples in Section 4). Taking products one concludes the same for torus actions as long as the dimension of the torus is small compared to the dimension of the manifold. We note that the absence of strong implications on the bordism type for torus actions is underlined by the following result of Buchstaber and Ray (cf. \[BuRa\]): The complex bordism ring is generated by toric manifolds. In particular, any $`2n`$-dimensional stable almost complex manifold is bordant to a manifold which admits an action by the $`n`$-dimensional torus. Since bordism-finiteness fails for abelian actions it is natural to look at semi-simple group actions next. Our results involve actions of $`S^3`$ with the following property: $$S^3\text{ acts with fixed point and the }S^1\text{-action has isolated fixed points}$$ ($``$) for some (and hence every) fixed subgroup $`S^1S^3`$. In turns out that for such actions the question whether bordism-finiteness holds depends on properties of the first Pontrjagin class. A special case is the following ###### Theorem 1.1. Let $`C`$ and $`m`$ be natural numbers. The class of connected $`m`$-dimensional oriented manifolds with vanishing first Pontrjagin class and Euler characteristic $`C`$ which admit an $`S^3`$-action satisfying $`()`$ contains only finitely many oriented bordism types. The conclusion holds in the more general situation where a negative multiple of the first Pontrjagin class $`p_1`$ is a sum of squares (see Th. 3.4). Moreover bordism-finiteness holds if one refines the bordism ring taking into account the condition on $`p_1`$ (see Th. 3.6). To prove Theorem 1.1 and its refinements we combine the Lefschetz fixed point formula in equivariant index theory with the following well known properties of $`S^3`$ (which hold for any semi-simple group): In any fixed dimension the group $`S^3`$ admits only finitely many non-equivalent representations. Any one-dimensional $`S^3`$-representation is trivial. These properties imply that the local geometry of the induced $`S^1`$-action at the $`S^3`$-fixed point is determined up to finite ambiguity. The condition on the first Pontrjagin class guarantees that this is also true at the other $`S^1`$-fixed points. Finally one applies the Lefschetz fixed point formula in equivariant index theory to conclude that the oriented bordism type (resp. its refinement) is determined up to finite ambiguity. Regarding the problem of diffeomorphism-finiteness in a given homotopy type we use similar arguments to show ###### Theorem 1.2. In a fixed dimension $`2n`$ there are only finitely many homotopy complex projective spaces with an $`S^3`$-action satisfying $`()`$. By simply-connected surgery theory one knows that in any dimension $`2n6`$ there are infinitely many homotopy complex projective spaces distinguished by their Pontrjagin classes (cf. \[Hs\]). According to Theorem 1.2 almost all of them do not admit an $`S^3`$-action satisfying $`()`$. It is interesting to compare the result above with a conjecture of Petrie (cf. \[Pe\]) which states that the total Pontrjagin class of a homotopy $`P^n`$ with non-trivial $`S^1`$-action is standard. The conjecture is known to be true if $`n4`$ (cf. \[Ja\]; for related results cf. for example the survey \[Do\] as well as \[De2\] and the references therein). It implies diffeomorphism-finiteness for homotopy complex projective spaces with non-trivial $`S^1`$-action. As indicated before our methods are rather classical, involving index theory and properties of semi-simple groups, and work best if a negative multiple of the first Pontrjagin class is a sum of squares. On the other hand the theory of elliptic genera led to applications for $`Spin^c`$-manifolds with $`S^3`$-action if the first Pontrjagin class is a sum of squares (cf. \[De2\], \[De3\], \[Li\]). Combining both methods one obtains further information on manifolds with $`S^3`$-action. This paper is structured in the following way. Section 2 deals with tangential weights and weights of equivariant complex line bundles for manifolds with admit an $`S^3`$-action with fixed point. We give conditions in terms of the first Pontrjagin class which guarantee that these weights are determined up to finite ambiguity. In Section 3 we extend Theorem 1.1 to the case that a negative multiple of the first Pontrjagin class is a sum of squares (see Theorem 3.4). To prove this theorem and related results we use the Lefschetz fixed point formula in equivariant index theory. In Section 4 we show, by example, that the bordism-finiteness result given in Theorem 3.4 fails if one relaxes the conditions on the $`S^3`$-action or the condition on the first Pontrjagin class. In Section 5 we prove the result on homotopy complex projective spaces (see Theorem 1.2 above). We use similar methods to show that in a given dimension there are only finitely many complete intersections which admit an $`S^3`$-action satisfying $`()`$. ## 2 Weights In this section we give some information on the tangential weights and the weights of complex line bundles over a $`2n`$-dimensional manifold $`M`$ which supports a non-trivial $`S^3`$-action with fixed point. By a result of Hattori and Yoshida (cf. \[HaYo\]) the $`S^3`$-action lifts uniquely to any complex line bundle $`LM`$ and we shall always do so. Fix $`S^1S^3`$. At an $`S^1`$-fixed point $`L`$ reduces to a complex one-dimensional $`S^1`$-representation with character $`\lambda \lambda ^a`$. We call $`a`$ the weight of $`L`$ at this point. Note that at the $`S^3`$-fixed point the weight has to vanish since any complex one-dimensional $`S^3`$-representation is trivial (one way to characterize semi-simple compact Lie groups). At an $`S^1`$-fixed point the tangent bundle $`TM`$ reduces to a real $`S^1`$-representation which we identify (non-canonically) with a complex representation with character $`\lambda \lambda ^{m_i}`$. We call $`m_1,\mathrm{},m_n`$ the tangential weights of $`M`$ at the fixed point. Note that the tangential weights are only well defined up to sign. At an $`S^3`$-fixed point $`pt`$ the $`S^1`$-representation $`TM_{\|pt}`$ is induced from an $`S^3`$-representation. Since there are only finitely many equivalence classes of such representations in a given dimension (another way to characterize semi-simple compact Lie groups) the tangential weights at $`pt`$ are determined by the dimension of $`M`$ up to finite ambiguity, i.e. the set of tangential weights at $`pt`$ belongs to a finite set which only depends on $`2n`$. The next lemma is used in the following section to derive the bordism-finiteness results mentioned in the introduction. ###### Lemma 2.1. Let $`L_1,\mathrm{},L_k`$ be $`S^3`$-equivariant complex line bundles over $`M`$ and $`N`$ a positive integer. Assume the first Pontrjagin class of the bundle $`E:=(L_1+\mathrm{}+L_k)+NTM`$ vanishes. Then the tangential weights at any $`S^1`$-fixed point are determined by the dimension $`2n`$ of $`M`$ up to finite ambiguity. The weights of $`L_i`$ are determined by $`N`$ and $`2n`$ up to finite ambiguity. Proof: Note that the weights of $`E`$ at an $`S^1`$-fixed point are just the weights of the line bundles $`L_j`$ and the tangential weights taken with multiplicity $`N`$. Let $`Y_1,\mathrm{},Y_C^{}`$ denote the connected components of the fixed point manifold $`M^{S^1}`$. We assume that the $`S^3`$-fixed point $`pt`$ belongs to $`Y_1`$. Let $`m_{s,1},\mathrm{},m_{s,n}`$ denote the tangential weights and let $`a_{s,j}`$ denote the weight of $`L_j`$ at $`Y_s`$. Since the $`S^3`$-equivariant vector bundle $`E`$ has vanishing first Pontrjagin class it follows from a spectral sequence argument (see the lemma below) that the weights of $`E`$ at $`Y_s`$ satisfy $$\underset{j=1}{\overset{k}{}}a_{s,j}^2+N\underset{i=1}{\overset{n}{}}m_{s,i}^2=C^{\prime \prime },$$ (1) where $`C^{\prime \prime }`$ is a constant which does not depend on $`Y_s`$. At the $`S^3`$-fixed point $`ptY_1`$ the weights of $`L_j`$ vanish and the tangential weights $`m_{1,i}`$ are determined by the dimension of $`M`$ up to finite ambiguity. Hence $`C^{\prime \prime }=N_{i=1}^nm_{1,i}^2`$ is bounded from above by $`N`$ times a constant which only depends on $`2n`$. By formula (1) the tangential weights (resp. the weights of $`L_j`$) at any $`S^1`$-fixed point are determined by $`2n`$ (resp. $`2n`$ and $`N`$) up to finite ambiguity. $`\mathrm{}`$ ###### Lemma 2.2. Let $`EM`$ be a $`2r`$-dimensional $`S^3`$-equivariant vector bundle with weights $`e_{s,1},\mathrm{},e_{s,r}`$ at the $`S^1`$-fixed point component $`Y_s`$. If $`p_1(E)=0`$ then $`_{i=1}^re_{s,i}^2`$ is independent of $`Y_s`$. Proof: We use the Leray-Serre spectral sequence for the Borel construction of $`EM`$ (for details cf. \[De1\], Lemma A.6). Let $`M_{S^3}:=ES^3\times _{S^3}M`$ and let $`\pi :M_{S^3}BS^3`$ be the fibration with fibre $`i:MM_{S^3}`$ associated to the universal $`S^3`$-principal bundle $`ES^3BS^3`$ and the $`S^3`$-space $`M`$. Since $`H^{}(BS^3;)`$ is concentrated in degree $`4`$ it follows from the Leray-Serre spectral sequence that the sequence $$H^4(BS^3;)\stackrel{\pi ^{}}{}H^4(M_{S^3};)\stackrel{i^{}}{}H^4(M;)$$ (2) is exact. Next consider the bundle $`E_{S^3}M_{S^3}`$ associated to the $`S^3`$-equivariant vector bundle $`EM`$ via the Borel construction. We denote its first Pontrjagin class by $`p_1(E)_{S^3}H^4(M_{S^3};)`$. Since $`i^{}(p_1(E)_{S^3})=p_1(E)=0`$ it follows from (2) that $`p_1(E)_{S^3}=\pi ^{}(\alpha )`$ for some $`\alpha H^4(BS^3;)`$. By naturality, $`p_1(E)_{S^1}=\pi ^{}(C^{\prime \prime }x^2)`$, where $`\pi `$ also denotes $`M_{S^1}BS^1`$, $`xH^2(BS^1;)`$ denotes a generator and $`C^{\prime \prime }x^2`$ is the image of $`\alpha `$ under the map induced by the inclusion $`S^1S^3`$. Note that $`p_1(E)_{S^1}`$ reduces at $`qY_s`$ to $`(_{i=1}^re_{s,i}^2)x^2H^4(BS^1;)H^4(q_{S^1};)`$. Hence $`_{i=1}^re_{s,i}^2=C^{\prime \prime }`$ for any $`Y_s`$. $`\mathrm{}`$ ## 3 Bordism-finiteness In this section we apply the Lefschetz fixed point formula to the result of the previous section to derive bordism-finiteness theorems. Let $``$ be an $`S^1`$-equivariant elliptic operator over $`M`$ with symbol $`\sigma K_{S^1}(TM)`$. By the fundamental work of Atiyah and Singer (cf. \[AtSiI\]) the equivariant index of $``$ is equal to the topological index $`t\text{-}ind(\sigma )`$, where $`t\text{-}ind:K_{S^1}(TM)R(S^1)`$ is defined as the push-forward in complex $`K`$-theory for $`Mpt`$. Let $`N`$ denote the normal bundle of $`i:M^{S^1}M`$ and $`\mathrm{\Lambda }_t:=_{i=0}^{\mathrm{}}\mathrm{\Lambda }^it^i`$. ###### Theorem 3.1 (Lefschetz fixed point formula \[AtSe\]). For an element $`\sigma K_{S^1}(TM)`$ and any topological generator $`\lambda S^1`$ $$t\text{-}ind(\sigma )(\lambda )=t\text{-}ind\left(\frac{i^{}(\sigma )(\lambda )}{\mathrm{\Lambda }_1(N_{})(\lambda )}\right).$$ (3) $`\mathrm{}`$ Note that the right hand side of formula (3) consists of a finite sum of local contributions at the fixed point components of the $`S^1`$-action. Since the set of topological generators is dense in $`S^1`$ the theorem above implies that the index function $`t\text{-}ind:K_{S^1}(TM)R(S^1)`$ vanishes identically if $`M^{S^1}`$ is empty. If the elliptic operator is geometrically defined using an $`H`$-structure (cf. \[AtSiIII\]), for example if $``$ is a twisted signature operator, then the local data in formula (3) may be expressed in terms of the tangent bundle and the vector bundles associated to the $`H`$-structure restricted to $`M^{S^1}`$. Next recall that the Pontrjagin numbers of $`M`$ are determined by twisted signatures, where the twist bundles are associated to the tangent bundle via some orthogonal representation. Since the $`S^1`$-action lifts canonically to these bundles one can consider the $`S^1`$-equivariant twisted signatures. If the $`S^1`$-action on $`M`$ has no fixed point these indices have to vanish since the index function $`t\text{-}ind`$ vanishes identically. In this case all Pontrjagin numbers of $`M`$ vanish (for an elementary proof see \[Bo\]) which implies that $`M`$ represents an element of order two in the oriented bordism ring. The following reformulation is convenient for our purposes. Given two $`S^1`$-manifolds $`M_1`$ and $`M_2`$ we say that they have the same local $`𝐒^\mathrm{𝟏}`$-geometry if there exists an $`S^1`$-equivariant orientation preserving diffeomorphism from the normal bundle of the fixed point manifold $`M_1^{S^1}`$ to the normal bundle of $`M_2^{S^1}`$. In this case one can glue $`M_1`$ and $`M_2`$ along the fixed point manifold together to obtain a manifold with fixed point free $`S^1`$-action which is bordant to $`M_1M_2`$ (cf. for example \[CoFl\], Th. (22.1)). We summarize the discussion in ###### Corollary 3.2. Let $`M_1`$ and $`M_2`$ be $`S^1`$-manifolds with the same local $`S^1`$-geometry. Then $`M_1M_2`$ represents a torsion element in the oriented bordism ring. $`\mathrm{}`$ We shall combine the corollary with Lemma 2.1 to obtain bordism-finiteness for the class of manifolds for which a negative multiple of the first Pontrjagin class is a sum of squares. Let us call such manifolds semi-negative, i.e. a manifold $`M`$ is semi-negative if there exists a finite number of classes $`y_jH^2(M;)`$ and a positive integer $`N`$ such that $`Np_1(M)+y_j^2=0`$. Note that $`M`$ is semi-negative if and only if $`M`$ admits complex line bundles $`L_j`$ for which the first Pontrjagin class of $`E:=L_j+NTM`$ vanishes. Of course many manifolds such as $`3`$-connected manifolds with non-vanishing first rational Pontrjagin class (e.g. quaternionic projective space $`P^k`$, $`k>1`$) cannot be semi-negative. However the class of semi-negative manifolds is quite rich and contains some interesting families. ###### Remarks 3.3. The class of semi-negative manifolds contains 1. manifolds with $`p_1`$ torsion, e.g. $`BO8`$-manifolds or $`4`$-connected manifolds. 2. manifolds with cohomology ring in degree $`4`$ like a $`4`$-manifold with indefinite intersection form. 3. in a given dimension all but a finite number of complete intersections. 4. infinitely many homotopy complex projective spaces in any given dimension $`2n6`$. ###### Theorem 3.4. Let $`C`$ and $`m`$ be natural numbers. The class of semi-negative connected $`m`$-dimensional oriented manifolds with Euler characteristic $`C`$ and an $`S^3`$-action satisfying $`()`$ contains only finitely many oriented bordism types. Proof: The theorem follows from Corollary 3.2 once we know that the local $`S^1`$-geometry is determined by $`(m,C)`$ up to finite ambiguity. Since the oriented bordism group $`\mathrm{\Omega }_m^{SO}`$ is finite in dimension $`m0mod4`$ we may assume $`m=4l`$. Let $`M`$ be a $`4l`$-dimensional manifold with $`S^3`$-action satisfying $`()`$ and Euler characteristic $`\chi (M)=C^{}C`$. By the classical Lefschetz fixed point formula the induced $`S^1`$-action has $`C^{}`$ isolated fixed points. Next assume that $`M`$ is semi-negative, i.e. $`p_1(L_1+\mathrm{}+L_k+NTM)=0`$ for a positive integer $`N`$ and complex line bundles $`L_j`$. By Lemma 2.1 the tangential weights of $`M`$ are determined by the dimension of $`M`$ up to finite ambiguity. So all manifolds in the theorem with Euler characteristic equal to $`C^{}`$ represent only a finite set of local $`S^1`$-geometries. Since $`0C^{}C`$ the theorem follows from Corollary 3.2. $`\mathrm{}`$ The theorem immediately implies Theorem 1.1. It extends to the complex bordism ring if one assumes in addition that the induced $`S^1`$-action preserves the stable almost complex structure. A corresponding result for the $`Spin^c`$-bordism ring does not hold since an $`S^3`$-action always lifts to a given $`Spin^c`$-structure and there are just too many of them. However Theorem 3.4 admits the following refinement. Let $`X(k)`$ be the Cartesian product of $`BSO`$ and $`k`$ copies of $`P^{\mathrm{}}`$. For a fixed positive integer $`N`$ let $`f_N:X(k)K(,4)`$ be the map which classifies $`Np_1+_{j=1}^kx_j^2H^4(X(k);)`$. Here $`p_1H^4(BSO;)`$ denotes the universal first Pontrjagin class and $`x_j`$ denotes a generator for the $`j`$-th copy of $`H^2(P^{\mathrm{}};)`$. Let $`B(k,N)X(k)`$ be the pullback of the path fibration $`E(,4)K(,4)`$ via $`f_N`$ and let $`\pi :B(k,N)BSO`$ denote the projection. Fix a classifying map $`\nu :MBSO`$ for the stable normal bundle of $`M`$. A $`B(k,N)`$-structure for $`M`$ is the isotopy class $`h`$ of a lift $`MB(k,N)`$ of $`\nu `$ in the fibration $`\pi :B(k,N)BSO`$. For fixed $`k`$ and $`N`$ we call $`(M,h)`$ a $`𝐁`$-manifold. Note that $`(M,h)`$ determines $`k`$ complex line bundles $`L_j`$ over $`M`$ such that $`Np_1(M)+_{j=1}^kp_1(L_j)=0`$. In particular, $`M`$ is semi-negative. Any polynomial in $`p_i(M)`$ and $`c_1(L_j)`$ is called a characteristic class of the $`B`$-manifold $`(M,h)`$. The characteristic numbers are defined as the values of the characteristic classes on the fundamental cycle of $`M`$. By the Pontrjagin lemma the characteristic numbers only depend on the bordism class of $`(M,h)`$. The bordism group $`\mathrm{\Omega }_{}^B`$ of $`B`$-manifolds may be studied in terms of stable homotopy theory using the Pontrjagin-Thom construction. ###### Proposition 3.5. For fixed $`k`$ and $`N`$ the bordism group $`\mathrm{\Omega }_n^B`$, $`n`$, is finitely generated. The group $`\mathrm{\Omega }_n^B`$ is completely determined by the characteristic numbers. Proof: Since the argument is standard (cf. \[La\], \[St\]) we only sketch it. First apply the construction above to $`BSO(r)`$ for $`r`$. So $`X_r`$ is the product of $`BSO(r)`$ and $`k`$ copies of $`P^{\mathrm{}}`$, $`B_rX_r`$ is obtained as pullback of the path fibration $`E(,4)K(,4)`$ and $`\pi _r:B_rBSO(r)`$ is the projection. Note that $`B(k,N)=lim_r\mathrm{}B_r`$. Let $`\gamma _rB_r`$ be the pullback of the universal vector bundle over $`BSO(r)`$ via $`\pi _r`$ and let $`M(\gamma _r)`$ be its Thom space. The Pontrjagin-Thom construction gives an isomorphism $`\mathrm{\Phi }:\mathrm{\Omega }_n^B\pi _{n+r}(M(\gamma _r))`$ for $`rn`$, which defines an isomorphism $$\mathrm{\Phi }:\mathrm{\Omega }_n^B\underset{r\mathrm{}}{lim}\pi _{n+r}(M(\gamma _r))\pi _n(M(\gamma ))$$ between the bordism group of $`n`$-dimensional $`B`$-manifolds and the $`n`$-th homotopy of the associated spectrum $`M(\gamma )`$. Since $`M(\gamma _r)`$ is a CW-complex with finite skeletons $`\mathrm{\Omega }_n^B`$ is finitely generated. Next assume the $`n`$-dimensional $`B`$-manifold $`(M,h)`$ has vanishing characteristic numbers. For the second statement it suffices to show that the bordism class of $`(M,h)`$ vanishes in $`\mathrm{\Omega }_n^B`$. Note that the space of characteristic classes of $`(M,h)`$ is equal to $`\widehat{h}^{}(H^{}(B_r;))`$, $`rn`$, where $`\widehat{h}:MB_r`$ represents $`h`$. Since $`(M,h)`$ has vanishing characteristic numbers $`\widehat{h}_{}(\mu _M)`$ vanishes in $`H_n(B_r;)`$. Now apply the Thom isomorphism for the normal bundle of $`M`$ and the bundle $`\gamma _rB_r`$ to conclude that the composition of the Pontrjagin-Thom map $`\mathrm{\Phi }`$ and the rational Hurewicz homomorphism $`\mathrm{\Psi }`$ $$\mathrm{\Omega }_n^B\stackrel{\Phi }{}\pi _{n+r}(M(\gamma _r))\stackrel{\Psi }{}H_{n+r}(M(\gamma _r);),$$ $`rn`$, maps $`(M,h)`$ to zero. Since $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are isomorphisms $`(M,h)`$ vanishes in $`\mathrm{\Omega }_n^B`$. $`\mathrm{}`$ ###### Theorem 3.6. Let $`C`$ and $`m`$ be natural numbers. For fixed $`k`$ and $`N`$ the class of connected $`m`$-dimensional $`B`$-manifolds with Euler characteristic $`C`$ and an $`S^3`$-action satisfying $`()`$ contains only finitely many $`B`$-bordism types. Proof: By Proposition 3.5 we may assume that $`m=2n`$. Let $`(M,h)`$ be a $`2n`$-dimensional $`B`$-manifold with Euler characteristic $`C`$ and an $`S^3`$-action satisfying $`()`$. The map $`h`$ induces via the projection of $`B`$ to the $`k`$-fold product of $`P^{\mathrm{}}`$ a classifying map for $`k`$ complex line bundles $`L_j`$ which satisfy $`Np_1(M)+_{j=1}^kp_1(L_j)=0`$. We want to show that the characteristic numbers of $`(M,h)`$ are determined by $`(m,C,k,N)`$ up to finite ambiguity. To this end let $`J`$ denote the subring of $`K(M)`$ generated by the complex line bundles $`L_j`$ and vector bundles associated to the tangent bundle. Let $`ch:K(M)H^{}(M;)`$ be the Chern character. We note that $`ch(J)`$ is the subspace $`V`$ of $`H^{}(M;)`$ which is spanned by the characteristic classes of the $`B`$-manifold $`(M,h)`$. Next we identify the characteristic numbers of $`(M,h)`$ with certain twisted signatures. By the cohomological version of the index theorem (cf. \[AtSiIII\]) the index of the signature operator of $`M`$ twisted with $`FJ`$ is given by $$sign(M;F)=\underset{i=1}{\overset{n}{}}(u_i\frac{1+e^{u_i}}{1e^{u_i}})ch(F),\mu _M,$$ where $`\pm u_i`$ are the formal roots of $`M`$, $`\mu _M`$ is the fundamental cycle of $`M`$ and $`,`$ denotes the pairing between cohomology and homology. This implies that $`\{sign(M;F)FJ\}`$ spans the $``$-vector space $`V,\mu _M`$. Hence the characteristic numbers of $`(M,h)`$ are determined by twisted signatures, where the twist bundle is an element of $`J`$. We are now in the position to prove the theorem from the Lefschetz fixed point formula. Choose a finite set $`\{E_1,\mathrm{},E_r\}J`$ ($`r`$ depends on $`k`$ and the dimension of $`M`$) of vector bundles which span $`J`$. Here each $`E_i`$ is given by a universal polynomial (which only depends on $`(m,k)`$) in the complex line bundles $`L_j`$ and vector bundles associated to the tangent bundle. We view $`E_i`$ as an $`S^3`$-equivariant vector bundle by lifting the $`S^3`$-action (uniquely) to each $`L_j`$. By Lemma 2.1 the tangential weights and the weights of $`L_j`$ are determined by $`(m,N)`$ up to finite ambiguity. This implies the same for the weights of $`E_i`$. Next consider the signature operator twisted by $`E_i`$. By the discussion of the Lefschetz fixed point formula after Theorem 3.1 its $`S^1`$-equivariant index is equal to a sum of $`C`$ local contributions which only depend on the weights of $`E_i`$ and the tangential weights. Since these belong to a finite set which only depends on $`(m,k,N)`$ we conclude that the ordinary twisted signatures $`sign(M;E_i)`$, $`i=1,\mathrm{},r`$, are determined by $`(m,C,k,N)`$ up to finite ambiguity. This implies the same for the characteristic numbers of the $`B`$-manifold $`(M,h)`$. Now the theorem follows from Proposition 3.5. $`\mathrm{}`$ ## 4 Examples of semi-simple group actions In this section we show that Theorem 3.4 is sharp in the sense that bordism-finiteness fails if one weakens the assumptions on the $`S^3`$-action or the first Pontrjagin class. As explained in the introduction one cannot expect bordism-finiteness for $`S^1`$-actions if one allows arbitrary fixed point sets. We restrict to S<sup>1</sup>-actions with isolated fixed points which satisfy a prescribed upper bound for the number of fixed points. Also we assume that the action extends to an action of $`S^3`$. In this situation, by Theorem 3.4, bordism-finiteness holds if the $`S^3`$-action has a fixed point and the manifolds are semi-negative. The following two propositions show that both assumptions are necessary. ###### Proposition 4.1. There exist connected $`20`$-dimensional semi-negative manifolds $`M_l`$, $`l`$, with Euler characteristic equal to $`12`$ which represent distinct oriented bordism classes such that each $`M_l`$ supports a fixed point free $`S^3`$-action with isolated $`S^1`$-fixed points. ###### Proposition 4.2. There exist connected $`20`$-dimensional manifolds $`N_l`$, $`l`$, with Euler characteristic equal to $`23`$ which represent distinct oriented bordism classes such that each $`N_l`$ supports an $`S^3`$-action with fixed point and isolated $`S^1`$-fixed points. Note that by Theorem 3.4 almost all of the $`N_l`$ are not semi-negative. We remark that examples as in the propositions are necessarily of dimension $`4k8`$ since in dimension $`4`$ the upper bound on the number of isolated fixed points gives a bound on the absolute value of the signature. The remaining part of this section is devoted to the proof of the propositions above. We use a kind of induction to extend actions on the base of a fibre bundle to the total space if the fibration is associated to a principal torus bundle. Let $`G`$ be a compact connected Lie group which acts from the left on a connected manifold $`Z`$. Let $`F`$ be a connected manifold with left $`U(1)`$-action. Assume the first Betti number $`b_1(Z)`$ vanishes or $`G`$ is simply-connected. For $`yH^2(Z;)`$ let $`SZ`$ denote the $`U(1)`$-principal bundle with first Chern class equal to $`y`$ and let $`M:=S\times _{U(1)}F`$ be the associated fibre bundle. In \[HaYo\] it was shown that the $`G`$-action lifts to $`S`$ (in fact uniquely if $`G`$ is simply-connected). For a fixed lift $`G`$ acts on $`M`$ by $`g(s,f)_{}:=(gs,f)_{}`$. Note that the projection $`MZ`$ is $`G`$-equivariant. We now restrict to the case $`G=S^3`$ and fix $`S^1S^3`$. Let $`a_Y`$ denote the weight of the $`S^1`$-action on $`S`$ restricted to a connected component $`Y`$ of $`Z^{S^1}`$, i.e. at a point of $`Y`$ (and hence at any point of $`Y`$) the $`S^1`$-action on the fibre of $`S`$ has character $`\lambda \lambda ^{a_Y}`$. For further reference we note the elementary ###### Lemma 4.3. Assume none of the weights $`a_Y`$ vanish. Then the $`S^1`$-action on $`M`$ has isolated fixed points if this holds for the $`S^1`$-action on $`Z`$ and the $`U(1)`$-action on $`F`$. $`\mathrm{}`$ We now begin with the construction of the examples mentioned before. The first series consists of Cayley plane bundles over $`Z:=S^2\times S^2`$ which support an $`S^3`$-action with isolated $`S^1`$-fixed points but no $`S^3`$-fixed point. We take the left $`S^3`$-action on $`Z`$ induced from the homogeneous action on each copy of $`S^2`$. This action has isolated $`S^1`$-fixed points but no $`S^3`$-fixed points. Its principal stabilizer is equal to $`_2`$, the center of $`S^3`$. As fibre $`F`$ we take the Cayley plane $`Cl_2=F_4/Spin(9)`$. We fix an orientation of $`Cl_2`$. Next we choose an embedding $`j:U(1)T`$, where $`T`$ is a maximal torus of $`Spin(9)`$, such that the induced $`U(1)`$-action on $`Cl_2`$ is effective and has isolated fixed points. Let $`\gamma _iZ`$ denote the pullback of the Hopf bundle over $`S^2`$ under the projection of $`Z`$ on the $`i`$-th factor and let $`z_i:=c_1(\gamma _i)H^2(Z;)`$. Equip $`Z`$ with the orientation dual to $`z_1z_2`$. Note that the weights of the $`S^3`$-equivariant line bundle $`\gamma _1^a\gamma _2^b`$ at the $`S^1`$-fixed points of $`Z`$ have the form $`\pm a\pm b`$. To construct $`M_l`$ fix positive integers $`abmod2`$, let $`S_l`$, $`l`$, be the $`U(1)`$-principal bundle associated to $`(\gamma _1^a\gamma _2^b)^{2l+1}`$ and set $`M_l:=S_l\times _{U(1)}Cl_2`$. Note that $`M_l`$ comes with an orientation induced from the orientations of $`Z`$ and $`Cl_2`$. As explained above the $`S^3`$-action on $`Z`$ lifts to $`M_l`$. For the induced $`S^1`$-action some of the tangential weights of $`M_l`$ are odd since the weights of $`S_l`$ have the form $`(2l+1)(\pm a\pm b)`$ and the action of $`U(1)`$ on $`Cl_2`$ is effective. In particular, the principal isotropy group of the $`S^3`$-manifold $`M_l`$ is trivial. Proof of Proposition 4.1: The Euler characteristic $`\chi (M_l)`$ of the oriented $`20`$-dimensional manifold $`M_l`$ is just the product of the Euler characteristic of the base and the fibre, thus $`\chi (M_l)=12`$. The $`S^3`$-action on $`M_l`$ has no fixed points since $`p:M_lZ`$ is equivariant. By Lemma 4.3 the induced $`S^1`$-action on $`M_l`$ has isolated fixed points. Since $`H^{}(M_l;)`$ is isomorphic to $`H^{}(Z;)`$ in degree $`4`$ (the fibre is $`7`$-connected) each $`M_l`$ is semi-negative (see Remarks 3.3). We compute the Milnor number $`s_{10}(TM_l),\mu _{M_l}`$ to show that the manifolds $`M_l`$ represent distinct bordism classes (for a real vector bundle with formal roots $`\pm u_i`$ the class $`s_{2t}`$ is defined as $`u_i^{2t}`$). Let $`\pi :EBU(1)`$ denote the pullback of $`BSpin(9)BF_4`$ to $`BU(1)`$ via $`Bj`$ and let $`E^{\mathrm{}}E`$ denote the tangent bundle along the fibres. One computes using \[BoHi\] that $`\pi _!(s_{10}(E^{\mathrm{}}))=\alpha x^2`$, where $`\alpha 0`$ (cf. for example \[De1\], Section 4.2). Here $`x`$ is a generator of $`H^2(BU(1);)`$ and $`\pi _!`$ is the push-forward in cohomology. Next note that the map $`f:ZBU(1)`$ which classifies $`(2l+1)(az_1+bz_2)`$ is covered by $`\stackrel{~}{f}:M_lE`$ and the tangent bundle along the fibres of $`p:M_lZ`$ is isomorphic to $`\stackrel{~}{f}^{}(E^{\mathrm{}})`$. Now compute $$s_{10}(TM_l),\mu _{M_l}=s_{10}(\stackrel{~}{f}^{}(E^{\mathrm{}})p^{}(TZ)),\mu _{M_l}=p_!(s_{10}(\stackrel{~}{f}^{}(E^{\mathrm{}}))),\mu _Z=$$ $$f^{}(\pi _!(s_{10}(E^{\mathrm{}}))),\mu _Z=\alpha (2l+1)^2(az_1+bz_2)^2,\mu _Z=2\alpha (2l+1)^2ab.$$ The computation shows that the $`M_l`$ represent distinct bordism classes. $`\mathrm{}`$ The next series is constructed from the series above. Equip $`P^{10}`$ with the $`S^3`$-action induced by the direct sum of the trivial one-dimensional complex representation and the irreducible complex $`S^3`$-representation of dimension $`10`$. Note that $`S^3`$ acts on $`P^{10}`$ with fixed point, the principal isotropy group is trivial and the induced $`S^1`$-action has isolated fixed points. The principal isotropy group of $`M_l`$ is also trivial (see the discussion before the proof of Prop. 4.1). Now define $`N_l`$ by taking the equivariant connected sum of $`M_l`$ and $`P^{10}`$ along a principal orbit. Proof of Proposition 4.2: First note that the disjoint union of $`M_l`$ and $`P^{10}`$ is bordant to $`N_l`$. Also $`\chi (N_l)=\chi (M_l)+\chi (P^{10})`$. Hence, by Proposition 4.1 the manifolds $`N_l`$, $`l`$, represent distinct bordism classes and have Euler characteristic $`\chi (N_l)=23`$. By construction $`S^3`$ acts on $`N_l`$ with fixed point and the induced $`S^1`$-action has isolated fixed points. $`\mathrm{}`$ ## 5 Homotopy complex projective spaces and complete intersections In this section we prove the theorem on homotopy complex projective spaces given in the introduction and a related result for complete intersections (see Theorem 5.1 below). Let $`M`$ be a $`2n`$-dimensional closed manifold with $`H^{}(M;)H^{}(P^n;)`$ and non-trivial $`S^1`$-action. Let $`\gamma `$ denote the complex line bundle with $`c_1(\gamma )=x`$, where $`x`$ is a fixed generator of $`H^2(M;)`$. By \[HaYo\] the $`S^1`$-action lifts to $`\gamma `$. We denote the weight of $`\gamma `$ and the tangential weights at a connected component $`Y_s`$ of $`M^{S^1}`$ by $`a_s`$ and $`m_{s,1},\mathrm{},m_{s,n}`$, respectively. In the proof we use the following information on the weights: For fixed $`s`$ the weights of $`\gamma `$ are related to the tangential weights by $$\|\underset{ts}{}(a_ta_s)^{n_t+1}\|=\|\underset{m_{s,i}0}{}m_{s,i}\|,$$ (4) where $`n_t`$ denotes the complex dimension of $`Y_t`$. To prove this identity one either applies the localization theorem in $`K`$-theory to the $`S^1`$-action and induced $`_p`$-actions for $`p`$ a large prime number (cf. \[Pe\], Th. 2.8) or uses cohomological means (cf. \[Br\], Ch. VII, Th. 5.5). Proof of Theorem 1.2: For $`n=1`$ the theorem is trivial. For $`n=2`$ it follows for example from the classification of $`4`$-dimensional $`S^3`$-manifolds (cf. \[MePa\]). So assume $`n3`$. Let $`M`$ be an oriented homotopy $`P^n`$ with $`S^3`$-action satisfying $`()`$. Let $`Y_0,\mathrm{},Y_n`$ denote the isolated fixed points under the induced $`S^1`$-action on $`M`$. We assume that $`Y_0`$ is also fixed by $`S^3`$. Hence, $`a_0`$ vanishes and the tangential weights $`m_{0,1},\mathrm{},m_{0,n}`$ are determined by the dimension $`m=2n`$ up to finite ambiguity. Therefore $`\|_{i=1}^nm_{0,i}\|`$ is bounded from above by a constant which only depends on $`m`$. Now apply formula (4) twice to conclude that the same holds for the absolute value of all the weights $`a_s`$ and $`m_{s,i}`$. Next we argue as in the proof of Theorem 3.6 to show that the values of polynomials in the Pontrjagin classes and the generator $`xH^2(M;)`$ on the fundamental cycle belong to a finite set which only depends on $`m`$. Since the cohomology ring of $`M`$ is a truncated polynomial ring in $`x`$ the Pontrjagin classes are also determined up to finite ambiguity. By simply-connected surgery theory it follows that the diffeomorphism type of $`M`$ belongs to a finite set which only depends on the dimension $`m`$. $`\mathrm{}`$ Next we consider complete intersections. A complete intersection $`V_n^{(d_1,\mathrm{},d_r)}`$ of complex dimension $`n`$ and multidegree $`(d_1,\mathrm{},d_r)`$, $`d_i2`$, is defined as the transversal intersection of hypersurfaces of degree $`d_i`$, $`i=1,\mathrm{},r`$, in $`P^{n+r}`$. Thom showed that the diffeomorphism type of $`V_n^{(d_1,\mathrm{},d_r)}`$ is completely determined by $`n`$ and $`(d_1,\mathrm{},d_r)`$. ###### Theorem 5.1. In a fixed complex dimension there are only finitely many complete intersections with an $`S^3`$-action satisfying $`()`$. We remark that the $`\widehat{A}`$-vanishing theorem of Atiyah-Hirzebruch (cf. \[AtHi\]) may be used to show that there are only finitely many $`Spin`$-complete intersections with non-trivial $`S^1`$-action in a fixed even complex dimension. Also results of Hattori (cf. \[Ha\]) on the equivariant $`Spin^c`$-Dirac operator imply finiteness of the number of complete intersections with non-trivial $`S^1`$-action if one assumes that the action preserves the induced stable almost complex structure. Proof of Theorem 5.1: Assume $`M:=V_n^{(d_1,\mathrm{},d_r)}`$ admits an $`S^3`$-action satisfying $`()`$. Let $`\gamma `$ denote the pullback of the dual Hopf bundle over $`P^{n+r}`$ via $`i:MP^{n+r}`$ and let $`x:=c_1(\gamma )H^2(M;)`$. Recall from \[Hi\] that $`p(M)=(1+x^2)^{n+r+1}_{i=1}^r(1+d_i^2x^2)^1`$ and $`x^n,\mu _M=_{i=1}^rd_i`$. Hence, for all but a finite number of multidegrees $`M`$ is semi-negative with $`N=1`$ and $`L_j=\gamma `$. For semi-negative $`M`$ the tangential weights and the weights of $`\gamma `$ are determined by the complex dimension $`n`$ up to finite ambiguity by Lemma 2.1. We argue as before (see the proof of Th. 3.6) to conclude that the values of polynomials in the Pontrjagin classes and $`x`$ on the fundamental cycle $`\mu _M`$ belong to a finite set which only depends on $`n`$. In particular this holds for $`x^n,\mu _M=d_i`$. Since $`d_i2`$ the theorem follows. $`\mathrm{}`$ Anand Dessai Department of Mathematics University of Augsburg D-86159 Augsburg e-mail: dessai@math.uni-augsburg.de http://www.math.uni-augsburg.de/geo/dessai.html
warning/0003/hep-th0003115.html	ar5iv	text	# Glueball Spectrum for QCD from 𝐴⁢𝑑⁢𝑆 Supergravity Duality This work was supported in part by the Department of Energy under Contracts No. DE-FG02-91ER40676 and No. DE-FG02-91ER40688 ## 1 Introduction The Maldacena duality conjecture and its further extensions state that there is an exact equivalence between large N conformal field theories in d-dimensions and string theory in $`\mathrm{𝐀𝐝𝐒}^{𝐝+\mathrm{𝟏}}\times 𝐌`$. Subsequently Witten suggested how to break explicitly the conformal (and SUSY) symmetries to arrive at a dual gravity description for $`SU(N)`$ quarkless $`QCD_4`$. Thus we may have at last a definite proposal for the long sought “QCD string”. As anticipated by ’t Hooft the dual correspondence for the large $`1/N`$ expansion for $`SU(N)`$ Yang-Mills theory is the perturbative expansion of string theory. Still this theory is difficult to formulate, let alone solve. At present explicit calculations also require taking the strong coupling limit, $`g^2N\mathrm{}`$, where the string tension goes to infinity ($`\alpha ^{}0`$) and the dual theory is classical gravity. In this paper we complete the study of the glueball spectrum for the strong coupling dual description of $`QCD_4`$. For comparison the analogous spectrum calculation is presented for $`QCD_3`$, which shows a very similar pattern, which in qualitative terms can be traced to the underlying flat space T duality between type IIA and IIB string theories which in turn are the AdS duals to $`QCD_4`$ and $`QCD_3`$ respectively. The goal is to learn more about the AdS/Yang-Mills correspondence by comparing the AdS strong coupling spectrum with the rather well determined glueball spectrum in lattice gauge theory. Of course, the strong coupling expansion at best can provide a rough guide to the underlying physics. Nonetheless the correspondence to the continuum (i.e. weak coupling) limit of the lattice spectrum is surprisingly good. This comparison may prove useful to provide support for the conjectured Maldacena duality and to give specific information on the major strong coupling artifacts that must be removed as one approaches universality at the ultraviolet fixed point. In addition to new spectral calculations, we summarize the earlier work by many authors . In particular, we extend our earlier paper on the tensor glueball for $`QCD_3`$ on an $`AdS^5`$ black hole background to the physically relevant case of $`QCD_4`$. For $`QCD_3`$, we found that the tensor spectrum ($`2^{++}`$) was degenerate with dilaton($`0^{++}`$) and the axion ($`0^+`$). However, the mass gap set by a lower scalar glueball obeys the inequality, $$m(0^{++})<m(2^{++}).$$ (1) This scenario is repeated for $`QCD_4`$. The dilaton mode ($`0^{++}`$) coupling to $`Tr[F^2]`$ remains degenerate with the tensor ($`2^{++}`$), but the axion ($`0^+`$) is heavier, consistent with lattice results. Due to the trace anomaly, the lowest mass scalar ($`0^{++}`$) couples to the energy density $`T_{00}`$ and again it obeys this inequality above, Eq. (1). As in our earlier work, the goal is to see the details of the spin structure of the lowest glueball states, which we believe is most sensitive to the underlying gauge theory. We find that our analysis combined with all the earlier results can give a systematic and complete strong coupling glueball spectrum. It is useful to end this introduction with a rough overview of our results and the organization of the paper. The geometrical construction for $`QCD_4`$ is roughly as follows. One starts with 11 dimensional M theory on $`\mathrm{𝐀𝐝𝐒}^\mathrm{𝟕}\times 𝐒^\mathrm{𝟒}`$. The seven dimensional $`AdS^7`$ we take to have a radial co-ordinate $`r`$ and Euclidean space-time co-ordinates $`x_1,x_2,x_3,x_4,x_5`$ and $`x_{11}`$. The “eleventh” dimension is taken to be compact, reducing the theory to type IIA string theory. Matter at the center of this space ($`r=0`$) consists of N D4-branes (or NS 5-branes wrapping $`S^1`$ in the 11th coordinate) with world volume co-ordinates $`x_1\mathrm{}x_5`$. The 5-d Yang-Mills CFT “living on” the brane is dimensionally reduced to $`QCD_4`$ by raising the “temperature”, $`\beta ^1`$, in a direction $`x_5=\tau `$, parallel to the brane. The new metric is an $`AdS^7`$ black hole with $`x_{11}`$ compact . Compact directions on $`S^4`$ will be denoted by $`x_\alpha `$, $`\alpha =7,8,9,10.`$ The strong coupling glueball calculation consists of finding the normal modes for the bosonic components of the supergraviton multiplet in the $`\mathrm{𝐀𝐝𝐒}^\mathrm{𝟕}\times 𝐒^\mathrm{𝟒}`$ black hole background. We are only interested in excitations that lie in the superselection sector for $`QCD_4`$. So we can ignore modes for all non-trivial harmonics in $`S^5`$ that carry a non-zero R charge and all Kaluza-Klein (KK) modes in the two $`S^1`$ circles with U(1) KK charges. Imposing these restriction and exploiting symmetries of the background metric reduce the problem to six independent wave equations, referred to as $`S_4,T_4,V_4,N_4,M_4`$ and $`L_4`$ in the text. In Fig. 1 (left side), we plot the low mass states for each equation, labeling the quantum numbers for each level. We can identify the modes with the bosonic components of the zero mass sector of type IIA string theory: the graviton ($`G`$), the dilaton ($`\varphi `$), the NS-NS 2-form ($`B`$) and RR 1- & 3-forms ($`C_{(1)}`$, $`C_{(3)}`$). The spin degeneracy of the spectrum is due to a spurious $`O(4)`$ symmetry of the strong coupling approximation that combines the 11th and three spatial co-ordinates. However, as we explain in the text, all extra states not observed in the lattice data for the glueball spectrum, actually carry a discrete “$`\tau `$-parity” that places them, like KK modes, outside the QCD superselection sector. We have included in Fig. 1 (right side) the full spectrum for $`QCD_3`$. The entire analysis is very similar. Starting with type IIB strings in $`\mathrm{𝐀𝐝𝐒}^\mathrm{𝟓}\times 𝐒^\mathrm{𝟓}`$ which is dual to $`𝒩=4`$ SUSY Yang-Mills, one introduces a compact thermal circle forming an $`AdS^5`$ black hole. Again there are six independent wave equations, labeled by $`S_3,T_3,V_3,N_3,M_3`$ and $`L_3`$. They correspond to fluctuations for type IIB fields: the graviton ($`G`$), the dilaton ($`\varphi `$), the NS-NS 2-form ($`B`$) as before and RR 0- & 2-forms ($`C_{(0)}`$, $`C_{(2)}`$). For both cases we also include volume fluctuations in the compact sphere $`S^4`$ and $`S^5`$ for type IIA and IIB respectively. In Sec. 2, we give general arguments for the spin and degeneracy of glueball states for $`QCD_4`$ followed in Sec. 3 by the analysis for $`QCD_3`$. For each, we give the resultant six wave equations and numerical values for the first ten levels. (Derivations for these equations are explained further in Appendix A.) For all but the lowest eigenvalue, the glueball masses, $`m_n`$, are well approximated by the WKB expansions : $`m_n^2=\mu ^2(n^2+\delta n+\gamma )`$. For the first level ($`n=0`$), we provide a simple but reasonably tight variational upper bound . (See Appendix B for details.) In Sec. 4, we give the parity and charge conjugation quantum numbers of the glueball states using the Born-Infeld action to determine the quantum numbers of the couplings between gravity fields and gauge fields. The striking similarity between the $`QCD_4`$ and the $`QCD_3`$ spectra (see Fig. 1) can be understood qualitatively in terms of T-duality, which relates $`D4`$ branes in IIA to $`D3`$ branes in IIB. In Sec. 5, we compare the AdS strong coupling spectrum with the well determined levels from lattice $`QCD_4`$ and remark on the relationship to the constituent gluon picture. No extra states are present in the $`AdS`$ spectrum that couple to QCD operators, although the absence of the low mass $`2^+`$ state is noted. We also show how the strong coupling expansion for the Pomeron intercept may be used to provide an estimate of the coupling at the crossover between the strong and weak coupling regimes. ## 2 Glueball Spectrum for $`QCD_4`$ To approach $`QCD_4`$ one begins with M theory on $`\mathrm{𝐀𝐝𝐒}^\mathrm{𝟕}\times 𝐒^\mathrm{𝟒}`$. We compactify the “eleventh” dimension (on a circle of radius $`R_1`$) to reduce the theory to type IIA string theory and then following the suggestion of Witten raise the “temperature”, $`\beta ^1`$, with a second compact radius $`R_2`$ in a direction $`\tau `$, with $`\beta =2\pi R_2`$. On the second “thermal” circle, the fermionic modes have anti-periodic boundary conditions breaking conformal and all SUSY symmetries. This lifts the fermionic masses and also the scalar masses, through quantum corrections. The ’t Hooft coupling is $`g^2N=2\pi g_sNl_s/R_2`$, in terms of the closed string coupling, $`g_s`$ and the string length, $`l_s`$ . Therefore, in the scaling limit, $`g^2N0`$, if all goes as conjectured, there should be a fixed point mapping type IIA string theory onto $`SU(N)`$ pure Yang-Mills theory. We consider the strong coupling limit at large N, where the string theory becomes classical gravity in the $`AdS^7`$ black hole metric, $$ds^2=(r^2\frac{1}{r^4})d\tau ^2+r^2\underset{i=1,2,3,4,11}{}dx_i^2+(r^2\frac{1}{r^4})^1dr^2+\frac{1}{4}d\mathrm{\Omega }_4^2,$$ (2) with radius of curvature, $`R_{AdS}^3=8\pi g_sNl_s^3`$. We have removed all dimensionful parameters in the metric by a normalization setting $`R_{AdS}=1`$ and $`\beta =2\pi /3`$. ### 2.1 Spin and Degeneracy of Glueball States In M theory the supergraviton is a single multiplet in 11-d with two bosonic fields - a graviton, $`G_{MN}`$, and a 3-form field, $`A_{MNL}`$, as designated in Table 1. After restricting all indices and co-ordinate dependence to $`AdS^7`$, we have a graviton, $`G_{\mu \nu }`$, a dilaton $`\varphi `$, and an NS-NS tensor field $`B_{\mu \nu }`$. In addition there are two RR fields, a one-form $`C_\mu `$ and a three-form $`C_{\mu \nu \lambda }`$. Furthermore, we will also consider the scalar modes coming from “volume” fluctuations for $`S^4`$. The relationship between M theory and IIA string theory nomenclature, after restricting to the $`AdS^7`$ subspace, is presented in Table 1. The table gives the $`J^{PC}`$ quantum numbers for all glueball states. The pattern of degeneracy (explained below) is indicated by the rows ending with the lowest eigenvalue for each of the six wave equations, Eq. (14): $`T_4,V_4,`$ etc. The task is to find all the quadratic fluctuations in the $`AdS^7`$ black hole background that might survive for $`QCD_4`$ in the scaling (weak coupling) limit, ignoring any Kaluza-Klein mode in compact manifolds (compactified $`S^1`$ for $`x_{11}`$, for $`\tau `$ and the spheres $`S^4`$). They are charge states in their own superselection sector that are clearly absent in the putative target theory. Additional “spurious” states will be discussed in Sec. 4 where we treat discrete symmetries. To count the number of independent fluctuations for a field of given spin, we adopt the following method. We imagine harmonic plane waves propagating in the AdS radial direction, $`r`$, with Euclidean time, $`x_4`$. For example metric fluctuation, $$G_{\mu \nu }=\overline{g}_{\mu \nu }+h_{\mu \nu }(x),$$ (3) in the background, $`\overline{g}_{\mu \nu }`$, are taken to have the form, $`h_{\mu \nu }(r,x_4)`$. There is no dependence on the spatial co-ordinates, $`x_i=(x_1,x_2,x_3,x_{11})`$ and the compactified “temperature” direction $`\tau `$. #### 2.1.1 Metric fluctuations The four dimensional field theory lives on the hypersurface co-ordinates, $`x_1,x_2,x_3,x_4`$, (with $`x_5=\tau `$ and $`x_{11}`$ as the two compactified coordinates.) A graviton has two polarization indices. If we were in flat space time, we could go to a gauge where these indices took values only among $`(x_1,x_2,x_3,x_{11},\tau )`$ and not from the set $`(r,x_4)`$. The polarization tensor should also be traceless. This leaves $`(5\times 6)/21=14`$ independent components. In the AdS space time, we can count the number of graviton modes the same way, though the actual modes that we construct will have this form of polarization only at $`r\mathrm{}`$; for finite $`r`$, other components of the polarization will be constrained to acquire nonzero values . We identify the spin content of these 14 components in two steps. First note that the background metric is $`SO(4)`$ symmetric. (It is flat in the first four of these directions, $`\overline{g}_{11}=\overline{g}_{22}=\overline{g}_{33}=\overline{g}_{11,11}=r^2`$, while it is “warped” in the $`\tau `$ direction, $`\overline{g}_{\tau \tau }=r^21/r^4`$). The system therefore has $`SO(4)`$ symmetry leading to three distinct equations corresponding to 9, 4 and 1 dimensional irreducible representations under $`SO(4)`$. In Table 1, these are denoted by $`T_4`$, $`V_4`$, and $`S_4`$ respectively. These representations lead us to a degenerate spectrum of spins under the physical $`SO(3)`$ rotations in $`x_1,x_2,x_3`$, which we list below: * $`9`$-dimensional representation breaks into $`5+3+1`$ under $`S0(3)`$, $`G_{ij}:h_{ij}{\scriptscriptstyle \frac{1}{3}}\delta _{ij}h_{kk}0`$ $``$ $`\text{spin-2},`$ (4) $`C_i:h_{i,11}=h_{11,i}0`$ $``$ $`\text{spin-1},`$ (5) $`\varphi :h_{11,11}=3h_{11}=3h_{22}=3h_{33}0`$ $``$ $`\text{spin-0},`$ (6) * $`4`$-dimensional representation breaks into $`3+1`$ under $`S0(3)`$, $`G_{i\tau }:h_{\tau i}=h_{i\tau }0`$ $``$ $`\text{spin-1},`$ (7) $`C_\tau :h_{\tau ,11}=h_{0\tau }0`$ $``$ $`\text{spin-0},`$ (8) * singlet under $`S0(3)`$, $`G\tau \tau :h_{\tau \tau }=4h_{11}=4h_{22}=4h_{33}=4h_{11,11}0`$ $``$ $`\text{spin-0},`$ (9) where $`i,j,k=1,2,3`$. In addition there is a scalar field, $`G_\alpha ^\alpha `$, coming from the metric on the $`S^4`$ sphere , (with $`m_{AdS}^2=72`$), which is referred to as $`L_4`$ in Table 1. #### 2.1.2 Three-form fields The behavior of the three-form field is discussed briefly in the Appendix, but we recall some of the main features here. The wave equation has a topological mass term which results in the equation being factorized into two first-order equations, yielding upon iteration two second-order equations. One solution is a massless 3-form field, which has solutions that are pure gauge for the case when there is no dependence on the sphere $`S^4`$, and is thus to be ignored. The other field gives a second-order equation with $`m_{AdS}^2=36`$, but the fact that we have a first-order equation as the primary equation reduces the degrees of freedom effectively to those of a massless 3-form field. Let the propagation directions be again $`(x_4,r)`$. If we consider the component with indices $`A_{123}`$ then we get a specific nonzero value also for the components $`A_{r\tau ,11}`$ and $`A_{4\tau ,11}`$. We can count the independent degrees of freedom by looking only at components that do not have the propagation directions $`x_4,r`$ among the indices. Thus we get the fields listed in Table 1. Reducing the $`SO(4)`$ states under rotations in $`x_1,x_2,x_3`$ yields: * $`4`$-dimensional representation breaks into $`3+1`$ under $`S0(3)`$, $`B_{ij}:A_{ij,11}0`$ $`\text{spin-1},`$ (10) $`C_{123}:A_{ijk}0`$ $`\text{spin-0},`$ (11) * $`6`$-dimensional representation into $`3+3`$ under $`S0(3)`$, $`B_{i\tau }:A_{i\tau ,11}0`$ $`\text{spin-1},`$ (12) $`C_{ij\tau }:A_{ij\tau }0`$ $`\text{spin-1}.`$ (13) The field equations for these states have amplitudes $`N_4`$ and $`M_4`$ as listed in Table 1. ### 2.2 Wave Equations and $`QCD_4`$ Glueball Spectra The wave equations for the metric fluctuations for $`QCD_4`$ have been obtained in Ref. by analyzing the linearized Einstein equations about the $`\mathrm{𝐀𝐝𝐒}^\mathrm{𝟕}\times 𝐒^\mathrm{𝟒}`$ black hole background which leads to three independent equations, $`T_4`$, $`V_4`$ and $`S_4`$ . Here we complete the spectral analysis giving the numerical values for all glueball masses. Fluctuations $`N_4`$, $`M_4`$ and $`L_4`$ can be found similarly, leading again to all together six independent equations for $`QCD_4`$, expressed in a manifestly hermitian form: $``$ $`{\displaystyle \frac{d}{dr}}(r^7r){\displaystyle \frac{d}{dr}}T_4(r)(m^2r^3)T_4(r)=0,`$ (14) $``$ $`{\displaystyle \frac{d}{dr}}(r^7r){\displaystyle \frac{d}{dr}}V_4(r)(m^2r^3{\displaystyle \frac{9}{r(r^61)}})V_4(r)=0,`$ $``$ $`{\displaystyle \frac{d}{dr}}(r^7r){\displaystyle \frac{d}{dr}}S_4(r)(m^2r^3+{\displaystyle \frac{432r^5}{(5r^62)^2}})S_4(r)=0,`$ $``$ $`{\displaystyle \frac{d}{dr}}(r^7r){\displaystyle \frac{d}{dr}}N_4(r)(m^2r^327r^5+{\displaystyle \frac{9}{r}})N_4(r)=0,`$ $``$ $`{\displaystyle \frac{d}{dr}}(r^7r){\displaystyle \frac{d}{dr}}M_4(r)(m^2r^327r^5{\displaystyle \frac{9r^5}{r^61}})M_4(r)=0,`$ $``$ $`{\displaystyle \frac{d}{dr}}(r^7r){\displaystyle \frac{d}{dr}}L_4(r)(m^2r^372r^5)L_4(r)=0.`$ We shall provide a more detailed discussion on how these equations can be obtained in Appendix A, while concentrating here on establishing our normalization convention. Consider first metric perturbations of the form $$h_{\mu \nu }=ϵ_{\mu \nu }(r)e^{ik_4x_4},$$ (15) with all other fields set to zero. We shall further fix gauge to $`h_{4\mu }=0`$, and from the linearized Einstein’s equation, we determine the discrete spectrum with $`k_4=im`$. Because of the $`SO(4)`$ symmetry in $`x_1,x_2,x_3,x_{11}`$, the system is highly degenerate. Three distinct equations for various perturbations can be obtained by the following procedure. Tensor: There are five independent perturbations which form the spin-2 representations of $`SO(3)`$: $$h_{ij}=q_{ij}r^2T_4(r)e^{mx_4},$$ (16) where $`i,j=1,2,3`$ and $`q_{ij}`$ is an arbitrary constant traceless-symmetric $`3\times 3`$ matrix. Vector: Consider perturbations: $$h_{i\tau }=h_{\tau i}=q_i\sqrt{r^61}V_4(r)e^{mx_4},$$ (17) where $`i=1,2,3`$ and $`q_i`$ is an arbitrary constant 3-vector. Both equations for $`T_4`$ and $`V_4`$ have also been obtained in Ref. by considering the corresponding degenerate scalar modes. Scalar: The analogous scalar perturbation is $$h_{\tau \tau }=(r^2r^4)S_4(r)e^{mx_4}.$$ (18) Three-form and Volume Scalar: Next we turn to 3-form fields. It is sufficient to consider $`A_{ij,11}=B_{ij}`$ $`=`$ $`ϵ_{ijk}q_kr^2N_4(r)e^{mx_4},`$ (19) $`A_{i\tau ,11}=B_{i\tau }`$ $`=`$ $`q_i\sqrt{r^61}M_4(r)e^{mx_4},`$ (20) where $`q_i`$ is again an arbitrary constant 3-vector. Note that the metric in the direction $`x_{11}`$ is the same as that in the directions $`i,j,k`$, so the above functions $`N_4`$, $`M_4`$ will also give the solutions for $`C_{ijk}`$ and $`C_{ij\tau }`$ respectively. (Fluctuations for $`B_{ij}`$ have been considered previously in Ref. . ) Lastly, for the volume perturbation, we consider $$h_\alpha ^\alpha =L_4(r)e^{mx_4}.$$ To calculate the discrete spectrum for each of these equations, one must apply the correct boundary conditions at $`r=1`$ and $`r=\mathrm{}`$. This issue has been discussed in several earlier papers . The boundary conditions are found by solving the indicial equation. In all cases the appropriate boundary condition at $`r=1`$ is the one without the logarithmic singularity. At $`r=\mathrm{}`$ the least singular boundary is required to have a normalizable eigenstate. (See Appendix B for a listing of all boundary conditions.) Matching boundary conditions from $`r=1`$ and $`r=\mathrm{}`$ results in a discrete set of eigenvalues $`m_n^2`$, where $`n`$ is the number of zeros in the wave function inside the interval $`r(1,\mathrm{})`$. We solved the eigenvalue equations by the shooting method, integrating from $`r_11`$ to large $`r_{\mathrm{}}\mathrm{}`$. The resultant spectrum is given in Table 2. To further check our results, we have compared our numerical masses to the WKB approximations, $$m_n^2\mu _4^2(n^2+\delta n+\gamma )+0(\frac{1}{n}),$$ (21) where $`\mu _4^2=36\pi (\mathrm{\Gamma }(2/3)/\mathrm{\Gamma }(1/6))^2`$. For each equation, individual integer constant, $`\delta `$, was determined analytical and constant $`\gamma `$ was fit to the numerical data in Table 4. (See Appendix B). The fits to the WKB formula are accurate to better than 0.1 % for all but the lowest ($`n=0`$) mode. For each lowest mode, we have also carried out an independent simple variational estimate. We note that in case, our numerically calculated value for $`m_0^2`$ is always close and respects the variational estimate as an upper bound. (See Table 5.) ## 3 Glueball Spectrum for $`QCD_3`$ For $`QCD_3`$ the construction of the supergravity dual begins with Maldacena’s conjecture for type IIB string theory in a $`\mathrm{𝐀𝐝𝐒}^\mathrm{𝟓}\times 𝐒^\mathrm{𝟓}`$ background metric. Here the dual theory is conjectured to be the conformal field theory for $`𝒩=4`$ SUSY SU(N) Yang-Mills in 4-d. The $`AdS`$ curvature is induced by N units of charge on N coincident D3 branes giving rise to a constant volume 5-form $`F_{(5)}`$ in the product manifold. The co-ordinates in the $`AdS^5`$ space, we label by one “radial” co-ordinate $`r`$ and four space-time co-ordinates, $`x_\mu `$, $`\mu =1,2,3,4`$, parallel to the D3 branes. The remaining five co-ordinates in $`S^5`$ are labeled by $`x_\alpha `$, $`\alpha =6,7,8,9,10`$. Following the suggestion of Witten for obtaining a supergravity dual to $`QCD_3`$, we break conformal and SUSY symmetries, by introducing a compact “thermal” co-ordinate $`x_4=\tau `$ with anti-periodic boundary on $`S^1`$ for the fermionic modes. The resultant metric at high temperature is an $`AdS^5`$ black hole, $$ds^2=(r^2\frac{1}{r^2})d\tau ^2+r^2\underset{i=1,2,3}{}dx_i^2+(r^2\frac{1}{r^2})^1dr^2+d\mathrm{\Omega }_5^2,$$ (22) with radius of curvature, $`R_{AdS}^4=4\pi g_sNl_s^4`$ and 3-d Yang-Mills coupling, $`g_3^2N=2g_sN/R`$ in terms of the string coupling $`g_s`$, string length $`l_s`$ and compact $`S^1`$ circumference $`\beta =2\pi R`$. We have removed all dimensionful parameters from the metric by adopting a simple normalization with $`R_{AdS}=1`$ and $`\beta =\pi `$, for the circumference of the thermal circle. At high temperature (or equivalently low energies), IIB string theory in this background is conjectured to equivalent to $`QCD_3`$. ### 3.1 Spin and Degeneracy of Glueball States Type IIB string theory at low energy has a supergravity multiplet with several zero mass bosonic fields: a graviton, $`G_{\mu \nu }`$, a dilaton $`\varphi `$, an axion (or zero form RR field) $`C`$ and two tensors, the NS-NS and RR fields $`B_{\mu \nu }`$ and $`C_{\mu \nu }`$ respectively. In addition there is the 4-form RR field $`C_{(4)}`$ that is constrained to have a self-dual field strength, $`F_{(5)}=dC_{(4)}`$. Now the task is to find all the quadratic fluctuations in the above background metric whose eigen-modes correspond to the discrete glueball spectra for $`QCD_3`$ at strong coupling. We are only interested in excitations that lie in the superselection sector for $`QCD_3`$. Thus for example we can ignore all non-trivial harmonic in $`S^5`$ that carry a non-zero R charge and all Kaluza-Klein (KK) modes in the $`S^1`$ thermal circle with a U(1) KK charge. The result of these considerations, discussed in detail below are summarized in Table 3. To count the number of independent fluctuations for a field of given spin, we again imagine harmonic plane waves propagating in the AdS radial direction, $`r`$, with Euclidean time, $`x_3`$. For example, the metric fluctuations in $`AdS^5`$ $$G_{\mu \nu }=\overline{g}_{\mu \nu }+h_{\mu \nu }(x),$$ (23) in the fixed background $`\overline{g}_{\mu \nu }`$ are taken to be of the form $`h_{\mu \nu }(r,x_3)`$. There is no dependence on the spatial co-ordinates, $`x_i=(x_1,x_2)`$, or the compactified “temperature” direction, $`\tau `$. #### 3.1.1 Metric fluctuations A graviton has two polarization indices. If we were in flat space time, we could go to a gauge where these indices took values only among $`(x_1,x_2,\tau )`$ and not from the set $`(r,x_3)`$. The polarization tensor should also be traceless. This leaves $`(3\times 4)/21=5`$ independent components. In the AdS space time, we can count the number of graviton modes the same way, though the actual modes that we construct will have this form of polarization only at $`r\mathrm{}`$; for finite $`r`$, other components of the polarization will be constrained to acquire nonzero values . Therefore, a set of five independent polarization tensors can be characterized by the following non-vanishing components at $`r\mathrm{}`$. In the $`AdS^5`$ black hole background, $`\tau `$ is compact, so the rotations group is $`SO(2)`$ in $`(x_1,x_2)`$ and the five states are in 3 irreducible representations: A spin-2 doublet (helicities $`\pm 2`$), spanned by $$G_{ij}:h_{12}=h_{21}0,\mathrm{and}h_{11}=h_{22}0,$$ (24) a spin-1 doublet (helicities $`\pm 1`$), spanned by $$G_{i\tau }:h_{\tau 1}=h_{1\tau }0,\mathrm{and}h_{\tau 2}=h_{2,\tau }0$$ (25) and a spin-0 state, $$G_{\tau \tau }:h_{\tau \tau }=2h_{11}=2h_{22}0.$$ (26) These fluctuations are denoted by $`T_3`$, $`V_3`$, and $`S_3`$ respectively in Table 3. #### 3.1.2 Two-form fields Each 2-form field in $`AdS^5`$ satisfies a field equation that includes a topological mass term. The 2-forms $`B_{\mu \nu }`$ and $`C_{\mu \nu }`$ can be combined into one complex 2-form field $`\stackrel{~}{B}_{\mu \nu }=B_{\mu \nu }+iC_{\mu \nu }`$. The field equation for $`\stackrel{~}{B}_{\mu \nu }`$ can be factorized into two first order equations, and each can be iterated leading to a second order equation of the form $$\mathrm{Max}\stackrel{~}{B}_{\mu \nu }+m_{AdS}^2\stackrel{~}{B}_{\mu \nu }=0,$$ where $`\mathrm{Max}`$ is the Maxwell operator on 2-forms. For modes that have no dependence on the $`S^5`$ co-ordinates, one equation is massive, with $`m_{AdS}^2=16`$, and the other is massless. It can be shown that the massless equation has only pure gauge solutions; so they can be ignored. (See Appendix A for further details.) For the purpose of counting modes, polarizations for a massless 2-form in $`AdS^5`$ can also be restricted to be transverse, with fields depending only on $`(x_3,r)`$. Therefore the polarization tensor is an antisymmetric 2-tensor in the directions $`x_1,x_2,\tau `$, leading to 3 independent components. On the other hand, for a general massive 2-form field, longitudinal polarizations are allowed, so that the polarization is an antisymmetric tensor in the coordinates $`x_1,x_2,\tau ,r`$, with 6 independent components. Nevertheless, the number of independent components for our massive 2-form $`\stackrel{~}{B}_{\mu \nu }`$ is only 3 (complex), as if we are dealing with a massless case. This is due to the fact that the second order equation above stems from a first order equation, $`ϵ_{\mu \nu }{}_{}{}^{\sigma \lambda \rho }_{[\sigma }^{}\stackrel{~}{B}_{\lambda \rho ]}+4i\stackrel{~}{B}_{\mu \nu }=0`$, relating real and imaginary parts and leading to additional constraints. For instance, with fields depending on $`(x_3,r)`$, if we start with the polarization tensor having $`B_{12}=B_{21}0,`$ then the first order equation will constrain us to have specific values for $`C_{\tau r}=C_{r\tau }0`$ and $`C_{3r}=C_{r3}0`$, while allowing all other components of $`B`$ and $`C`$ to be zero. Thus we count as independent fields $`B_{12},B_{1\tau },B_{2\tau },C_{12},C_{1\tau },C_{2\tau }`$, i.e., there are three independent solutions for $`B_{\mu \nu }`$ and three solutions for $`C_{\mu \nu }`$. The non-vanishing polarizations of the $`B_{\mu \nu }`$ tensor can be grouped as: $`B_{ij}:B_{12}=B_{21}`$ $``$ $`0,`$ (27) corresponding to spin-0 and $`B_{i\tau }:B_{1\tau }=B_{\tau 1}`$ $``$ $`0,`$ (28) $`B_{2\tau }=B_{\tau 2}`$ $``$ $`0.`$ (29) corresponding to a spin-1 doublet. In Table 3, these fluctuations are denoted by $`N_3`$ and $`M_3`$ respectively. They are degenerate with ones for the R-R 2-form $`C_{\mu \nu }`$. #### 3.1.3 Scalar fields In general, for each scalar field, there is a unique field equation with the plane wave dependence which we have been considering. There are three such scalar modes: The volume fluctuations $`G_\alpha ^\alpha `$ in $`S^5`$, (with $`m_{AdS}^2=32`$), denoted by $`L_3`$ , and fluctuations for the dilaton $`\varphi `$ and the axion $`C`$. However, as we have shown in an earlier paper , the latter two spectra are degenerate with the $`2^{++}`$ tensor fluctuations. Therefore, separate equations are not required . ### 3.2 Wave Equations and $`QCD_3`$ Glueball Spectrum Metric fluctuations for $`QCD_3`$ have been obtained previously by analyzing the linearized Einstein equations about the $`\mathrm{𝐀𝐝𝐒}^\mathrm{𝟓}\times 𝐒^\mathrm{𝟓}`$ black hole background which leads to three independent equations, $`T_3`$, $`V_3`$ and $`S_3`$ . Fluctuations $`N_3`$, $`M_3`$ and $`L_3`$ can be found similarly, leading to all together six independent equations for $`QCD_3`$. From the equation of motion, we determine the discrete spectrum with $`k_3=im`$. (See Appendix A for details.) The full set of independent equations are: $``$ $`{\displaystyle \frac{d}{dr}}(r^5r){\displaystyle \frac{d}{dr}}T_3(r)(m^2r)T_3(r)=0,`$ (30) $``$ $`{\displaystyle \frac{d}{dr}}(r^5r){\displaystyle \frac{d}{dr}}V_3(r)(m^2r{\displaystyle \frac{4}{r(r^41)}})V_3(r)=0.`$ $``$ $`{\displaystyle \frac{d}{dr}}(r^5r){\displaystyle \frac{d}{dr}}S_3(r)(m^2r+{\displaystyle \frac{64r^2}{(3r^41)^2}})S_3(r)=0,`$ $``$ $`{\displaystyle \frac{d}{dr}}(r^5r){\displaystyle \frac{d}{dr}}N_3(r)(m^2r12r^3+{\displaystyle \frac{4}{r}})N_3(r)=0,`$ $``$ $`{\displaystyle \frac{d}{dr}}(r^5r){\displaystyle \frac{d}{dr}}M_3(r)(m^2r12r^3{\displaystyle \frac{4r^3}{r^41}})M_3(r)=0,`$ $``$ $`{\displaystyle \frac{d}{dr}}(r^5r){\displaystyle \frac{d}{dr}}L_3(r)(m^2r32r^3)L_3(r)=0.`$ Each equation can be expressed in a variety of forms, depending on the choice of normalization. The following choices have been made so that each equation takes on a manifestly hermitian form: Tensor: $$h_{ij}=q_{ij}r^2T_3(r)e^{mx_3},$$ (31) where $`i,j=1,2`$, with $`q_{ij}`$ an arbitrary constant traceless-symmetric $`2\times 2`$ matrix. Vector: $$h_{i\tau }=q_i\sqrt{r^41}V_3(r)e^{mx_3},$$ (32) where $`q_i`$ is a constant 2-vector. Scalar: This case has been treated carefully using several different gauge choices . Here we adopt the form suggested by Constable and Meyer with $$h_{\tau \tau }=(r^2r^2)S_3(r)e^{mx_3}.$$ (33) Two-forms and volume scalar: For $`B_{12}`$, consider perturbations of the form $$B_{12}=r^2N_3(r)e^{ik_3x_3},$$ (34) with $`B_{1\tau }=B_{2\tau }=0`$. Alternatively, we consider $$B_{i\tau }=q_i\sqrt{r^41}M_3(r)e^{mx_3},$$ (35) with $`B_{12}=0`$, where $`q_i`$ is an arbitrary constant 2-vector. Lastly, for the volume perturbation, we consider $$h_\alpha ^\alpha =L_3(r)e^{mx_3}.$$ To calculate the discrete spectrum for each of these equations one must again apply the correct boundary conditions at $`r=1`$ and $`r=\mathrm{}`$ as mentioned in the case of $`QCD_4`$ earlier. (See Appendix B for a listing of boundary conditions.) The resultant spectrum is given in Table 4. Similarly, we have also compared our numerical $`QCD_3`$ masses to the WKB approximations, $$m_n^2\mu _3^2(n^2+\delta n+\gamma )+0(\frac{1}{n}),$$ (36) where $`\mu _3^2=16\pi (\mathrm{\Gamma }(3/4)/\mathrm{\Gamma }(1/4))^2`$, with integer constants, $`\delta `$, listed in Table 6, determined analytically. Again, the constants $`\gamma `$ were fits to the numerical data. (See Appendix B). ## 4 Parity and Charge Conjugation assignments Next we determine how the supergravity fields and therefore the glueballs couple to the boundary gauge theory. This allows us to unambiguously assign the correct parity and charge quantum numbers to the glueball states. For this purpose we consider the effective Born-Infeld action on the branes. ### 4.1 $`QCD_4`$ The 4-d gauge theory for $`QCD_4`$ is obtained by dimensional reduction from a 5-d gauge theory, which is the low energy dynamics of D4-branes in 10-dimensional Type IIA string theory. Although this 10-d theory may itself be regarded as a dimensional reduction of 11-d M-theory for membranes, it is sufficient and more convenient to consider the 10-d theory itself. Since supergravity fields can be thought of as coupling constants for gauge theory operators, their quantum numbers can be assigned by the parity and charge conjugation invariance of the overall action, (supergravity field times composite operator). For simplicity let us consider the coupling of a supergravity field to just one D4-brane — this coupling is given by a Born-Infeld action plus a Wess-Zumino term, $$S=d^5xdet[G_{\mu \nu }+e^{\varphi /2}(B_{\mu \nu }+F_{\mu \nu })]+d^4x(C_1FF+C_3F+C_5),$$ (37) where $`\mu ,\nu =1,2,3,4,\tau `$. Later we will argue that our quantum number assignment is correct also for the non-abelian case of N coincident D branes. In the 5-d field theory, we have the space-time world volume co-ordinate $`x_1,x_2,x_3,x_4,\tau `$ with $`\tau `$ compactified on $`S^1`$. The Euclidean time coordinate we take to be $`x_4`$. After dimensional reduction the physical fields will be characterized by their representation under the little group $`SO(3)`$ of rotations on the spatial co-ordinates $`x_i`$, $`i=1,2,3`$, in the 4-d theory. For the 5-d gauge fields, we define parity by $`P`$ $`:`$ $`A_i(x_i,x_4,\tau )A_i(x_i,x_4,\tau ),`$ (38) $`P`$ $`:`$ $`A_4(x_i,x_4,\tau )A_4(x_i,x_4,\tau ),`$ (39) $`P`$ $`:`$ $`A_\tau (x_i,x_4,\tau )A_\tau (x_i,x_4,\tau ),`$ (40) for $`x_ix_i`$, $`x_4x_4`$, and $`\tau \tau `$. For the Euclidean $`𝐑^\mathrm{𝟓}`$ space, this is the only discrete symmetry. However after compactification to $`𝐑^\mathrm{𝟒}\times 𝐒^\mathrm{𝟏}`$, we can define another parity (not related by a 5-d proper Lorentz transformation) by inverting the $`\tau `$ co-ordinate on $`𝐒^\mathrm{𝟏}`$. Thus we define a separate discrete $`\tau `$parity transformation $`P_\tau :`$ $`\tau \tau `$, $`P_\tau `$ $`:`$ $`A_i(x_i,x_4,\tau )A_i(x_i,x_4,\tau ),`$ (41) $`P_\tau `$ $`:`$ $`A_4(x_i,x_4,\tau )A_4(x_i,x_4,\tau ),`$ (42) $`P_\tau `$ $`:`$ $`A_\tau (x_i,x_4,\tau )A_\tau (x_i,x_4,\tau ).`$ (43) Charge conjugation for a non-abelian gluon field is $$C:{\scriptscriptstyle \frac{1}{2}}T_aA_\mu ^a(x){\scriptscriptstyle \frac{1}{2}}T_a^{}A_\mu ^a(x)$$ (44) where $`T^a`$ are the Hermitian generators of the group. In terms of matrix fields ($`A{\scriptscriptstyle \frac{1}{2}}T_aA^a`$), $`C:A_\mu (x)A_\mu ^T(x).`$ This leads to a subtlety. For example consider the transformation of a trilinear gauge invariant operators, $$C:Tr[F_{\mu _1\nu _1}F_{\mu _2\nu _2}F_{\mu _3\nu _3}]Tr[F_{\mu _3\nu _3}F_{\mu _2\nu _2}F_{\mu _1\nu _1}].$$ (45) The order of the fields is reversed. Hence the symmetric products, $`d^{abc}F_1^aF_2^bF_3^c`$, have $`C=1`$ and the antisymmetric products, $`f^{abc}F_1^aF_2^bF_3^c`$, $`C=+1`$. Of course using a single brane, we can only find symmetric products. For reasons explained further in Sec. 5, we will only encounter symmetric traces over polynomials in F, designate by $`Sym\mathrm{Tr}[F_{\mu \nu }\mathrm{}]`$. Even polynomials have $`C=+1`$ and odd polynomials $`C=1`$. #### 4.1.1 Graviton couplings Expanding the Born-Infeld action, we can now read off the $`J^{PC}(P_\tau )`$ assignments: ¿From the coupling, $`G_{\mu \nu }T^{\mu \nu }G_{\mu \nu }\mathrm{Tr}(F_{\mu \lambda }F_\nu ^\lambda )+\mathrm{},`$ we obtain $$G_{ij}2^{++}(P_\tau =+),G_{i\tau }1^+(P_\tau =+)G_{\tau \tau }0^{++}(P_\tau =+).$$ (46) Under compactification of 11-d supergravity theory, $`G_{\mu ,11}`$ becomes the Ramond-Ramond 1-form $`C_\mu `$, which couples as $`ϵ^{\mu \nu \lambda \kappa \eta }C_\mu Sym\mathrm{Tr}[F_{\nu \lambda }F_{\kappa \eta }W],`$ where $`W`$ is an an even power of fields $`F`$. Consequently, the coupling, $`ϵ^{ijk}C_iSym\mathrm{Tr}[F_{\tau j}F_{k4}W]+\mathrm{}`$ leads to $$C_i1^{++}(P_\tau =).$$ (47) Similarly, $`ϵ^{ijk}C_\tau \mathrm{Tr}(F_{ij}F_{4k}W)+\mathrm{}`$ gives $$C_\tau 0^+(P_\tau =+),$$ (48) and $`G_{11,11}`$ leads to the dilaton $`\varphi `$ with coupling $`\varphi \mathrm{Tr}F^2,`$ $$\varphi 0^{++}(P_\tau =+).$$ (49) #### 4.1.2 Two-form, three-form fields, and volume scalar Consider first the NS-NS 2-form field $`B_{\mu \nu }`$. This field in the Type IIA theory arises from the 3-form field of the 11-d supergravity theory when the components of the 3-form field are $`A_{\mu \nu ,11}`$. For $`U(1)`$ gauge theory in leading order this field couples as $`B_{\mu \nu }F^{\mu \nu }`$. More generally in the $`SU(N)`$ gauge theory, we must have a multi-gluon coupling, $`B_{\mu \nu }Sym\mathrm{Tr}[F_{\mu \nu }W]`$. (Again $`W`$ is an an even power of fields $`F`$ and the trace is symmetrized.) To determine the parity, assume for the supergravity modes that we are in a gauge where the indices of the 2-form do not point along $`x_4,x_{11},r`$. With $`i,j=1,2,3`$, this leads to coupling $`B_{ij}Sym\mathrm{Tr}[F^{ij}W]`$ with $$B_{ij}1^+(P_\tau =+),$$ (50) and coupling $`B_{i\tau }Sym\mathrm{Tr}[F^{i\tau }W]`$ with $$B_{i\tau }1^{}(P_\tau =).$$ (51) An analogous analysis can also be carried out for 3-form fields, $`C_{ijk}`$ and $`C_{ij\tau }`$. The coupling, $`C_{123}Sym\mathrm{Tr}[F^{4\tau }W]`$ leads to $$C_{ijk}0^+(P_\tau =),$$ (52) and coupling $`ϵ^{ijk}C_{ij\tau }Sym\mathrm{Tr}[F^{4k}W]`$ leads to $$C_{ij\tau }1^{}(P_\tau =+).$$ (53) Lastly, the volume scalar couples as $`h_\alpha ^\alpha Sym\mathrm{Tr}F^4+\mathrm{}`$ giving $$h_\alpha ^\alpha 0^{++}(P_\tau =+).$$ (54) The complete parity and charge conjugation assignments are given in Table 1. ### 4.2 $`QCD_3`$ The 3-d gauge theory is obtained by dimensional reduction from a 4-d gauge theory. To find the symmetries of the interactions, we consider the Born-Infeld action plus Wess-Zumino term, describing the coupling of a supergravity field to a single D3-brane, $$S=d^4xdet[G_{\mu \nu }+e^{\varphi /2}(B_{\mu \nu }+F_{\mu \nu })]+d^4x(C_0FF+C_2F+C_4),$$ where $`\mu ,\nu =1,2,3,\tau `$. As in the case of $`QCD_4`$, we will find the charge conjugation and parity assignments with the help of the symmetries of the 4-d gauge theory and then take the dimensional reduction to the 3-d theory after compactification of the coordinate $`\tau `$. The Euclidean time is taken to be $`x_3`$ and the spatial co-ordinates, $`x_i`$, $`i=1,2`$. After dimensional reduction the physical fields will be characterized by their representation under the little group of the space co-ordinates of the 3-d theory: this is the group $`SO(2)`$ rotations in the $`x_1,x_2`$ plane. However, unlike the case of $`QCD_4`$, the usual spatial inversion, $`x_ix_i`$ with $`x_3x_3`$ and $`\tau \tau `$, is a rotation $`R(\theta )`$ in $`SO(2)`$ with $`\theta =\pi `$; it therefore does not lead to a discrete symmetry. A discrete symmetry can be defined by $`x_1x_1`$ and $`x_2x_2`$, as was pointed out in the lattice studies by Harte and Philipsen . However, the sole manifestation of this symmetry is the helicity “doublets”, $`\lambda =\pm J`$, for states with spin $`J>0`$. We have already taken this degeneracy into account. On the other hand, $`\tau `$-parity remains a discrete symmetry of the action, as is the case for $`QCD_4`$. In this paper, we shall define “parity” for $`QCD_3`$ as $`PR(\pi )\times P_\tau `$ where $`x_ix_i`$, $`x_3x_3`$, and $`\tau \tau `$. This of course is precisely the parity for the uncompactified 4-d theory with $`x_3`$ treated as Euclidean time. #### 4.2.1 Graviton, dilaton and axion states The graviton $`G_{\mu \nu }`$ couples as $`G_{\mu \nu }T^{\mu \nu }`$ as in the case of $`QCD_4`$. Because an even number of gluons occur in the field operators, the charge conjugation for all such states are $`C=+`$. For parity, we assume we are in a gauge where the indices of $`G_{\mu \nu }`$ do not point along $`x_3,r`$. From the coupling, $`G_{\mu \nu }\mathrm{Tr}[F^{\mu \lambda }F_\lambda ^\nu ]+\mathrm{}`$, we get states $$G_{ij}2^{++}(P_\tau =+),G_{i\tau }1^{++}(P_\tau =),G_{\tau \tau }0^{++}(P_\tau =+).$$ (55) The dilaton couples as $`\varphi \mathrm{Tr}F^2`$, leading to $$\varphi 0^{++}(P_\tau =+),$$ (56) and for an axion coupling $`C_0\mathrm{Tr}(F_{12}F_{\tau 3})`$, $$C_00^+(P_\tau =).$$ (57) #### 4.2.2 Two-form fields and volume scalar Consider first the NS-NS 2-form field $`B_{\mu \nu }`$. For parity, again assume that we are in a gauge where the indices of the 2-form do not point along $`x_3,r`$. With $`i,j=1,2`$, the coupling $`B_{ij}Sym\mathrm{Tr}[F^{ij}W]`$ leads to $$B_{ij}0^+(P_\tau =+),$$ (58) and $`B_{i\tau }Sym\mathrm{Tr}[F^{i\tau }W]`$ to $$B_{i\tau }1^+(P_\tau =).$$ (59) Finally for the Ramond-Ramond 2-form $`C_{\mu \nu }`$, we have the coupling $`C_{12}Sym\mathrm{Tr}[F_{\tau 3}W]`$, so $$C_{12}0^{}(P_\tau =),$$ (60) and $`ϵ^{ij}C_{\tau i}Sym\mathrm{Tr}[F_{j3}W]`$ giving $$C_{i\tau }1^{}(P_\tau =+).$$ (61) Finally, as in the case of $`QCD_4`$, the volume scalar couples as $`h_\alpha ^\alpha \mathrm{Tr}F^4+\mathrm{}`$ so that $$h_\alpha ^\alpha 0^{++}(P_\tau =+).$$ (62) The complete parity and charge conjugation assignments are given in Table 3. ## 5 Discussion Lastly we turn to the question of how well the strong coupling limit for the Maldacena dual theory of QCD represents the infrared physics probed by the glueball spectra. Happily we now have a rather definitive lattice glueball spectrum by Morningstar and Pearson with which to make comparisons (See right side of Fig. 2 below). ### 5.1 Comparison with Lattice Glueball Spectrum Originally, claims were made about accurate comparisons to a few percent for isolated (scalar) mass ratios. As we pointed out in Ref. for $`QCD_3`$, the lowest mass scalar comes from the gravitational multiplet, not the dilaton. A similar spectrum is observed for $`QCD_4`$. Consequently such accurate mass ratios were a misconception. This should not be regarded as a failure, since any reasonable expectation of a strong coupling approximation should not give quantitative results. On the other hand, there is a rather remarkable correspondence of the overall mass and spin structure between our strong coupling glueball spectrum and the lattice results at weak coupling for $`QCD_4`$ (see Fig. 2 below.) Apparently the spin structure of type IIA supergravity does resemble the low mass glueball spin splitting. The correspondence is sufficient to suggest that the Maldacena duality conjecture may well be correct and that further efforts to go beyond strong coupling are worthy of sustained effort. Note that for each value of $`PC=(++,+,+,)`$, the lowest state is present in approximately the right mass range. In addition, the exact $`2^{++}/0^{++}`$ degeneracy for AdS strong coupling corresponds to a relatively small splitting in the lattice calculations. Finally, there is a radial excitation of the pseudoscalar $`0^+`$ that suggests that even this effect is approximated. This is intriguing because in the supergravity description the radial mode is a standing wave in the extra 5th dimension whereas in the lattice it is a conventional radial mode. Apparently scale changes in 4-d are being represented by the distance into the extra “warped” 5th axis. At higher masses the discrepancies increase. One reason is the obvious fact that on the supergravity side all orbital excitations of higher spin states are pushed to infinity in strong coupling by virtue of the divergent string tension, $$\sigma \frac{1}{2\pi \alpha ^{}}=\frac{16\pi g^2N}{27}[\mathrm{\hspace{0.33em}1}+0(\frac{1}{g^2N})].$$ (63) For example, the $`3^{++}`$ state is a purely stringy effect outside of the classical limit of supergravity. Finally, we must emphasize that our comparison is premised on the neglect of many “spurious” states in the strong coupling limit that are in the wrong superselection sector to survive in the conjectured weak coupling limit of QCD. For example, all the Kaluza-Klein modes in the compact thermal $`S^1`$ manifold and the sphere $`S^4`$ (or $`S^5`$) have masses at the cut-off scale. (The first mode on the thermal circle has a KK mass scale $`m_{KK}^2=4`$ in type IIB and $`m_{KK}^2=9`$ in type IIA in the units used in Table 4 and 2 respectively.) But these spurious KK modes all carry conserved $`U(1)`$ or $`R`$ charges that are absent in the target theory. We assume they will disappear in the continuum limit. A subtler situation occurs in the $`QCD_4`$ example. Because normal modes in the extreme strong coupling limit do not distinguish between the compact 11th dimension and the spatial co-ordinates $`x_1,x_2,x_3`$ on the brane, the spectrum actually has an exact $`SO(4)`$ symmetry. Thus there are additional states (listed in Table 1 but ignored in Fig. 2) exactly degenerate with the physically reasonable $`2^{++}/0^{++},\mathrm{\hspace{0.33em}0}^+,\mathrm{\hspace{0.33em}1}^+`$ and $`1^{}`$ states. They are all odd under the discrete symmetry of reflecting the thermal circle (i.e. $`P_\tau =`$) so they also lie in another superselection sector. A major challenge is to understand how this $`SO(4)`$ symmetry is lifted and if the unwanted states remain at the cut-off in weak coupling. Physically the 11th axis is very different. The membranes of 11-d M-theory wraps this axis. Another possibility worth exploring is modifying the background metric with an orbifold that projects directly onto the even $`\tau `$ parity sector for QCD. ### 5.2 Constituent Gluon Picture The basic idea behind the AdS/CFT correspondence in the context of glueballs is similar to an observation made much earlier by Fritzsch and Minkowski , by Bjorken and by Jaffe, Johnson and Ryzak . Namely that the low mass glueball spectrum can be qualitatively understood in terms of local gluon interpolating operators of minimal dimension. For example, Ref. lists all gauge invariant operators for dimension $`d=4,5`$ and $`6`$. Eliminating operators that are zero by the classical equation of motion and states that decouple because of the conservation of the energy momentum tensor, the operators are in rough correspondence with all the low mass glueballs states, as computed in a constituent gluon or bag model. Indeed more recently Kuti has pointed out that a more careful use of the spherical cavity approximation even gives a rather good quantitative match to the lowest 11 states in the lattice spectrum. Consequently it is interesting to compare this set of operators with the supergravity model. We list below all the operators for $`d6`$ except the operators with explicit derivatives (e.g. $`Tr[FDF]`$) and $`Tr[FDDF]`$ ): \| Dimension \| State \| Operator \| Supergravity \| \| --- \| --- \| --- \| --- \| \| d = 4 \| $`0^{++}`$ \| $`Tr(FF)=\stackrel{}{E}^a\stackrel{}{E}^a\stackrel{}{B}^a\stackrel{}{B}^a`$ \| $`\varphi `$ \| \| d = 4 \| $`2^{++}`$ \| $`T_{ij}=E_i^aE_j^a+B_i^aB_j^a\text{trace}`$ \| $`G_{ij}`$ \| \| d = 4 \| $`0^+`$ \| $`Tr(F\stackrel{~}{F})=\stackrel{}{E}^a\stackrel{}{B}^a`$ \| $`C_\tau `$ \| \| d = 4 \| $`0^{++}`$ \| $`2T_{00}=\stackrel{}{E}^a\stackrel{}{E}^a+\stackrel{}{B}^a\stackrel{}{B}^a`$ \| $`G_{\tau \tau }`$ \| \| d = 4 \| $`2^+`$ \| $`E_i^aB_j^a+B_i^aE_j^a\text{trace}`$ \| Absent \| \| d = 4 \| $`2^{++}`$ \| $`E_i^aE_j^aB_i^aB_j^a\text{trace}`$ \| Absent \| \| d = 6 \| $`(1,2,3)^\pm `$ \| $`Tr(F_{\mu \nu }\{F_{\rho \sigma },F_{\lambda \eta }\})d^{abc}F^aF^bF^c`$ \| $`B_{ij},C_{ij\tau }`$ \| \| d = 6 \| $`(0,1,2)^{\pm +}`$ \| $`Tr(F_{\mu \nu }[F_{\rho \sigma },F_{\lambda \eta }])f^{abc}F^aF^bF^c`$ \| Absent \| In this table we have used a Minkowski metric. The classification is parallel to our discussion of couplings in the Born-Infeld action for D4 branes (see Sec. 4.1), except that all components and derivatives in the 5th, i.e., $`\tau `$, direction are zero and therefore all $`P_\tau =1`$ states are absent (See Table 1). The column on the right lists the supergravity mode that couples (after dropping $`\tau `$ components) to each operator. Several observations are in order. For the d = 4 operators, there is complete agreement on the 3 lowest quantum numbers: $`J^{PC}=0^{++},2^{++},0^+`$. Also there is agreement on the absence of a low mass $`1^+`$ state that Ref. attributes to the conservation law that decouples the operator corresponding to the momentum tensor $`T_{0i}=\stackrel{}{E}^a\times \stackrel{}{B}^a`$. In this context, it is worth commenting on the two independent sources for $`0^{++}`$ states — the condensate, $`\mathrm{Tr}(FF)`$, and the energy density, $`T_{00}={\scriptscriptstyle \frac{1}{2}}(\stackrel{}{E}^a\stackrel{}{E}^a+\stackrel{}{B}^a\stackrel{}{B}^a)`$. Naively one might conclude the second one coming from the conserved energy momentum tensor should be dropped, for the same reason we dropped the operator for momentum conservation, $`T_{0i}`$. However in fact for QCD because of the conformal anomaly for the trace of the energy momentum tensor it is easy to show that the decoupling argument fails. A bag-like model circumvents the decoupling because the bag itself implies a scale breaking “vacuum” (empty bag) thus introducing an extra four-vector, the bag velocity $`u_\mu `$. Our AdS black hole background, which is a key ingredient in our approach, also breaks conformal invariance. Consequently all agree that there is an extra low mass $`0^{++}`$ state in addition to the one which in our case is degenerate with the tensor $`2^{++}`$. The lattice data clearly favors this low mass $`0^{++}`$ state, in agreement with our AdS spectrum. Finally, one rather low mass state, the $`2^+`$, is missing in the AdS spectrum. This state is clearly present in the lattice spectrum and is identified in the bag model. At d = 6, we have two states identified in the C = -1 symmetric trace, $`Tr(F_{\mu \nu }\{F_{\rho \sigma },F_{\lambda \eta }\})`$: the $`1^+`$ state for operator $`d^{abc}\stackrel{}{B}_a(\stackrel{}{E}_b\stackrel{}{E}_c)`$ and the $`1^{}`$ state for the operator $`d^{abc}\stackrel{}{E}_a(\stackrel{}{E}_b\stackrel{}{E}_c)`$. These states are clearly related to the field content of a IIA supersymmetric multiplet. The higher spin representation are not present at strong coupling. Moreover, we have no states corresponding to the antisymmetric trace operators, $`Tr(F[F,F])`$. It appears that in the limit where we restrict to supergravity modes and ignore the massive stringy states, we will be able to obtain only the glueball state with the symmetric $`d^{abc}`$ coupling between the group indices, and not the state with the antisymmetric $`f^{abc}`$ coupling. If we consider the chiral primaries of the 4-d Yang-Mills theory (which was the theory on the boundary before we took the $`\tau `$ direction to be compact), then we find that these had the form $`Sym\mathrm{Tr}(X^iX^jX^k)`$ \- i.e, we have a symmetric trace over the fields. Other operators that couple to the supergravity fields will be supersymmetry descendents of these chiral primaries, but the symmetry in the trace would be maintained<sup>1</sup><sup>1</sup>1We thank W. Taylor for a discussion on this point.. This fact may be related to the observation of Tseytlin that generalizing the Born-Infeld action to a non-abelian case gives rise to symmetric trace operators in the field theory. One legitimate point of view is simply to suppose that all missing states must by definition be stringy effects that will be restored in weak coupling. However we prefer to look on this as a possible clue on constructing a better initial geometry for the supergravity/QCD duality proposals. ### 5.3 Strong coupling Expansion for Pomeron Intercept We shall end this discussion with a comment on the slope of the leading glueball trajectory as way to estimate the crossover value for the bare coupling, where continuum physics might begin to hold. The Pomeron is the leading Regge trajectory passing through the lightest glueball state with $`J^{PC}=2^{++}`$. In a linear approximation, it can be parameterized by $$\alpha _P(t)=2+\alpha _P^{}(tm_T^2),$$ (64) where we can use the strong coupling estimate for the lightest tensor mass<sup>2</sup><sup>2</sup>2 We have adopted the normalization in the $`AdS`$-black hole metric to simplify the coefficients, e.g., for $`AdS^7`$, $`\overline{g}_{\tau \tau }=r^2r^4`$. This corresponds to fixing the “thermal-radius” $`R_1=1/3`$ so that $`\beta =2\pi R_1=2\pi /3`$., $$m_T[9.86+0(\frac{1}{g^2N})]\beta ^1.$$ (65) Moreover if we make the standard assumption that the closed string tension is twice that between two static quark sources , we also have a strong coupling expression for the Pomeron slope, $$\alpha _P^{}[\frac{27}{32\pi g^2N}+0(\frac{1}{g^4N^2})]\beta ^2.$$ (66) Putting these together, we obtain a strong coupling expansion for the Pomeron intercept, $$\alpha _P(0)20.66(\frac{4\pi }{g^2N})+0(\frac{1}{g^4N^2}).$$ (67) Turning this argument around, we can estimate a crossover value between the strong and weak coupling regimes by fixing $`\alpha _P(0)1.2`$ at its phenomenological value . In fact this yields for $`QCD_4`$ at $`N=3`$ a reasonable value for $`\alpha _{strong}=g^2/4\pi =0.176`$ for the crossover. Much more experience with this new approach to strong coupling must be gained before such numerology can be taken seriously. However, similar crude argument have proven to be a useful guide in the crossover regime of lattice QCD. One might even follow the general strategy used in the lattice cut-off formulations. Postpone the difficult question of analytically solving the QCD string to find the true UV fixed point. Instead work at a fixed but physically reasonable cut-off scale (or bare coupling) to calculate the spectrum. If one is near enough to the fixed point, mass ratios should be reliable. After all, the real benefit of a weak/strong duality is to use each method in the domain where it provides the natural language. On the other hand, clearly from a fundamental point of view, finding analytical tools to understand the renormalized trajectory and prove asymptotic scaling within the context of the gauge invariant QCD string would also be a major achievement — an achievement that presumably would include a proof of confinement itself. Acknowledgments: We would like to acknowledge useful conversations with R. Jaffe, A. Jevicki, D. Lowe, J. Kuti, J. Minahan, J. M. Maldacena, H. Ooguri, W. Taylor and U-J Wiese. ## Appendix A Wave Equations In this appendix we outline the derivation of the wave equations that were used to find the energy levels in the supergravity theory. First we take the case of $`QCD_3`$, for which we have the metric, $$ds^2=(r^2\frac{1}{r^2})d\tau ^2+r^2\underset{i=1,2,3}{}dx_i^2+(r^2\frac{1}{r^2})^1dr^2+d\mathrm{\Omega }_5^2,$$ (A.1) where $`x_3`$ is the Euclidean time direction. The simplest equation is the scalar wave equation for the dilaton and the axion. At the linear perturbation level, both satisfy $$\varphi _{,\mu }{}_{}{}^{;\mu }=\frac{1}{\sqrt{g}}[\varphi _{,\mu }g^{\mu \nu }\sqrt{g}]_{,\nu }=0.$$ (A.2) We introduce a plane wave ansatz, $$\varphi =T_3(r)e^{ik_3x_3},$$ (A.3) with zero momentum and mass, $`m=ik_3`$ providing the equation (8) for $`T_3`$ in the text. Fluctuations in the volume of the sphere $`S^5`$, which is a fiber at each point of space time, provides another scalar mode. This scalar has an AdS mass squared equal to $`32`$. The field equation is $$\frac{1}{\sqrt{g}}[\varphi _{,\mu }g^{\mu \nu }\sqrt{g}]_{,\nu }32\varphi =0,$$ (A.4) $`\varphi =L_3(r)e^{ik_3x_3}`$, which reduces to the equation (8) for $`L_3`$ in the text. Next let us address the case of the two-form fields, $`B_{\mu \nu }`$ and $`C_{\mu \nu }`$, which at the linear level are conveniently combined into a single complex field, $`\stackrel{~}{B}_{\mu \nu }=B_{\mu \nu }+iC_{\mu \nu }`$. It was shown in that this field satisfies the equation, $$(\mathrm{Max}k(k+4))\stackrel{~}{B}_{\mu \nu }+2iϵ_{\mu \nu }{}_{}{}^{\rho \lambda \sigma }_{\rho }^{}\stackrel{~}{B}_{\lambda \sigma }=0,$$ (A.5) for $`k=0,1,\mathrm{}`$ harmonics on $`S^5`$. The Maxwell operator is defined by $$\mathrm{Max}\stackrel{~}{B}_{\mu \nu }=H_{\mu \nu \lambda }{}_{}{}^{;\lambda }.$$ (A.6) in terms of the field strength, $$H_{\mu \nu \lambda }=_\mu \stackrel{~}{B}_{\nu \lambda }+_\nu \stackrel{~}{B}_{\lambda \mu }+_\lambda \stackrel{~}{B}_{\mu \nu }.$$ (A.7) Since the second order differential operator factorizes into two first order operators, solutions fall into two classes, $$(2kI+iD)\stackrel{~}{B}_{\mu \nu }^{(1)}=0,(2(k+4)IiD)\stackrel{~}{B}_{\mu \nu }^{(2)}=0,$$ (A.8) where $`(DB)_{\mu \nu }=ϵ_{\mu \nu }{}_{}{}^{\rho \lambda \sigma }_{\rho }^{}\stackrel{~}{B}_{\lambda \sigma }`$ and $`I`$ is the identity matrix. It is convenient to iterate these first order equations to get the second order equations, $$(\mathrm{Max}k^2)\stackrel{~}{B}_{\mu \nu }^{(1)}=0,(\mathrm{Max}(k+4)^2)\stackrel{~}{B}_{\mu \nu }^{(2)}=0.$$ (A.9) We are interested in fields with no dependence on the coordinates of the sphere, so we can take $`k=0`$. It can be shown that the first class of solutions, $`\stackrel{~}{B}_{\mu \nu }^{(1)}`$, are pure gauge, so we are only interested in the second class, which has an effective mass squared of $`16`$ for the field $`\stackrel{~}{B}_{\mu \nu }^{(2)}`$. As explained in the text, one must of course check that solutions to the second order equation for $`\stackrel{~}{B}_{\mu \nu }^{(2)}`$, actually are valid solution to the original wave equation. This reduce the number of independent tensor fields, $`\stackrel{~}{B}_{\mu \nu }`$, from 6 to 3. For example with the ansatz, $$\stackrel{~}{B}_{12}=N_3(r)r^2e^{ik_3x_3}$$ (A.10) the first order equations will determine $`\stackrel{~}{B}_{\tau 3}`$ and $`\stackrel{~}{B}_{r\tau }`$ once we have a solution of the second order equation for $`\stackrel{~}{B}_{12}`$. This does not place any constraints on the solution for $`\stackrel{~}{B}_{12}`$ itself. We have defined $`\stackrel{~}{B}_{12}`$ in terms of the normalized coefficient $`N_3(r)`$ to obtain a hermitian field equation for $`N_3`$ similar to our earlier scalar mode $`T_3`$. In a similar manner we can find the equation for $`\stackrel{~}{B}_{1\tau }`$. We can solve the second order equation for this fluctuation without constraint, and then the requirement arising from the associated first order equations determines corresponding values of $`\stackrel{~}{B}_{23}`$ and $`\stackrel{~}{B}_{2r}`$. After adopting the ansatz indicated in the main text, we again obtain a hermitian equation for $`M_3(r)`$. The graviton perturbations arise from the Einstein action with a cosmological constant, expanded around the given background. We write $`G_{\mu \nu }=\overline{g}_{\mu \nu }+h_{\mu \nu }`$. the equation for the perturbation $`h_{\mu \nu }`$ is $$\frac{1}{2}h_{\mu \nu ;\lambda }{}_{}{}^{;\lambda }\frac{1}{2}h_{\lambda ;\mu \nu }^\lambda +\frac{1}{2}h_{\mu \lambda ;\nu }{}_{}{}^{;\lambda }+\frac{1}{2}h_{\nu \lambda ;\mu }{}_{}{}^{;\lambda }+4h_{\mu \nu }=0$$ (A.11) Near the boundary at $`r=\mathrm{}`$, we can choose a gauge to make the perturbations transverse to $`r,x_3`$ and traceless. It turns out that we can maintain this condition for all $`r`$ for perturbations of the form $`h_{12}`$, and for perturbations of the form $`h_{1\tau }`$. In these cases the above equation for $`h_{\mu \nu }`$ gives immediately the wave equations to be solved. But keeping in mind the decomposition in spin eigenstates in the $`x_1x_2`$ plane, we also find that we have to consider a spin-0 perturbation which at infinity has the form $`h_{\tau \tau }`$, with $`h_{11}=h_{22}=\frac{1}{2}h_{\tau \tau }`$. In this case we can choose a gauge to make $`h_{3\mu }=0`$, but for finite $`r`$ we will find in this gauge that $`h_{rr}0`$ and also that the part transverse to $`r,x_3`$ is not traceless. Thus we have to keep $`h_{\tau \tau },h_{11}=h_{22},h_{rr}`$ as independent coupled functions in the analysis. It was shown in how these equations can be reduced to one effective equation which can then be solved in the same way as the equations for the other fields. A nicer choice of gauge was used in which led to an equivalent but simpler equation. We will use the latter source for the equations, especially since the results there include all dimensions, and so can be used for the case of $`QCD_4`$ as well. Now we turn to the case of $`\mathrm{𝐀𝐝𝐒}^\mathrm{𝟕}\times 𝐒^\mathrm{𝟒}`$, which is very similar. The metric is $$ds^2=(r^2\frac{1}{r^4})d\tau ^2+r^2\underset{i=1,2,3,4,11}{}dx_i^2+(r^2\frac{1}{r^4})^1dr^2+\frac{1}{4}d\mathrm{\Omega }_4^2,$$ (A.12) where we have $`x_4`$ as the time direction and we note that the radius of the $`S^4`$ is half the curvature radius of the AdS space: this will affect the masses arising from the deformations of the sphere. There is a three-form field $`A_{\mu \nu \lambda }`$ which behaves in a manner similar to the two-form field $`\stackrel{~}{B}_{\mu \nu }`$ discussed above. Its field equation can again be factorized into two first order equations, which we iterated to second order equations. One factor at $`k=0`$ corresponds to pure gauge, while the other at $`k=0`$ has the value $`m_{AdS}^2=4(k+3)^2=36`$. For the scalar mode due to fluctuations of the volume of the $`S^4`$, we get a $`m_{AdS}^2=4\times 18=72`$. For these considerations we arrive at the wave equations (22) given in the text. Finally for comparison with P. van Nieuwenhuizen Ref , we note that they have scaled the radius of the $`S^4`$ to unity, instead of $`{\scriptscriptstyle \frac{1}{2}}`$. Our choice was made to keep the radius of the $`AdS^7`$ equal to unity, to make the comparison between $`AdS^5`$ and $`AdS^7`$ more natural. This change in metric scales the squared masses by a factor of $`4`$ relative to Ref. . ## Appendix B WKB and Variational Estimates First we change variable from $`r`$ to $`xr^2`$ and express all twelve equations in the standard Sturm-Liouville form, $$\{\frac{d}{dx}\tau (x)\frac{d}{dx}+w(x)\}\varphi _n(x)=m_n^2\sigma (x)\varphi _n(x),$$ where $`\tau (x)`$, $`w(x)`$, and $`\sigma (x)`$ are generalized “tension”, “external force”, and “mass-density” respectively. In our case, we have for $`QCD_4`$: $`\tau _4(x)=(x^4x)`$ and $`\sigma _4(x)=\frac{x}{4}`$; for $`QCD_3`$: $`\tau _3(x)=(x^3x)`$ and $`\sigma _3(x)=\frac{1}{4}`$. Force densities $`w(x)`$ for all 12 cases are listed in Table 5. We shall use this as our starting point for carrying out variational and WKB analyses. ### B.1 Variational Estimates for $`m_0^2`$: Solving for eigenstates, $`\{\varphi _n(x)\}`$ and their corresponding eigenvalues, $`\{m_n^2\}`$, is equivalent to finding stationery points of the following functional, $$\mathrm{\Gamma }[\varphi ]\frac{_1^{\mathrm{}}𝑑x[\tau (x)\varphi ^{}(x)^2+w(x)\varphi (x)^2]}{_1^{\mathrm{}}𝑑x\sigma (x)\varphi (x)^2},$$ (B.1) with $`m_n^2=\mathrm{\Gamma }[\varphi _n]`$. In particular, the square of the mass for the lowest state, $`m_0^2=\mathrm{\Gamma }[\varphi _0]`$, is the absolute minimum of $`\mathrm{\Gamma }[\varphi ]`$. To be properly defined as a Sturm-Liouville problem, it is necessary to impose boundary conditions: $`\tau (x)\varphi (x)\varphi ^{}(x)0,`$ for $`x1`$ and for $`x\mathrm{}`$, in accord with the boundary conditions stated earlier for our numerical solutions. Explicit limiting behaviors for all 12 cases are listed in Table 5. Given a trial wave function, $`\varphi (x)`$, Eq. B.1 provides a variational upper bound for $`m_0^2`$. As we have shown in Ref. , accurate variational estimates for ground-state masses can be obtained with minimum efforts. Here, we shall only attempt to obtain simple estimates by choosing trial functions so that integrals in the equation above can be evaluated analytically. The simplest possible trial wave function for each case can be chosen as a product of $`\tau (x)`$ and $`x^1`$, as indicated in Table 5. Indeed, our variational approach has served us well by providing a useful consistency check for our numerical efforts along the way. These are also summarized in Table 5. ### B.2 WKB As explained in Ref., we begin a WKB analysis by first bringing our differential equations from the Sturm-Liouville form into a radial-Schroedinger form , $$\left(\frac{d^2}{dx^2}+V(x;m^2)\right)\psi (x)=E\psi (x),$$ (B.2) by scaling $`\psi (x)\sqrt{\tau (x)}\varphi (x)`$. Eigenvalues $`m_n^2`$ are found by solving for zero-energy bound states from below, $`E0^{}`$. The potential $`V(x;m^2)`$ is given by $$V(x;m^2)=\frac{m^2\sigma (x)+w(x)}{\tau (x)}+{\scriptscriptstyle \frac{1}{2}}(\frac{\tau }{\tau }){\scriptscriptstyle \frac{1}{4}}(\frac{\tau ^{}}{\tau })^2.$$ (B.3) To obtain the desired WKB estimate, we simply need to evaluate in the large $`m^2`$ limit the following integral: $`(n+1/2)\pi =_1^{\mathrm{}}𝑑x\sqrt{\stackrel{~}{V}(x;m^2)},`$ i.e., we seek a WKB condition in the form, $`(n+\frac{1}{2})\pi =s_0m+s_1+0(\frac{1}{m}),`$ where coefficients $`s_0`$ and $`s_1`$ can be explicitly evaluated. As explained in Ref. , the “effective potential”, $`\stackrel{~}{V}(x,m^2)`$, is $`\stackrel{~}{V}(x;m^2)=V(x;m^2)+1/(4(x1)^2).`$ To isolate the $`m^2`$-dependence, let us write $`\stackrel{~}{V}(x,m^2)m^2V_0(x)+\stackrel{~}{V}_1(x)`$, where $`V_0(x)=\sigma (x)/\tau (x),`$ i.e., $`V_0^{(4)}(x)=1/(4(x^31))`$, and $`V_0^{(3)}(x)=1/(4x(x^21))`$. For instance, one immediately finds that $`s_0=_1^{\mathrm{}}𝑑x\sqrt{V_0(x)}`$. For $`QCD_4`$ and $`QCD_3`$, they are $`s_0^{(4)}=\frac{\sqrt{\pi }}{6}\mathrm{\Gamma }(\frac{1}{6})/\mathrm{\Gamma }(\frac{2}{3})`$, and $`s_0^{(3)}=\frac{\sqrt{\pi }}{4}\mathrm{\Gamma }(\frac{1}{4})/\mathrm{\Gamma }(\frac{3}{4})`$. The remaining piece, $`\stackrel{~}{V}_1(x)`$, is listed in Table 6. To find coefficient $`s_1`$, we need to know the behavior of the ratio of $`V_0(x)`$ to $`\stackrel{~}{V}_1(x)`$ near $`x1`$ and $`x\mathrm{}`$. The dominant behavior of $`\stackrel{~}{V}_1(x)`$, in these limits, can be characterized by two indices, $`\stackrel{~}{V}_1(x){\displaystyle \frac{\delta _l^2}{4(x1)^2}}+0({\displaystyle \frac{1}{(x1)}}),`$ $`for`$ $`x1,`$ (B.4) $`\stackrel{~}{V}_1(x){\displaystyle \frac{\delta _r^2}{4x^2}}+0({\displaystyle \frac{1}{x^3}}),`$ $`for`$ $`x\mathrm{}.`$ (B.5) One finds that $`s_1=(\frac{\pi }{2})(\delta _l+\delta _r),`$ and $`m^2=(\frac{\pi }{s_0})^2\{n^2+\delta n+\gamma \}+0(n^1)`$ where $$\delta =1+\delta _l+\delta _r,$$ (B.6) as exhibited in Table 6. The values for $`\gamma `$ are found by a fit to the numerical values of $`m_n^2`$.
warning/0003/astro-ph0003411.html	ar5iv	text	# Electron-Ion Recombination Rate Coefficients and Photoionization Cross Sections for Astrophysically Abundant Elements IV. Relativistic calculations for C IV and C V for UV and X-ray modeling ## 1 INTRODUCTION Electron-ion recombination with H- and He-like ions is of particular interest in X-ray astronomy (Proc. X-ray Symposium, 2000). X-ray emission in the K$`\alpha `$ complex of He-like ions, such as C V, from the $`n=21`$ transitions yields perhaps the most useful spectral diagnostics for temperature, density, ionization balance, and abundances in the plasma source (Gabriel 1972, Mewe and Schrijver 1978, Pradhan and Shull 1981, Pradhan 1985). Li-like C IV is of considerable importance in UV emission spectra from active galactic nuclei and quasars (e.g. Laor et al. 1994), as well as absorption in AGN (Crenshaw and Kraemer 1999). In addition, the C IV and other Li-like ionization states are valuable tracers of the plasma in the ‘hot interstellar medium’ (Spitzer 1990, Spitzer and Fitzpatick 1993, Martin and Bowyer 1990, Bregman and Harrington 1986). The primary sets of atomic data needed for accurate calculations of ionization fractions are for photoionization and recombination. Theoretical models of spectral formation also require excitation cross sections and transition probabilities. A considerable amount of atomic data is being computed for these atomic processes under the Iron Project (IP; Hummer et al. 1993) for electron impact excitation and radiative transition probabilities for astrophysically abundant elements using the Breit-Pauli R-matrix (BPRM) method that includes relativistic fine structure in intermediate coupling (Berrington et al. 1995). The present work is an extension of the IP work to photoionization and recombination. The ionization balance equations usually correspond to photoionization equilibrium $$_{\nu _0}^{\mathrm{}}\frac{4\pi J_\nu }{h\nu }N(X^z)\sigma _{PI}(\nu ,X^z)𝑑\nu =\underset{j}{}N_eN(X^{z+1})\alpha _R(X_j^z;T),$$ (1) and collisional equilibrium $$C_I(T,X^z)N_eN(X^z)=\underset{j}{}N_eN(X^{z+1})\alpha _R(X_j^z;T),$$ (2) where the $`\sigma _{PI}`$ are the photoionization cross sections, and the $`\alpha _R(X_j^z;T)`$ are the total electron-ion recombination rate coefficients of the recombined ion of charge $`z`$, $`X_j^z`$, to state j at electron temperature T. The $`C_I`$ are the rate coefficients for electron impact ionization that can be reliably obtained from experimental measurements (Bell et al. 1983). On the other hand, the (e + ion) recombination cross sections and rates are difficult to compute or measure. However, several experimental measurements of electron-ion recombination cross sections using ion storage rings have been carried out in recent years (e.g. Wolf et al. 1991, Kilgus et al. 1990,1993; Mannervik et al. 1997). The experimental cross sections exhibit detailed resonance structures observed at very high resolution in beam energy, and measure absolute cross sections. Therefore they provide ideal tests for theoretical methods, as well as the physical effects included in the calculations. Many of these experimental measurements have been for recombination with H- and He-like C and O. Among the recent theoretical developments is a self-consistent method for calculations for photoionization and (e + ion) recombination, as described in previous papers in this series. An identical eigenfunction expansion for the ion is employed in coupled channel calculations for both processes, thus ensuring consistently accurate cross sections and rates in an ab initio manner. The theoretical treatment of (e + ion) recombination subsumes both the non-resonant recombination (i.e. radiative recombination, RR), and the resonant recombination (i.e. di-electronic recombination, DR) processes in a unified scheme. In addition to the total, unified recombination rates, level-specific recombination rates and photoionization cross sections are obtained for a large number of atomic levels. The calculations are carried out in the close coupling approximation using the R-matrix method. Although the calculations are computationally intensive, they yield nearly all photoionization and recombination parameters needed for astrophysical photoionization models with higher precision than hitherto possible. Previous calculations of unified (e + ion) recombination cross sections and rates, reported in the present series on photoionization and recombination, were carried out in LS coupling (Nahar and Pradhan 1997, paper I; Nahar 1999). There were two reasons. First, the calculations are extremely complex and involve both radiative photoionization and collisional electron-ion scattering calculations; the full intermediate coupling relativistic calculations are many times more computationally intensive than the LS coupling ones. Second, the effect of relativistic fine structure was expected to be small for these light elements. For the highly charged H- and He-like ions, however, subsequent calculations showed that results including the relativistic effects are signficantly more accurate not only in terms of more detailed resonance structure, but also to enable a full resolution of resonances necessary to include radiative damping (Pradhan and Zhang 1997, Zhang et al. 1999, and references therein). The relativistic Breit-Pauli R-matrix (BPRM) method is now extended to calculate the total and level-specific recombination rate coefficients in the self consistent unified manner. In this paper we describe the first of a series of full-scale BPRM calculations of photoionization and photo-recombination, as inverse processes, and DR, to obtain total, unified (e + ion) recombination rates of He- and Li-like Carbon, C IV and C V. ## 2 THEORY The extension of the close coupling method to electron-ion recombination is described in earlier works (Nahar & Pradhan 1994, 1995), together with the details of the unified treatment. Here we present a brief description of the theory relevant to the calculations of electron recombination cross sections with H-like and He-like ions. The calculations are carried out in the close coupling (CC) approximation employing the $`R`$-matrix method in intermediate coupling with the BP Hamiltonian. The target ion is represented by an $`N`$-electron system, and the total wavefunction expansion, $`\mathrm{\Psi }(E)`$, of the ($`N`$+1) electron-ion system of symmetry $`SL\pi `$ or $`J\pi `$ may be represented in terms of the target eigenfunctions as: $$\mathrm{\Psi }(E)=A\underset{i}{}\chi _i\theta _i+\underset{j}{}c_j\mathrm{\Phi }_j,$$ (3) where $`\chi _i`$ is the target wavefunction in a specific state $`SL\pi `$ or level $`J_i\pi _i`$ and $`\theta _i`$ is the wavefunction for the ($`N`$+1)-th electron in a channel labeled as $`S_iL_i(J_i)\pi _ik_i^2\mathrm{}_i(J\pi ))`$; $`k_i^2`$ being its incident kinetic energy. $`\mathrm{\Phi }_j`$’s are the correlation functions of the ($`N`$+1)-electron system that account for short range correlation and the orthogonality between the continuum and the bound orbitals. In the relativistic BPRM calculations the set of $`SL\pi `$ are recoupled to obtain (e + ion) levels with total $`J\pi `$, followed by diagonalisation of the (N+1)-electron Hamiltonian, $$H_{N+1}^{BP}\mathrm{\Psi }=\mathrm{E}\mathrm{\Psi }.$$ (4) The BP Hamiltonian is $$H_{N+1}^{\mathrm{BP}}=H_{N+1}+H_{N+1}^{\mathrm{mass}}+H_{N+1}^{\mathrm{Dar}}+H_{N+1}^{\mathrm{so}},$$ (5) where $`H_{N+1}`$ is the nonrelativistic Hamiltonian, $$H_{N+1}=\underset{i=1}{\overset{N+1}{}}\left\{_i^2\frac{2Z}{r_i}+\underset{j>i}{\overset{N+1}{}}\frac{2}{r_{ij}}\right\},$$ (6) and the additional terms are the one-body terms, the mass correction term, the Darwin term and the spin-orbit term respectively. Spin-orbit interaction, $`H_{N+1}^{so}`$, splits the LS terms into fine-structure levels labeled by $`J\pi `$, where $`J`$ is the total angular momentum. The positive and negative energy states (Eq. 4) define continuum or bound (e + ion) states. $$\begin{array}{c}E=k^2>0continuum(scattering)channels\hfill \\ E=\frac{z^2}{\nu ^2}<0boundstates,\hfill \end{array}$$ (7) where $`\nu `$ is the effective quantum number relative to the core level. If $`E<`$ 0 then all continuum channels are ‘closed’ and the solutions represent bound states. The photoionization cross section can be obtained as $$\sigma _{PI}=\frac{1}{g}\frac{4\pi ^2}{3c}\omega 𝐒,$$ (8) where $`g`$ is the statistical weight factor of the bound state, $`𝐒`$ is the dipole line strength, $$𝐒=\|<\mathrm{\Psi }_B\|\|𝐃\|\|\mathrm{\Psi }_F>\|^2$$ (9) and $`𝐃`$ is the dipole operator (e.g. Seaton 1987). For highly charged ions (such as the H- and the He-like) radiative transition probabilities in the core ion are very large and may be of the same order of magnitude as autoionization probabilities. Autoionizing resonances may then undergo significant radiative decay and the photoionization process may be written as $$h\nu +X^+(X^+)^{}\{\begin{array}{c}(i)e+X^{++}\\ (ii)h\nu ^{}+X^+\end{array}$$ (10) Branch (ii) represents radiation damping of autoionizing resonances. In the present work this radiative damping effect is included for all near-threshold resonances, up to $`\nu 10`$, using a resonance fitting procedure (Sakimoto et al. 1990, Pradhan and Zhang 1997, Zhang et al. 1999). Recombination of an incoming electron to the target ion may occur through non-resonant, background continuum, usually referred to as radiative recombination (RR), $$e+X^{++}h\nu +X^+,$$ (11) which is the inverse process of direct photoionization, or through the two-step recombination process via autoionizing resonances, i.e. dielectronic recombination (DR): $$e+X^{++}(X^+)^{}\{\begin{array}{c}(i)e+X^{++}\\ (ii)h\nu +X^+\end{array},$$ (12) where the incident electron is in a quasi-bound doubly-excited state which leads either to (i) autoionization, a radiation-less transition to a lower state of the ion and the free electron, or to (ii) radiative stabilization to recombining ion states predominantly via decay of the ion core (usually to the ground state) and the bound electron. In the unified treatment the photoionization cross sections, $`\sigma _{\mathrm{PI}}`$, of a large number of low-$`n`$ bound states – all possible states with $`nn_{\mathrm{max}}10`$ – are obtained in the close coupling (CC) approximation as in the Opacity Project (Seaton 1987). Coupled channel calculations for $`\sigma _{\mathrm{PI}}`$ include both the background and the resonance structures (due to the doubly excited autoionizing states) in the cross sections. The recombination cross section, $`\sigma _{\mathrm{RC}}`$, is related to $`\sigma _{\mathrm{PI}}`$, through detailed balance (Milne relation) as $$\sigma _{\mathrm{RC}}(ϵ)=\frac{\alpha ^2}{4}\frac{g_i}{g_j}\frac{(ϵ+I)^2}{ϵ}\sigma _{\mathrm{PI}}$$ (13) in Rydberg units; $`\alpha `$ is the fine structure constant, $`ϵ`$ is the photoelectron energy, and $`I`$ is the ionization potential. In the present work, it is assumed that the recombining ion is in the ground state, and recombination can take place into the ground or any of the excited recombined (e+ion) states. Recombination rate coefficients of individual levels are obtained by averaging the recombination cross sections over the Maxwellian electron distribution, $`f(v)`$, at a given temperature as $$\alpha _{RC}(T)=_0^{\mathrm{}}vf(v)\sigma _{RC}𝑑v.$$ (14) The contributions of these bound states to the total $`\sigma _{\mathrm{RC}}`$ are obtained by summing over the contributions from individual cross sections. $`\sigma _{\mathrm{RC}}`$ thus obtained from $`\sigma _{\mathrm{PI}}`$, including the radiatively damped autoionizing resonances (Eq. (10)), corresponds to the total (DR+RR) unified recombination cross section. The recombination cross section, $`\sigma _{RC}`$ in Megabarns (Mb), is related to the collision strength, $`\mathrm{\Omega }_{\mathrm{RC}}`$, as $$\sigma _{RC}(ij)(Mb)=\pi \mathrm{\Omega }_{RC}(i,j)/(g_ik_i^2)(a_o^2/1.\times 10^{18}),$$ (15) where $`k_i^2`$ is the incident electron energy in Rydberg. As $`\sigma _{RC}`$ diverges at zero-photoelectron energy, the total collision strength, $`\mathrm{\Omega }`$, is used in the recombination rate calculations. Recombination into the high-$`n`$ states must also be included, i.e. $`n_{\mathrm{max}}<n\mathrm{}`$, (Fig. 1 of Nahar & Pradhan 1994). To each excited threshold $`S_iL_i(J_i)\pi _i`$ of the $`N`$-electron target ion, there corresponds an infinite series of ($`N`$+1)-electron states, $`S_iL_i(J_i)\pi _i\nu \mathrm{}`$, to which recombination can occur, where $`\nu `$ is the effective quantum number. For these states DR dominates the recombination process and the background recombination is negligibly small. The contributions from these states are added by calculating the collision strengths, $`\mathrm{\Omega }_{\mathrm{DR}}`$, employing the precise theory of radiation damping by Bell and Seaton (1985, Nahar & Pradhan 1994). Several aspects related to the application of the theory to the calculation of DR collision strengths are described in the references cited. General details of the theory and close coupling BPRM calculations are described in paper I and Zhang et al. (1999, and references therein). ## 3 COMPUTATIONS The electron-ion recombination calculations entail CC calculations for photoionization and electron-ion scattering. Identical eigenfunction expansion for the target (core) ion is employed for both processes; thus enabling inherently self-consistent photoionization/recombination results in an ab initio manner for a given ion. The total recombination cross sections, $`\sigma _{\mathrm{RC}}`$, are obtained from the photoionization cross sections, $`\sigma _{\mathrm{PI}}`$, and DR collision strengths, $`\mathrm{\Omega }_{\mathrm{DR}}`$, are calculated as descriped in paper I and Zhang et al. (1999). However, the computations for the cross sections are repeated with a much finer energy mesh in order to delineate the detailed resonance structures as observed in the experiments. Computations of photoionization cross sections, $`\sigma _{\mathrm{PI}}`$, in the relativistic BPRM intermediate coupling approximations are carried out using the package of codes from the Iron Project (Berrington et al. 1995; Hummer et al. 1993). However, radiation damping of resonances up to $`n=10`$ are included through use of the extended codes of STGF and STGBF (Pradhan & Zhang 1997). The $`R`$-matrix calculations are carried out for each total angular momentum symmetry $`J\pi `$, corresponding to a set of fine structure target levels $`J_t`$. In the energy region from threshold up to about $`\nu =\nu _{\mathrm{max}}`$ = 10 ($`\nu `$ is the effective quantum number of the outer orbital of the recombined ion bound state), detailed photorecombination cross sections are calculated as is Eq. (12). The electrons in this energy range recombine to a large number of final (e + ion) states and recombination cross sections are computed for all coupled symmetries and summed to obtain the total $`\sigma _{\mathrm{RC}}`$. The number of these final recombined states in the BPRM case is larger, owing to more channels involving fine structure, than the $`LS`$ coupling case. In the higher energy region, $`\nu _{\mathrm{max}}<\nu <\mathrm{}`$ below each threshold target level, where the resonances are narrow and dense and the background is negligible, we compute the detailed and the resonance averaged DR cross sections. The DR collision strengths in BPRM are obtained using extensions of the $`R`$-matrix asymptotic region codes (paper I; Zhang and Pradhan 1997) respectively. It is necessary to use extremely fine energy meshes in order to delineate the resonance structures belonging to each $`n`$-complex. The level specific recombination rate coefficients are obtained using a new computer program, BPRRC (Nahar and Pradhan 2000). The program extends the photoionization cross sections at the high energy region, beyond the highest target threshold in the CC wavefunction expansion of the ion, by a tail from Kramers fit of $`\sigma _{PI}(E)=\sigma _{PI}^o(E_o^3/E^3)`$. The level specific rates are obtained for energies going up to infinity. These rates include both non-resonant and resonant contributions up to energies $`z^2/\nu _{\mathrm{max}}^2`$; Contributions from all autoionizing resonances up to $`\nu \nu _{\mathrm{max}}10`$ are included. The program BPRRC sums up the level specific rates, which is added to the contributions from the high-n DR, to obtain the total recombination rates. As an additional check on the numerical calculations, the total recombination rate coefficients, $`\alpha _R`$, are also calculated from the total recombination collision strength, obtained from all the photoionization cross sections, and the DR collision strengths. The agreement between the two numerical approaches is within a few percent. The background contribution from the high-n states, ($`10<n\mathrm{}`$), to the total recombination is also included as the ”top-up” part (Nahar 1996). This contribution is important at low temperatures, but is negligible at high temperatures. The rapid rise in $`\alpha _R`$ toward very low temperatures is due to low energy recombination into the infinite number of these high-n states, at electron energies not usually high enough for resonant excitations. Below we describe the calculations individually for the ions under consideration. ### 3.1 e + C V $``$ C IV The fine structure levels of the target ion, C V, included in the wavefunction expansion for C IV are given in Table 1. The 13 fine structure levels of C V up to $`3p`$ correspond to configurations $`1s^2`$, $`1s2s`$, $`1s2p`$, $`1s3s`$, $`1s3p`$ (correlation configurations also involve the n = 4 orbitals). Although calculated energies are within a few percent of the observed ones, the latter are used in the computations for more accurate positions of the resonances. The bound channel wavefunction, second term in $`\mathrm{\Psi }`$ in Eq. (3), contains configuration $`3d^2`$. Levels of angular momentum symmeties $`1/2J9/2`$ are considered. With largest partial wave of the outer electron is $`l=9`$, these correspond to $`0L7`$ in doublet and quartet spin symmetries. The R-matrix basis set is represented by 30 continuum functions. It is necessary to represent the wavefunction expansion in the inner region of the R-matrix boundary with a relatively large number of terms in order to avoid some numerical problems that result in slight oscillation in computed cross sections. ### 3.2 e + C VI $``$ C V The wavefunction expansion of C V is represented by 9 fine structure levels (Table 1) of hydrogenic C VI from $`1s`$ to $`3d`$. Correlation orbitals $`4s`$, $`4p`$, $`4d`$, and $`4f`$ are also included. The $`SL\pi `$ symmetries consist of $`0L7`$ of singlet and triplet spin symmetries, for even and odd parities. All levels of C V with total angular momentum symmetry $`0J6`$ are included. The R-matrix basis set consist of 40 terms to reduce numerical instablities that might otherwise result in slight oscillations in the cross sections. ## 4 RESULTS AND DISCUSSION Results for photoionization and recombination are presented below, followed by a discussion of the physical features and effects. ### 4.1 Photoionization The ground state cross sections are needed for various astrophysical models, such as in determination of ionization fraction in photoionization equilibrium of plasma. Figs. 1 and 2 present the ground state photoionization cross section for C IV ($`1s^22s(^2S_{1/2})`$) and C V $`(1s^2(^1S_0))`$. Plots (a,b) in each figure show the total cross section (a), and the partial cross section (b) into the ground level of the residual ion. The total cross sections Fig. 1(a) and Fig. 2(a) show the K-shell ionization jump at the n = 2 target levels, i.e. inner-shell photoionization $$h\nu +CIV(1s^22s,2p)e+1s2s(,2p),$$ and $$h\nu +CV(1s^2,1s2p)e+2s(,2p).$$ In X-ray photoionization models these inner-shell edges play an important role in the overall ionization efficiency. For both C IV and C V, the first excited target n = 2 threshold(s) lie at a high energy and the cross sections show a monotonic decrease over a relatively large energy range. (Slight oscillations are seen in the ground level of C V due to some numerical instability; as mentioned earlier, such oscillations are reduced using larger R-matrix basis set.) The resonances at high energies belong to Rydberg series of $`n=2,3`$ levels. ## 5 Recombination cross sections and rate coefficients Figs. 3(a,b) present the total recombination cross sections, $`\sigma _{RC}`$, for C IV and C V. In contrast to the earlier presentation for the small energy range (Zhang et al. 1999) to compare with experiments, the figures display the cross sections from threshold up to the energy of the highest target threshold, $`3d`$, considered in the present work. The resonances belong to different n-complexes; these are in close agreeement with experimental data (Zhang et al. 1999). Fig. 4 presents total unified recombination rate coefficients for $`e+CVCIV`$. The solid curve is the present $`\alpha _R`$ in relativistic BP approximation, and short-long dashed is earlier unified rates in LS coupling and where radiation damping effect was not included (Nahar & Pradhan 1997). In the high temperature region, the earlier LS rates significantly overestimat the recombination rate. We compare the present BPRM rates with several other available sets of data, e.g. ‘experimentally derived DR’ rates (dot-long dashed curve, Savin 1999; which in fact include both the RR+DR contribution - see below), and previous theoretical DR rate coefficients in LS coupling (dot-dashed curve, Badnell et al. 1990). Zhang et al (1999) compared in detail the BPRM cross sections with experimental data from ion storage rings for $`e+CVCIV`$, with close agreement in the entire range of measurements for both the background (non-resonant) cross sections and resonances. The reported experimental data is primarily in the region of low-energy resonances that dominate recombination (mainly DR) with H- and He-like ions. The recombination rate coefficients, $`\alpha _R`$, obtained using the cross sections calculated by Zhang et al. (dotted curve) agree closely with those of Savin (1999) (dot-long dash curve) who used the experimental cross sections to obtain ’experimentally derived DR rates’. However, these rates do not include contributions from much of the low energy non-resonant RR and very high energy regions. The total unified $`\alpha _R(T)`$ (solid curve) which include all possible contributions is, therefore, somewhat higher than that obtained from limited energy range. The LS coupling DR rate by Badnell et al. (1990) (dot-dash curve) is lower than the others. The dashed and the long-dashed curves in the figure are RR rates by Aldrovandi & Pequignot (1973), and Verner and Ferland (1996); the latter agrees with the present rates at lower temperatures. In Figs. 5 and 6 we show the level-specific recombination rate coefficients into the lowest, and the excited, bound levels of C IV, for the $`1s^2ns,np`$ Rydberg series up to n = 10. These are the first such calculations; level-specific data have been obtained for all $`\mathrm{}9`$ and associated $`J\pi `$ symmetries. The behavior of the level-specific rates mimics that of the total (this is only true for the simple systems under consideration; in general the level specifc rates show significantly different structure for complex ions, as seen in our previous works). The only distinguishing feature is the DR bump. Since the numerical computations are enormously involved, particularly related to the resolution of resonances, the absence of any unidentified features in the level-specific rates is re-assuring. Total $`\alpha _R(T)`$ for $`e+CVICV`$ are given in Table 2, and are plotted in Fig. 7. The total unified recombination rate coefficients in the present BPRM calculations are plotted in the solid curve. The short-long dashed curve represents earlier rates obtained in LS coupling and with no inclusion of radiation damping of low-n autoionizing resonanes (Nahar & Pradhan 1997). The earlier LS rates overestimate the recombination rates at high temperatures. The dotted curve shows the rate coefficient computed using the Zhang et al. cross sections in a limited energy range with resonances, i.e. mainly DR. The DR rate by Shull & Steenberg (1982; dot-dash curve) agrees closely with the dotted curve. The dashed and the long-dashed curves are RR rates by Aldrovandi & Pequignot (1973), and Verner and Ferland (1996); they agrees with the present rates at lower temperatures. Table 2 presents the unified total BPRM recombination rate coefficients of C IV and C V averaged over a Maxwellian distribution. #### 5.0.1 Level-specific recombination rate coefficients Fig. 8 shows the level-specific rate coefficients for the ground and the excited n = 2 levels that are of considerable importance in X-ray spectroscopy, as they are responsible for the formation of the w,x,y,z lines from the 4 transitions $`1s^2(^1S_0)1s2p(^1P_1^o),1s2p(^3P_2^o),1s2p(^3P_1^o),1s2s(^3S_1)`$. The present work is particularly relevant to the formation of these X-ray lines since recombination-cascades from excited levels play an important role in determining the intensity ratios in coronal equilibrium and non-equilibrium plasmas (Pradhan 1985). The rates in Fig. 8 differ considerably from those by Mewe & Schrijver (1978, hereafter MS) that have been widely employed in the calculation of X-ray spectra of He-like ions (e.g. Pradhan 1982). We compare with the direct (RR + DR) rates separately calculated by MS using approximate Z-scaled RR and DR rates for the individual n = 2 levels of He-like ions. Their RR rates were from Z-scaled recombinaton rate of He<sup>+</sup> given by Burgess & Seaton (1960); the LS coupling data were divided according to the statistical weights of the fine structure levels. Their DR rates were obtained using averaged autoionization probabilities, Z-scaled from iron (Z = 26) and calculated with hydrogenic wavefunctions, together with radiative decay probabilities of the resonant $`2s2p,2p^2,(2p3s,2p3p,2p3d)`$ levels, decaying to the final n = 2 levels of He-like ions. The present work on the other hand includes DR contributions from all resonances up to $`2pn\mathrm{};n10,\mathrm{}n1`$. Figs. 3 and 4 of Zhang et al. (1999) show the detailed photorecombination cross sections for these resonance complexes. But the present rate coefficients are much lower (Fig. 8). It is surprising that the MS rates are much higher than the present ones. If we consider, for example, the level-specific rate for the $`1s2s(^3S_1)`$ level, the MS value includes contributions from only $`2s2p,3s2p,3p2p`$ autoionizing levels. That the MS values overestimate the rates is also indicated by the fact that, at Log T = 6.4 (the DR-peak temperature, Fig. 7), the sum of their individual n = 2 level-specific rates is 1.5 $`\times 10^{12}`$ cm<sup>-3</sup> sec<sup>-1</sup>, compared to our unified total $`\alpha _R`$ = 2.27 $`\times 10^{12}`$ cm<sup>-3</sup> sec<sup>-1</sup> (Table 2; C V). That would imply that the MS rates for the n = 2 levels alone account for 2/3 of the total recombination (RR + DR) rate for C V; which is unlikely. Fig. 9 presents level specific recombination rate coefficients of $`1sns(^3S)`$ Rydberg series of C V levels up to $`n`$ = 10. The features are similar to those of C IV. Although resolution of resonances in each cross section is very cumbersome, the sum of the level-specific rate coefficients, together with the DR contribution, agrees within a few percent of the total recombination rate coefficient. Recombination-cascade matrices may now be constructed for C IV and C V, and effective recombination rates into specific levels obtained accurately, using the direct recombination rates into levels with $`n10,\mathrm{}n1`$ (Pradhan 1985). The present data is more than sufficient for extrapolation to high-n,$`\mathrm{}`$ necessary to account for cascade contributions. Also needed are the radiative transition probabilities for all fine structure levels of C IV and C V, up to the n = 10 levels. They have also been calculated using the BPRM method, and will be available shortly (Nahar, in preparation). These data will be similar to that for Fe XXIV and Fe XXV calculated earlier under the Iron Project (Nahar and Pradhan 1999). We discuss below some of the important atomic effects relavant to the present calculations in particular, and electron-ion recombination in general. ### 5.1 Resolution and radiation damping of resonances It is important to resolve near-threshold resonances ($`\nu 10`$) at an adequately fine energy mesh in order to (a) compute accurately their contribution to the rate coefficient or the averaged cross section, and (b) to determine the radiative and autoionization rates through the fitting procedure referred to eariler (Sakimoto et al. 1990, Pradhan and Zhang 1997, Zhang et al. 1999). Resonances that are narrower than the energy intervals chosen have very low autoionization rates and are mostly damped out; their contribution to (a) should be small. ### 5.2 Intereference between resonant (DR) and non-resonant (RR) recombination In general there is quantum mechanical interference between the resonant and the non-resonant components of the wavefunction expansion. Close coupling photoioniation calculations for strongly coupled near-neutral atomic systems cross sections considerable overlap between members of several Rydberg series of resonances that coverge on to the excited, coupled target levels. Unified (e + ion) rates have been calculated for over 40 atoms and ions of astrophysical interest. The resonance structures in a number of these are extremely complicated and show considerable interference. The large number of photoionization calculations under the Opacity Project, for the ground state and the excited states (typically, a few hundred excited states for each atom or ion), show overlapping resonances to often dominate the cross sections. These are archived in the Opacity Project database TOPBASE and may be accessed on-line via the Website: http://heasarc.gsfc.nasa.gov, or via the link from www.astronomy.ohio-state.edu/$``$pradhan). Furthermore, the excited state cross sections of excited metastable states may exhibit even more extensive resonance structures than the ground state (e.g. in photoionization of O III; Luo et al. 1989). Strictly speaking, one needs to consider the partial photoionization cross sections for each of the three cases (dipole photoionization): $`\mathrm{\Delta }S=0,\mathrm{\Delta }l=0,\pm 1`$ in LS coupling, or $`\mathrm{\Delta }J=0,\pm 1`$ in BPRM calculations. Although these are not often tabulated or displayed, an examination thereof reveals that for near-neutrals each partial cross section shows overlapping resonances. As the ion charge increases, the resonances separate out, the intereference decreases, and isolated resonance approximations may be used (Pindzola et al. 1992, Zhang 1998). Therefore, in general, we expect interference between the non-resonant RR and resonant DR recombination processes, particularly for many-electron systems where a separation between the two processes is unphysical and imprecise. Of course experimentally such a division is artificial and is not possible. For example, the recent experimental measurements on electron recombination with Fe XVIII to Fe XVII (Savin et al. 1999) clearly show the near threshold cross section to be dominated by the non-resonant recombination (RR) towards E $``$ 0, with superimposed resonance structures (the DR contribution) at higher energies. In a recent work Savin (1999) cited the work of Pindzola et al. (1992) to state that the effect of interferences is small (Pindzola et al. did not carry out close coupling calculations such as in the Opacity Project). Although this is not true in general, for highly ionized few-electron ions one expects sufficient resonance separation so that an independent treatment of RR and DR may be accurate; such is the case for recombination with H-like and He-like ions. The present unified treatment accounts for interference effects in an ab initio manner, and (e + ion) rates have been calculated for over 40 atoms and ions of astrophysical interest (heretofore, in LS coupling). ### 5.3 Comparision with experimental data and uncertainties Although experimental results are available for relatively few ions in limited energy ranges, and mostly for simple atomic systems such as the H-like and He-like ions, they are very useful for the calibration of theoretical cross sections. There is very good agreement with experimentally measured cross sections for electron recombination to C IV, C V, and O VII, as discussed in detail by Zhang et al. (1999). However, the energy range of experimental measurements is much smaller; the theoretical calculations are from E = 0 to very high energies necessary to obtain rate coefficients up to T = $`10^9`$K. For recombination with H- and He-like ions most (but not all) of the rates depends on relatively few low-n complex of resonances in the low-energy region covered by experiments. In a recent work Savin (1999) used the experimental cross sections for recombination with C V to C IV, and O VIII to O VII, to obtain ‘experimentally derived DR’ rate coeffcients and compared those with several sets of theoretical data including those in papers I and II. (Strictly speaking these rates include ‘RR + DR’ contributions since the experimental cross sections always include both, and which the unified method aims to obtain). Whereas the previous LS coupling results for C IV (paper I) were 43% higher than Savin’s values, the results for O VII were in reasonable agreement, $``$20% higher, within estimated uncertainties. The agreement between the unified rates and the experimentally derived DR rates is within our 10-20% in the region (i.e. at temperatures) where DR contribution peaks, around Log(T) = 6.3 for C IV (Fig. 2). As the reported experimental data did not extend to low energies, where the non-resonant RR contribution dominates, the unified rates are higher towards lower temperatures from the DR peak, and deviate in a predictably straightforward manner from the DR-only results. The ‘experimentailly derived’ data by Savin (1999) was obtained with a limited energy range, while our present results include a much larger range. Therefore our results are higher. The dotted curve in Fig. 4, which we also obtained over a limited energy range, agrees well with the dot-long-dashed curve (Savin 1999). The differences are within our estimate of uncertainty in the present results, up to 20%. Given that the experimental cross sections are also likely to be uncertain to about this range, the agreement seems remarkably good. While the present unified cross sections can be compared directly with experimental measurements, and the new rate coefficients are in good agreement with the experimentaly derived DR rates for recombination with simple ions such as the H- and He-like, the experimental data may represent a lower bound on the field-free theoretical recombination rates owing to (a) high-n and $`\mathrm{}`$ ionization reducing the ’DR’ peak, and (b) limited energy range in experiments. ### 5.4 General features of (e + ion) recombination rates The non-resonant ‘RR’ recombination peaks as E, T $``$ 0. This is due to the dominant contribution from an infinite number of high Rydberg states of the (e + ion) system into which the slow moving electron may recombine. At low-E and T, the total $`log_{10}(\alpha _R(T)`$) is shown as a straight line on the Log-Log scale due to the exponential Maxwellian damping factor exp(-E/kT). It is not entirely trivial to compute the low-E and T contributions (that we refer to as “high-n top-up”). We adapt the accurate numerical procedure developed by Storey & Hummer (1992) to calculate the n,$`\mathrm{}`$ hydrogenic photoionization cross section for $`11n\mathrm{}`$ (Nahar 1996). It is noted that the high-n top-up also represents the otherwise missing background contribution due to high-n resonant recombination (DR). Although this background contribution is small (negligibly so for the H- and He-like ions), it is included in the unified treatment. The resonant contribution (DR) peaks at higher E and T corresponding to the excitation energies and temperatures of the strong dipole transition(s) in the core ion. This is the broad peak in $`\alpha _R(T)`$. ### 5.5 Ionization fractions of Carbon Fig. 10 presents coronal ionization fractions of C using the new BPRM recombinaton rates for C IV and C V (solid lines). Also given are the results (dashed lines) from Arnaud and Rothenflug (1985), and previous results (dotted lines) using LS coupling rates from (Nahar and Pradhan 1997). Differences with both sets of data may be noted for C IV, C V, C VI, and C VII. The most significant change (enhancement) is for C VI, owing to the decrease in C V recombination rate, and the new ionization fractions appear to be in better agreement with Arnaud and Rothenflug (1985) than the Nahar and Pradhan (1997) results. Also discernible is the steeper fall-off in the C V ionization fraction on the high temperatures side. ## 6 CONCLUSION New relativistic calculations are presented for the total, unified (e + ion) rates coefficients for C IV and C V of interest in X-ray astronomy. As the photo-recombination cross sections in the dominant low-energy region have earlier been shown to be in very good agreeement with experiments (Zhang et al. 1999), it is expected that the present rates should be definitive, with an uncertainty that should not exceed 10–20%. The unified theoretical formulation and experimental measurements both suggest that the unphysical and imprecise division of the recombination process into ’radiative recombination (RR)’ and ’di-electronic recombination (DR)’ be replaced by ’non-resonant’ and ’resonant’ recombination, since these are naturally inseparable. Further calculations are in progress for Oxygen (O VI and O VII) and Iron (Fe XXIV and FeXXV). The available data includes: (A) Photoionization cross sections for bound fine structure levels of C IV and C V up to n = 10 – both total and partial (into the ground level of the residual ion). !!!! (B) Total, unified recombination rates for C IV and C V, and level-specific recombination rate coefficients for levels up to n = 10. All photoionization and recombination data are available electronically from the first author at: nahar@astronomy.ohio-state.edu. The total recombination rate coefficients are also available from the Ohio State Atomic Astrophysics website at: www.astronomy.ohio-state.edu/$``$pradhan. This work was supported partially by grants from NSF (AST-9870089) and NASA (NAG5-8423). The computational work was carried out on the Cray T94 at the Ohio Supercomputer Center.
warning/0003/hep-th0003245.html	ar5iv	text	# 1 Introduction ## 1 Introduction Field theories with N=2 extended supersymmetry are known to possess remarkable properties that sometimes allow one to obtain their exact (non-perturbative) low-energy solutions, pioneered by Seiberg and Witten . The natural way to produce such exact results is to relate the field theory problem to an integrable system (or Whitham dynamics in the Seiberg-Witten case) . The origin of an elliptic curve behind the exact solution also becomes apparent in this approach. Another (related) interpretation of the Seiberg-Witten result is possible from the viewpoint of anomalous breaking of N=2 superconformal symmetry and uniformization theory on Riemann surfaces . It is not very surprising since N=2 superconformal symmetry is well-known to be instrumental in deriving the classical structure of all (non-conformal) N=2 supersymmetric field theories, including N=2 supergravity, in four spacetime dimensions. It is natural to examine first the constraints imposed by N=2 superconformal invariance and then study compensation of unwanted N=2 superconformal symmetries. Or one can start from a classical field theory that is N=2 superconformally invariant, and then investigate its superconformal anomalies that can be developed in quantum field theory. Much work in the recent past was devoted to investigating the constraints imposed by N=4 superconformal symmetry on the correlation functions of the N=4 super-Yang-Mills theory in the context of AdS/CFT correspondence . More recently, similar constraints on the correlation functions were studied in the context of four-dimensional N=2 conformal supersymmetry . We recall that all N=2 supersymmetric field theories can be brought into the manifestly N=2 supersymmetric form, by using off-shell (unconstrained) N=2 superfields in Harmonic Superspace (HSS) . It makes the effective field theory methods to be very efficient in the N=2 case contrary to the N=4 case where only on-shell N=4 supersymmetry is possible. The next obvious step is to make manifest N=2 superconformal invariance of the N=2 supersymmetric quantum effective action, and then study its superconformal anomalies. Truly non-perturbative results are expected to be obtained along these lines. As regards quantum field theories with rigid N=2 supersymmetry, their building blocks are given by N=2 vector multiplets and hypermultiplets. At the level of the Low-Energy Effective Action (LEEA), on the (abelian) N=2 vector multiplet side we have to deal with the Seiberg-Witten-type action specified by a holomorphic potential and the associated special Kähler geometry. This topic is well-known, and we are going to use it as our basic pattern to follow. Hypermultiplets (e.g., if they are magnetically charged) may also develop the non-trivial LEEA that takes the form of a hyper-Kähler Non-Linear Sigma-Model (NLSM) by N=2 supersymmetry. After being formulated in HSS, the N=2 NLSM also possess an (analytic) potential. Our purpose in this paper is to impose and make manifest N=2 superconformal invariance in the N=2 NLSM, and then formulate an anomalous N=2 superconformal Ward identity on the hyper-Kähler potential. The paper is organized as follows. In sect. 2 we briefly review the N=2 supercurrents, the anomalous N=2 superconformal Ward identities in the N=2 gauge sector, and their relation to the solution of Seiberg and Witten . In sect. 3 we construct in HSS the most general N=2 superconformal NLSM that gives a general solution to the special hyper-Kähler geometry. The anomalous N=2 superconformal Ward identity on the hypermultiplet LEEA (N=2 NLSM) is found for the first time in sect. 4, by using the N=2 supergravity compensators in HSS. In sect. 5 we give our conclusions. In Appendix we briefly review a construction of N=2 supergravity in HSS. ## 2 N=2 supercurrent and Ward identities An N=2 supercurrent is the irreducible representation of N=2 supersymmetry in four spacetime dimensions, having superspin one. The independent field components of the N=2 supercurrent (of some N=2 matter system) include the energy-momentum tensor, the N=2 supersymmetry current, the central charge current, the axial current, the $`SU(2)`$ current of R-symmetry and some auxiliary field components of lower dimension. The relevant field components of the N=2 supercurrent were first identified in ref. by analyzing the free field theory of a massive (Fayet-Sohnius) hypermultiplet. The systematic way of derivation of the N=2 supercurrent superfield is provided by a construction of the irreducible N=2 superfields in the conventional (flat) N=2 superspace $`\{𝒵\}=(x^m,\theta _\alpha ^i,\overline{\theta }_j^\stackrel{_{{}_{}{}^{}}}{\alpha })`$, where $`m=0,1,2,3`$, $`\alpha =1,2`$ and $`i,j=1,2`$. All irreducible N=2 superprojectors were found in ref. , whereas all irreducible (scalar) N=2 superfields were explicitly derived in ref. . Amongst the irreducible N=2 superfields, comprising a general N=2 real scalar superfield, one finds almost all off-shell N=2 supermultiplets that usually appear in any discussion of N=2 supersymmetry (with a finite number of the auxiliary fields). In particular, an N=2 (restricted) chiral superfield $`\mathrm{\Phi }(x,\theta ,\overline{\theta })`$ is defined by the N=2 superspace (off-shell) constraints $$\overline{D}_\stackrel{_{{}_{}{}^{}}}{\alpha }^i\mathrm{\Phi }=0,D^4\mathrm{\Phi }=\mathrm{}\overline{\mathrm{\Phi }},$$ $`(2.1)`$ where $`(D_i^\alpha ,\overline{D}_\stackrel{_{{}_{}{}^{}}}{\alpha }^j)`$ are the usual (flat) N=2 superspace covariant derivatives. <sup>2</sup><sup>2</sup>2We use the notation $`D_{\alpha \beta }=D_{i\alpha }D_\beta ^i`$, $`D^{ij}=D^{i\alpha }D_\alpha ^j`$ and $`D^4=\frac{1}{12}D_{ij}D^{ij}`$, and similarly for the conjugated quantities. As a consequence of the constraints (2.1), the N=2 restricted chiral superfield $`\mathrm{\Phi }`$ possess a two-form $`F=F_{mn}dx^mdx^n`$ satisfying the ‘Bianchi identity’ $`dF=0`$. A solution to the ‘Bianchi identity’, $`F=dA`$, in terms of the one-form $`A`$ subject to the gauge transformations $`\delta A=d\lambda `$, allows one to represent the N=2 superfield strength $`W`$ of an abelian N=2 vector multiplet by a restricted chiral N=2 superfield too. The N=2 restricted chiral superfield $`\mathrm{\Phi }`$ is dual to an N=2 linear superfield $`L^{ij}`$. The latter is symmetric with respect to its $`SU(2)`$ indices and satisfies the off-shell constraints $$D_\alpha ^{(i}L^{jk)}=\overline{D}{}_{\stackrel{_{{}_{}{}^{}}}{\alpha }}{}^{(i}L_{}^{jk)}=0,\overline{(L^{ij})}=\epsilon _{ik}\epsilon _{jl}L^{kl}.$$ $`(2.2)`$ The duality relation in N=2 superspace is just given by $$L^{ij}=D^{ij}\mathrm{\Phi }.$$ $`(2.3)`$ Both superfields $`\mathrm{\Phi }`$ and $`L^{ij}`$ represent the irreducible N=2 multiplets of superspin zero and superisospin zero. Similarly, an irreducible N=2 scalar superfield $`R`$ of superspin zero and superisospin one is defined by the constraints $$D_{\alpha \beta }R=\overline{D}_{\stackrel{_{{}_{}{}^{}}}{\alpha }\stackrel{_{{}_{}{}^{}}}{\beta }}R=iD_\alpha ^j,\overline{D}_{j\stackrel{_{{}_{}{}^{}}}{\alpha }}R=0.$$ $`(2.4)`$ The irreducible N=2 scalar superfield $`R`$ is dual to an N=2 projective superfield $`T^{ijkl}`$, $$T^{ijkl}=D^{(ij}\overline{D}^{kl)}R,$$ $`(2.5)`$ where $`T^{ijkl}`$ is totally symmetric with respect to its $`SU(2)`$ indices and obeys the (off-shell) constraints $$D_\alpha ^{(i}T^{jklm)}=\overline{D}{}_{\stackrel{_{{}_{}{}^{}}}{\alpha }}{}^{(i}T_{}^{jklm)}=0,\overline{(T^{i_1i_2i_3i_4})}=\epsilon _{i_1j_1}\epsilon _{i_2j_2}\epsilon _{i_3j_3}\epsilon _{i_4j_4}T^{j_1j_2j_3j_4}.$$ $`(2.6)`$ The N=2 supercurrent $`J`$ is also in the list of the irreducible N=2 scalar superfields, being defined by the constraints $$D_{ij}J=\overline{D}_{ij}J=0.$$ $`(2.7)`$ The N=2 superspace constraints (2.7) imply that the energy-monentum tensor is symmetric, conserved and traceless, whereas all the vector currents are conserved. In other words, $`J`$ is a multiplet of N=2 superconformal currents. An N=2 superconformal anomaly amounts to breaking the N=2 supercurrent conservation relations (2.7) by an N=2 anomaly multiplet of lower superspin, i.e. of superspin zero. This is equivalent to activating an irreducible superspin-zero superfield in the N=2 scalar superfield $`J`$. As is clear from the above discussion, there are potentially two ways of assigning the N=2 anomaly multiplet with an N=2 irreducible superspin-zero superfield: either $`\mathrm{\Phi }`$ (or, equivalently, $`L^{ij}`$) or $`R`$ (or, equivalently, $`L^{ijkl}`$). The main difference between the two choices is the fact that $`L^{ij}`$ still contains a conserved vector current (associated with unbroken central charge transformations), whereas the vector current in $`L^{ijkl}`$ is not conserved. The first choice yields the N=2 superconformal anomaly relation in the standard form $$\frac{i}{4}D_{ij}J=L_{ij}.$$ $`(2.8)`$ For example, the N=2 supercurrent conservation law in the quantum N=2 SYM theory takes the form $$\frac{i}{4}D_{ij}J=\overline{D}_{ij}\overline{S},S=\frac{c}{2}\mathrm{tr}W^2,$$ $`(2.9)`$ where the (Lie algebra-valued) N=2 SYM superfield strength $`W`$ has been introduced. The constant $`c`$ is proportional to the one-loop renormalization group beta-function. Though $`\mathrm{tr}W^2`$ is merely a chiral (not a restricted chiral) N=2 superfield, eq. (2.9) can be easily brought into the form (2.8) by a local shift of the supercurrent, $`JJ4iS`$. The most general N=2 supersymmetric Ansatz for the LEEA of some number $`(r)`$ of abelian N=2 vector multiplets is governed by a holomorphic potential $`(W)`$ of the N=2 (restricted chiral) superfield strengths $`W_p`$ ,<sup>3</sup><sup>3</sup>3Our normalization differs from that used in ref. by a factor $`i/(16\pi )`$. $$I[W]=d^4xd^4\theta (W_p)+\mathrm{h}.\mathrm{c}.$$ $`(2.10)`$ The superfield $`W`$ has conformal weight $`+1`$, in accordance with its canonical dimension, whereas the N=2 chiral superspace measure in eq. (2.10) has conformal weight $`(2)`$. It is, therefore, clear that the N=2 superconformal Ward identity for the LEEA (2.10) is given by $$\underset{p=1}{\overset{r}{}}W_p\frac{}{W_p}2=0.$$ $`(2.11)`$ The N=2 superconformal solution to the holomorphic potential $`(W)`$ is thus given by a homogeneous (of degree two) function. The (rigid) N=2 superconformal invariance of the action (2.10) is the necessary pre-requisite for its coupling to N=2 conformal supergravity . Given a non-trivial renormalization flow (like in the N=2 SYM theory), the N=2 superconformal Ward identity (2.11) is going to be broken by the anomaly $`S`$, $$\underset{p=1}{\overset{r}{}}W_p\frac{}{W_p}2=4S.$$ $`(2.12)`$ This equation is just the anomalous N=2 superconformal Ward identity for (abelian) N=2 vector multiplets, which was found in ref. by ‘averaging’ the anomaly relation (2.9) with respect to the quantum effective action (2.10). A simple derivation of eq. (2.12) by using the N=2 supergravity compensators is discussed in sect. 4. Equation (2.12) can be applied to a derivation of the Seiberg-Witten solution provided that one knows the anomaly $`S`$ as the function of $`W`$, in the context of the N=2 SYM theory based on the gauge group $`SU(2)`$ spontaneously broken to $`U(1)`$. The original derivation made use of the electric-magnetic duality and renormalization flow. Since the anomaly $`S`$ is N=2 chiral, gauge-invariant and of dimension two, its vacuum expectation value has to be proportional to the order parameter $`u=\frac{1}{2}\mathrm{tr}W^2`$, with the coefficient being dictated by the one-loop beta-function $`\beta _1`$ — see eq. (2.9). A comparision with ref. yields $$c=2\pi i\beta _1.$$ $`(2.13)`$ To close eq. (2.12), as the equation on $``$, one needs a relation between $`a=W`$ and $`u`$. It was obtained in ref. by using the modular invariance of $`u=u(a)`$, in the form of a non-linear differential equation, $$(1u^2)u^{\prime \prime }+\frac{1}{4}au^{}{}_{}{}^{3}=0.$$ $`(2.14)`$ A connection to integrable systems arises after identifying the moduli space of the Coulomb branch in the Seiberg-Witten model with the moduli space of complex structures on an elliptic curve. The Seiberg-Witten solution then appears to be a classical solution to the equations of motion of a particular spin chain system , while the origin of the elliptic curve underlying the dynamics becomes apparent in this approach. <sup>4</sup><sup>4</sup>4The origin of the elliptic curve in the Seiberg-Witten exact solution is also explained by brane technology in the context of M-theory . Generalizations to the larger gauge groups and the presence of N=2 matter are straightforward, in principle. Our main goal, however, is to explore what can happen on the hypermultiplet side, at the level of the LEEA. Unlike the N=2 vector multiplets, the universal and most symmetric off-shell formulation of hypermultiplets is only possible in HSS, with the infinite number of the auxiliary fields. ## 3 Rigid N=2 superconformal symmetry in HSS In the HSS approach the standard N=2 superspace coordinates $`\{𝒵\}=(x^m,\theta _\alpha ^i,\overline{\theta }_j^\stackrel{_{{}_{}{}^{}}}{\alpha })`$ are extended by bosonic harmonics (or twistors) $`u^{\pm i}`$, $`i=1,2`$, belonging to the group $`SU(2)`$ and satisfying the unimodularity condition $$u^{+i}u_i^{}=1,\overline{u^{i+}}=u_i^{}.$$ $`(3.1)`$ The hidden analyticity structure of the N=2 superspace constraints defining both N=2 vector multiplets and FS hypermultiplets, as well as their solutions in terms of unconstrained N=2 superfields, can be made manifest in HSS . Instead of an explicit parametrization of the twistor sphere $`S^2=SU(2)/U(1)`$, the $`SU(2)`$-covariant HSS approach deals with the equivariant functions of harmonics, having the definite $`U(1)`$ charges defined by $`U(u_i^\pm )=\pm 1`$. The simple harmonic integration rules, $$𝑑u=1\mathrm{and}𝑑uu^{+(i_1}\mathrm{}u^{+i_m}u^{j_1}\mathrm{}u^{j_n)}=0\mathrm{otherwise},$$ $`(3.2)`$ are similar to the (Berezin) integration rules in superspace. In particular, any harmonic integral over a $`U(1)`$-charged quantity vanishes. The harmonic covariant derivatives, preserving the defining equations (3.1) in the original (central) basis, are given by $$^{++}=u^{+i}\frac{}{u^i},^{}=u^i\frac{}{u^{+i}},^0=u^{+i}\frac{}{u^{+i}}u^i\frac{}{u^i}.$$ $`(3.3)`$ They satisfy an $`su(2)`$ algebra and commute with the standard (flat) N=2 superspace covariant derivatives $`D_i^\alpha `$ and $`\overline{D}_\stackrel{_{{}_{}{}^{}}}{\alpha }^j`$. The operator $`^0`$ measures $`U(1)`$ charges. The key feature of HSS is the existence of the analytic subspace parametrized by $$(\zeta ^M;u)=\left\{\begin{array}{c}x_{\mathrm{analytic}}^{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}=x^{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}4i\theta ^{i\alpha }\overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }j}u_{(i}^+u_{j)}^{},\theta _\alpha ^+=\theta _\alpha ^iu_i^+,\overline{\theta }_\stackrel{_{{}_{}{}^{}}}{\alpha }^+=\overline{\theta }_\stackrel{_{{}_{}{}^{}}}{\alpha }^iu_i^+;u_i^\pm \end{array}\right\},$$ $`(3.4)`$ which is invariant under N=2 (rigid) supersymmetry : $$\delta x_{\mathrm{analytic}}^{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}=4i\left(\epsilon ^{i\alpha }\overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }+}+\theta ^{\alpha +}\overline{\epsilon }^{\stackrel{_{{}_{}{}^{}}}{\alpha }i}\right)u_i^{}4i\left(\epsilon ^\alpha \overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }+}+\theta ^{\alpha +}\overline{\epsilon }^\stackrel{_{{}_{}{}^{}}}{\alpha }\right),$$ $$\delta \theta _\alpha ^+=\epsilon _\alpha ^iu_i^+\epsilon _\alpha ^+,\delta \overline{\theta }_\stackrel{_{{}_{}{}^{}}}{\alpha }^+=\overline{\epsilon }_\stackrel{_{{}_{}{}^{}}}{\alpha }^iu_i^+\overline{\epsilon }_\stackrel{_{{}_{}{}^{}}}{\alpha }{}_{}{}^{+},\delta u_i^\pm =0.$$ $`(3.5)`$ The analytic dependence includes $`\theta _{\widehat{\alpha }}^+`$ but not $`\theta _{\widehat{\alpha }}^{}`$, where $`\widehat{\alpha }=(\alpha ,\stackrel{_{{}_{}{}^{}}}{\alpha })`$. The usual complex conjugation does not preserve analyticity. However, it does, after being combined with another (star) conjugation that only acts on the $`U(1)`$ indices as $`(u_i^+)^{}=u_i^{}`$ and $`(u_i^{})^{}=u_i^+`$. One has $`\stackrel{}{\overline{u^{\pm i}}}=u_i^\pm `$ and $`\stackrel{}{\overline{u_i^\pm }}=u^{\pm i}`$. Analytic superfields $`\varphi ^{(q)}(\zeta (𝒵,u),u)`$ of any positive (integral) $`U(1)`$ charge $`q`$ in HSS are defined by the off-shell constraints (cf. the definition of N=1 chiral superfields) $$D_\alpha ^+\varphi ^{(q)}=\overline{D}_\stackrel{_{{}_{}{}^{}}}{\alpha }^+\varphi ^{(q)}=0,\mathrm{where}D_\alpha ^+=D_\alpha ^iu_i^+\mathrm{and}\overline{D}_\stackrel{_{{}_{}{}^{}}}{\alpha }^+=\overline{D}_\stackrel{_{{}_{}{}^{}}}{\alpha }^iu_i^+.$$ $`(3.6)`$ The analytic measure reads $`d\zeta ^{(4)}dud^4x_{\mathrm{analytic}}^md^2\theta ^+d^2\overline{\theta }^+du`$, and it is of $`U(1)`$ charge $`(4)`$. The covariant derivatives in the analytic basis (3.4) receive certain connection terms. For example, the harmonic derivative $`^{++}`$ in the analytic subspace is replaced by $$D^{++}=^{++}4i\theta ^{\alpha +}\overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }+}\frac{}{x^{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}}.$$ $`(3.7)`$ This derivative preserves analyticity and permits integration by parts. Similarly, one easily finds the $`U(1)`$ charge operator in the analytic subspace reads $$D^0=u^{+i}\frac{}{u^{+i}}u^i\frac{}{u^i}+\theta ^{\alpha +}\frac{}{\theta ^{\alpha +}}+\overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }+}\frac{}{\overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }+}}.$$ $`(3.8)`$ In what follows we always use the analytic basis and the associated HSS covariant derivatives denoted by capital $`D`$, without making explicit references. The use of harmonics also gives us control over the (linearly realised) $`SU(2)_R`$ symmetry (or its absence), in the context of manifest N=2 supersymmetry (see ref. for more details). Since the translational and Lorentz symmetries, as well as N=2 supersymmetry, are manifestly realized in HSS, the latter provides us with the natural arena for a study of ‘truly’ N=2 superconformal symmetries on the top of N=2 non-conformal (rigid or Poincaré) supersymmetry. The superfield transformation rules with respect to dilatations (with the infinitesimal parameter $`\rho `$) are dictated by conformal weights $`w`$ of the superfields, together with the weights of the N=2 superspace coordinates, $$w[x]=1,w[\theta ]=w[\overline{\theta }]=\frac{1}{2},w[u]=0.$$ $`(3.9)`$ The non-trivial part of the N=2 superconformal transformations is given by $`SU(2)_{\mathrm{conf}.}`$ internal rotations with the parameters $`l^{ij}`$, special conformal transformations with the parameters $`k_{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}`$, and N=2 special supersymmetry with the parameters $`\eta _\alpha ^i`$ and $`\overline{\eta }_i^\stackrel{_{{}_{}{}^{}}}{\alpha }`$. The N=2 superconformal extension of the spacetime conformal transformations, $$\delta x^{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}=\rho x^{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}+k_{\beta \stackrel{_{{}_{}{}^{}}}{\beta }}x^{\alpha \stackrel{_{{}_{}{}^{}}}{\beta }}x^{\beta \stackrel{_{{}_{}{}^{}}}{\alpha }},$$ $`(3.10)`$ is dictated by the requirement of preserving the unimodularity and analyticity conditions in eqs. (3.1) and (3.6), respectively. As regards the non-trivial part of the N=2 superconformal transformation laws, one finds \| $`\delta x^{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}`$ \| $`=4i\lambda ^{ij}u_i^{}u_j^{}\theta ^{\alpha +}\overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }+}+k_{\beta \stackrel{_{{}_{}{}^{}}}{\beta }}x^{\alpha \stackrel{_{{}_{}{}^{}}}{\beta }}x^{\beta \stackrel{_{{}_{}{}^{}}}{\alpha }}+4i\left(x^{\alpha \stackrel{_{{}_{}{}^{}}}{\beta }}\overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }+}\overline{\eta }_\stackrel{_{{}_{}{}^{}}}{\beta }^{}x^{\stackrel{_{{}_{}{}^{}}}{\alpha }\beta }\theta ^{\alpha +}\eta _\beta ^{}\right),`$ \| \| --- \| --- \| \| $`\delta \theta ^{\alpha +}`$ \| $`=\lambda ^{ij}u_i^+u_j^{}\theta ^{\alpha +}+k_{\beta \stackrel{_{{}_{}{}^{}}}{\beta }}x^{\alpha \stackrel{_{{}_{}{}^{}}}{\beta }}\theta ^{\beta +}2i(\theta ^{\beta +}\theta _\beta ^+)\eta ^\alpha +x^{\alpha \stackrel{_{{}_{}{}^{}}}{\beta }}\overline{\eta }_\stackrel{_{{}_{}{}^{}}}{\beta }^+,`$ \| \| $`\delta \overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }+}`$ \| $`=\stackrel{}{\overline{(\delta \theta ^{\alpha +})}},`$ \| \| $`\delta u_i^+`$ \| $`=\left[\lambda ^{kj}u_k^+u_j^++4ik_{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}\theta ^{\alpha +}\overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }+}+4i\left(\theta ^{\alpha +}\eta _\alpha ^++\overline{\eta }_\stackrel{_{{}_{}{}^{}}}{\alpha }^+\overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }+}\right)\right]u_i^{},`$ \| \| $`\delta u_i^{}`$ \| $`=0.`$ \| $`(3.11)`$ Since the building blocks of invariant actions in HSS are given by the measure, analytic superfields and HSS covariant derivatives, only their transformation properties under the rigid ‘truly’ N=2 superconformal transformations are needed. It follows from eq. (3.11) that $$\mathrm{Ber}\frac{(\zeta ^{},u^{})}{(\zeta ,u)}=12\mathrm{\Lambda },\mathrm{or}\delta [d\zeta ^{(4)}du]=2\mathrm{\Lambda }[d\zeta ^{(4)}du],$$ $`(3.12)`$ where the HSS superfield parameter $$\mathrm{\Lambda }=\left(\rho +k_{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}x^{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}\right)+\left(\lambda ^{ij}+4i\theta ^{\alpha i}\eta _\alpha ^j+4i\overline{\eta }_\stackrel{_{{}_{}{}^{}}}{\alpha }^j\overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }i}\right)u_i^+u_j^{}$$ $`(3.13)`$ has been introduced. Similarly, one easily finds that $$(D^{++})^{}=D^{++}(D^{++}\mathrm{\Lambda })D^0\mathrm{and}(D^0)^{}=D^0.$$ $`(3.14)`$ The truly (rigid) N=2 superconformal infinitesimal parameters can, therefore, be encoded into the single scalar harmonic superfield $`\mathrm{\Lambda }`$ that is subject to the constraint $$(D^{++})^2\mathrm{\Lambda }=0,$$ $`(3.15)`$ and the reality condition $$\stackrel{}{\overline{(\mathrm{\Lambda }^{++})}}=\mathrm{\Lambda }^{++},\mathrm{where}\mathrm{\Lambda }^{++}D^{++}\mathrm{\Lambda }.$$ $`(3.16)`$ The transformations rules of the harmonics, $$\delta u_i^+=\mathrm{\Lambda }^{++}u_i^{},\delta u_i^{}=0,$$ $`(3.17)`$ together with eqs. (3.12), (3.14), (3.15) and (3.16) yield the very simple and convenient description of rigid N=2 conformal supersymmetry (on the top of N=2 Poincaré supersymmetry) in N=2 HSS. The special hyper-Kähler geometry of the N=2 (rigidly) superconformal NLSM in components was investigated in ref. . A general solution to the special hyper-Kähler geometry in HSS was described in our recent paper . We use the pseudo-real $`Sp(1)`$ notation for a Fayet-Sohnius (FS) hypermultiplet superfield, $$q_a^+=(\stackrel{}{\overline{q}}{}_{}{}^{+},q^+),a=1,2,q^{a+}=\epsilon ^{ab}q_b^+,$$ $`(3.18)`$ which can be easily generalized to the case of several FS hypermultiplets, $`q^{a+}q^{A+}`$ and $`q_A^+=\mathrm{\Omega }_{AB}q^{B+}`$, with a constant (antisymmetric) $`Sp(k)`$-invariant metric $`\mathrm{\Omega }_{AB}`$, $`A,B=1,\mathrm{},2k`$. First, we recall that the most general (rigidly) N=2 supersymmetric NLSM can be formulated in terms of the FS hypermultiplet superfields, $$I_{\mathrm{NLSM}}[q]=\frac{1}{\kappa ^2}𝑑\zeta ^{(4)}𝑑u\left[\frac{1}{2}q_A^+D^{++}q^{A+}+𝒦^{(+4)}(q^{A+},u_i^\pm )\right],$$ $`(3.19)`$ where the real analytic function $`𝒦^{(+4)}=\stackrel{}{\overline{𝒦^{(+4)}}}`$ of $`U(1)`$ charge $`(+4)`$ is known as a hyper-Kähler (pre-)potential <sup>5</sup><sup>5</sup>5The HSS superfields $`q`$ are dimensionless. The dimensionality of the measure in the action (3.19) is compensated by the coupling constant $`\kappa `$ of dimension of length. By manifest N=2 supersymmetry of the NLSM action (3.19), the NLSM metric must be hyper-Kähler for any choice of $`𝒦^{(+4)}`$. Unfortunately, an explicit general relation between a hyper-Kähler potential and the corresponding hyper-Kähler metric is not available (see, however, refs. for the explicit hyper-Kähler potentials of (ALF) multi-Taub-NUT and Atiyah-Hitchin metrics, and their derivation from the NLSM (3.19) in HSS, and ref. for a review or a general introduction into the supersymmetric NLSM). Eq. (3.19) formally solves the hyper-Kähler constraints on the NLSM metric in terms of an arbitrary function $`𝒦^{(+4)}`$ that may be considered as the (analytic) hypermultilet counterpart to the (holomorphic) potential $``$ of abelian N=2 vector superfields in eq. (2.10). It is, therefore, natural to impose extra N=2 superconformal invariance on the action (3.19), in order to determine a general solution to the special hyper-Kähler geometry, since the free part of the action (3.19) is N=2 superconformally invariant . The FS superfields $`q^+`$ have conformal weight one, $`\delta q^+=\mathrm{\Lambda }q^+`$. Together with eqs. (3.12), (3.17) and (3.19) it implies two constraints : $$\frac{𝒦^{(+4)}}{q^{A+}}q^{A+}=2𝒦^{(+4)}\mathrm{and}\frac{𝒦^{(+4)}}{u_i^+}=0.$$ $`(3.20)`$ This means that the special hyper-Kähler potentials are given by homogeneous (of degree two) functions $`𝒦^{(+4)}(q^{A+},u_i^{})`$ of $`q^{A+}`$. There is no restriction on the dependence of $`𝒦^{(+4)}`$ upon $`u_i^{}`$, though there should be no dependence upon $`u_i^+`$. The HSS description of the N=2 superconformal hypermultiplet actions depending upon the FS superfields in terms of a homogeneous (degree 2) potential is thus formally the same as that of the N=2 superconformal (abelian) vector multiplet actions in the standard N=2 (chiral) superspace (sect. 2). However, the special Kähler geometry in the target space of the NLSM arising in the scalar sector of the N=2 superconformal action (2.10) is very different from the special hyper-Kähler geometry arising from the N=2 superconformal NLSM action (3.19) in components. ## 4 Ward identities and supergravity compensators The best way of derivation of the superconformal anomaly relations is based on the use of the supergravity (SG) compensators . To compensate the unwanted (local) N=2 superconformal symmetries in N=2 conformal SG, one needs two compensators . This also implies the existence of two anomaly relations in the N=2 case (cf. sect. 2). In the case of N=2 superfield SG, its most universal formulation is provided by HSS — see Appendix for a short introduction and our notation. Having an N=2 matter action $`I`$ coupled to N=2 SG, one can naturally define an N=2 supercurrent $`𝒥`$ by a variation of the action $`I`$ with respect to the N=2 conformal SG potential $`𝒢`$, $$𝒥=\frac{\delta I}{\delta 𝒢},$$ $`(4.1)`$ in the flat limit where all N=2 conformal SG fields vanish. The N=2 conformal SG potential $`𝒢`$ is defined by eq. (A.4), where it is introduced as the general N=2 harmonic real superfield, subject to the pre-gauge transformations (A.5) and the gauge transformations (A.7) whose linearized form reads $$\delta 𝒢=D^{++}l^{}.$$ $`(4.2)`$ The N=2 superconformal anomalies are also naturally defined in HSS by variational derivatives of the N=2 matter action with respect to the N=2 SG compensators $`v_5^{++}(\zeta ,u)`$ and $`\omega (\zeta ,u)`$ , $$L^{++}=\frac{\delta I}{\delta v_5^{++}},$$ $`(4.3)`$ and $$𝒯^{++++}=\frac{\delta I}{\delta \omega },$$ $`(4.4)`$ where $`v_5^{++}`$ is the real analytic gauge superfield, associated with the central charge and transforming under the (linearized) gauge transformations (A.13) as $$\delta v_5^{++}=D^{++}\lambda _5,$$ $`(4.5)`$ whereas the analytic density $`\omega `$ transforms according to eq. (A.14) that implies (in the linearized approximation) $$\delta \omega =D^{}\mathrm{\Lambda }^{++}=D^{}(D^+)^4l^{},$$ $`(4.6)`$ where we have used the last eq. (A.6). The HSS superfields, $`L^{++}`$ and $`𝒯^{(+4)}𝒯^{++++}`$, representing the N=2 superconformal anomalies, are analytic by their definition. Given the N=2 matter action $`I`$ that is entirely formulated in the ordinary N=2 superspace without harmonics (it is certainly the case in the N=2 gauge sector, without FS hypermultiplets), one can make a connection to the standard N=2 anomaly relation (2.8). The invariance of the action $`I`$ with respect to the gauge transformations (4.2) obviously yields $$D^{++}𝒥=0,$$ $`(4.7)`$ whereas the invariance of the same action with respect to the gauge transformations (4.5) implies $$D^{++}L^{++}=0.$$ $`(4.8)`$ Equation (4.7) means that $`𝒥`$ is independent upon harmonics too, whereas eq. (4.8) implies that $`L^{++}=u_i^+u_j^+L^{ij}(x,\theta ,\overline{\theta })`$, where $`L^{ij}`$ satisfies the N=2 linear multiplet constraints (2.2) due to the analyticity of $`L^{++}`$. Finally, the invariance of the action $`I`$ with respect to the pre-gauge transformations (A.5) and (A.11) yields $$\frac{i}{4}(D^+)^2𝒥=L^{++},$$ $`(4.9)`$ which is equivalent to eq. (2.8), as expected. The HSS results about N=2 SG potentials and compensators imply the natural definitions of the latter in the ordinary N=2 superspace by linear relations, $$G=𝑑u𝒢(\zeta ,u,\overline{\theta }^{})$$ $`(4.10)`$ and $$\mathrm{\Phi }=𝑑u(\overline{D}^{})^2v_5^{++}(\zeta ,u).$$ $`(4.11)`$ The N=2 real superfield $`G(x,\theta ,\overline{\theta })`$ gives the N=2 conformal SG potential (cf. ref. ), whereas the abelian N=2 (restricted chiral) superfield strength squared, $`\mathrm{\Phi }^2`$, can serve as an N=2 (unrestricted) chiral density, i.e. as the N=2 chiral compensator (cf. ref. ). The anomaly relation (2.8) is then the direct consequence of the definitions $$J=\frac{\delta I}{\delta G}\mathrm{and}S=\frac{\delta I}{\delta \mathrm{\Phi }^2}.$$ $`(4.12)`$ Having identified the compensator $`\mathrm{\Phi }`$ with one of the (abelian) N=2 gauge superfield strengths, $`\mathrm{\Phi }=W_1`$, taking the N=2 gauge LEEA (2.10) to represent the N=2 matter action $`I`$ above results in the anomalous N=2 superconformal Ward identity (2.12). The second compensator decouples from the effective N=2 gauge matter action, so that it does not have any impact on its anomaly structure. The situation is just the opposite one in the case where the effective N=2 matter action represents a selfinteraction of FS hypermultiplets. The analytic compensator $`\omega `$ can be considered as a part (density) of the hypermultiplet matter, while the invariance of the most general N=2 matter action in HSS with respect to the gauge transformations (4.2) and (4.6) gives rise to the second anomaly relation in the form $$D^{++}𝒥=D^{}𝒯^{(+4)}.$$ $`(4.13)`$ Being applied to the hypermultiplet LEEA $`I`$ in the form of the N=2 NLSM (3.19) in HSS, eq. (4.13) gives rise to the following anomalous N=2 superconformal Ward identity: $$\underset{A}{}q^{A+}\frac{𝒦^{(+4)}}{q^{A+}}2𝒦^{(+4)}=𝒯^{(+4)}.$$ $`(4.14)`$ This is the key equation in our paper. A generic anomaly $`𝒯^{(+4)}(q,u)`$ is analytic, while it has to be invariant under the unbroken gauge symmetries, internal symmetries, and modular transformations, if any. The anomalous Ward identity (4.14) then gives us the equation on the hyper-Kähler potential $`𝒦^{(+4)}`$ of the effective N=2 NLSM (LEEA). ## 5 Examples To ‘close’ the anomalous N=2 superconformal Ward identity (4.14), as the equation on the effective hyper-Kähler potential $`𝒦(q,u)`$, one has to know the anomaly $`𝒯`$ as a function of $`(q,u)`$ explicitly. In the absence of a general solution for the anomaly (at least, we are unaware ot it), it is worthy to discuss some examples. The non-anomalous symmetries play the major rôle in determining the form of the anomaly $`𝒯(q,u)`$. The crucial simplification arises when the $`SU(2)_R`$ automorphisms of N=2 supersymmetry algebra are not broken (together with the N=2 supersymmetry that we always assume). Since the $`SU(2)_R`$ transformations are linearly realised in HSS, the R-invariance of the hypermultilet LEEA amounts to the independence of the anomaly $`𝒯`$ (and, hence, of the hyper-Kähler potential $`𝒦`$) upon harmonics. Since both have $`U(1)`$ charge $`(+4)`$, the most general (analytic) invariant ‘Ansatz’ is given by a real quartic polynomial of the analytic FS superfields $`q^{A+}`$ , $$𝒯^{(+4)}(q)𝒦^{(+4)}(q)=\lambda _{ABCD}q^{A+}q^{B+}q^{C+}q^{D+},$$ $`(5.1)`$ whose coefficients $`\lambda _{(ABCD)}`$ are totally symmetric, being subject to the reality condition, $`\stackrel{}{\overline{𝒦}}{}_{}{}^{(+4)}=𝒦^{(+4)}`$. Not all of the coefficients in eq. (5.1) are really significant since the FS kinetic terms in eq. (3.19) have the manifest global $`Sp(n)`$ symmetry. It may be not accidental that this $`Sp(n)`$ symmetry coincides with the maximal $`Sp(n)`$ holonomy group of the hyper-Kähler manifolds in $`4n`$ real dimensions. In the case of a single FS hypermultiplet, eq. (5.1) is simplified to $$𝒯^{(+4)}𝒦^{(+4)}=\frac{\lambda }{2}(\stackrel{}{\overline{q}}{}_{}{}^{+})^2(q^+)^2+[\gamma \stackrel{}{\overline{(q^+)}}{}_{}{}^{4}+\beta \stackrel{}{\overline{(q^+)}}{}_{}{}^{3}q_{}^{+}+\mathrm{h}.\mathrm{c}.]$$ $`(5.2)`$ with one real $`(\lambda )`$ and two complex $`(\beta ,\gamma )`$ parameters. The $`Sp(1)`$ transformations of $`q_a^+`$ leave the form of eq. (5.2) invariant, but not its coefficients, which can be used to reduce the number of coupling constants in the family of the hyper-Kähler metrics described by the hyper-Kähler potential (5.2) from five to two. In addition, eq. (5.2) implies the (on-shell) conservation laws $$D^{++}𝒦^{(+4)}=D^{++}𝒯^{(+4)}=0,$$ $`(5.3)`$ which are valid on the equations of motion of the hypermultiplet (FS) superfield, $$D^{++}\stackrel{}{\overline{q}}{}_{}{}^{+}=𝒦^{(+4)}/q^+\mathrm{and}D^{++}q^+=𝒦^{(+4)}/\stackrel{}{\overline{q}}{}_{}{}^{+}.$$ $`(5.4)`$ To understand the physical significance of eq. (5.2), it is instructive to consider first a simpler case, by assuming the additional (translational) $`U(1)_T`$ symmetry that acts on the complex superfields $`(q^+,\stackrel{}{\overline{q}}{}_{}{}^{+})`$ by phase rotations (with a constant parameter $`\alpha `$), $$q^+e^{i\alpha }\theta ^+,\stackrel{}{\overline{q}}{}_{}{}^{+}e^{i\alpha }\stackrel{}{\overline{q}}{}_{}{}^{+},$$ $`(5.5)`$ but does not move the hyper-Kähler structure in the target space of the N=2 NLSM (3.19). It happens, e.g., in the N=2 supersymmetric QED with a single charged hypermultiplet, or in the Coulomb branch of the Seiberg-Witten model . In geometrical terms, the $`U(1)_T`$ symmetry amounts to the existence of a tri-holomorphic (translational) isometry in the N=2 NLSM target space. It is worth mentioning that the $`SU(2)_R`$ isometries are not tri-holomorphic but rotational: they rotate three independent complex structures in the N=2 NLSM hyper-Kähler target space. Given the $`SU(2)_R\times U(1)_T`$ isometry of the N=2 NLSM target space, it must be the symmetry of the NLSM hyper-Kähler potential, which implies further restrictions in eq. (5.2). The unique, $`SU(2)_R\times U(1)_T`$ invariant, hyper-Kähler potential is obviously given by $$𝒦_{\mathrm{Taub}\mathrm{NUT}}^{(+4)}=\frac{\lambda }{2}\left(\stackrel{}{\overline{q}}{}_{}{}^{+}q_{}^{+}\right)^2,$$ $`(5.6)`$ while it is known as the hyper-Kähler potential of the Taub-NUT metric with the mass parameter $`M=\frac{1}{2}\lambda ^{1/2}`$ . The induced coupling constant $`\lambda `$ in the one-loop approximation is determined by the HSS graph depictured in Fig. 1. As was shown in ref. , this graph does lead to the non-vanishing anomalous contribution having the form (5.6), provided that the matter hypermultiplet has a non-vanishing central charge. The wave lines in Fig. 1 denote the analytic propagators of the N=2 (abelian) vector superfields $`V^{++}`$ (in N=2 supersymmetric Feynman gauge) , $$iV^{++}(1)V^{++}(2)=\frac{1}{\mathrm{}_1}(D_1^+)^4\delta ^{12}(𝒵_1𝒵_2)\delta ^{(2,2)}(u_1,u_2),$$ $`(5.7)`$ where $`\delta ^{(2,2)}(u_1,u_2)`$ stands for the harmonic delta-function . The hypermultiplet analytic propagators (the solid lines in Fig. 1) with non-vanishing central charges are given by $$iq^+(1)q^+(2)=\frac{1}{\mathrm{}_1+m^2}\frac{(D_1^+)^4(D_2^+)^4}{(u_1^+u_2^+)^3}e^{\tau _3[v(2)v(1)]}\delta ^{12}(𝒵_1𝒵_2),$$ $`(5.8)`$ where the ‘bridge’ $`v`$ satisfies an equation $`𝒟_Z^{++}e^v=0`$, and $`m^2=\left\|Z\right\|^2`$ stands for the hypermultiplet (BPS) mass. We find $$\lambda =\frac{g^4}{\pi ^2}\left[\frac{1}{m^2}\mathrm{ln}\left(1+\frac{m^2}{\mathrm{\Lambda }^2}\right)\frac{1}{\mathrm{\Lambda }^2+m^2}\right],$$ $`(5.9)`$ where the gauge coupling constant $`g`$ and the IR-cutoff $`\mathrm{\Lambda }`$ have been introduced. It is not difficult to check that $`\lambda 0`$ only if $`Z0`$. The naive ‘non-renormalization theorem’ usually forbids the quantum corrections given by the integrals over a subspace of the full N=2 superspace, like the one in eq. (3.19). However, this ‘theorem’ does not apply here, because of the non-vanishing central charges that give rise to the explicit dependence of the superpropagators (5.8) upon the N=2 superspace Grassmann (anticommuting) coordinates via the bridges $`v(\theta ,\overline{\theta })`$ that are responsible for the N=2 superconformal anomaly. The more general $`SU(2)_R`$-invariant anomaly (5.2) cannot be generated in N=2 perturbation theory, but it can be generated non-perturbatively, due to instanton contributions . Of course, in an abelian N=2 supersymmetric field theory no instantons exist. This means that the N=2 perturbative anomaly described by the Taub-NUT metric is exact in the abelian case. If, however, the underlying N=2 supersymmetric quantum field theory has a non-abelian gauge group of rank larger than one (say, $`SU(3)`$), then one may expect the nonperturbative contributions to the hypermultiplet LEEA (in the Higgs branch) from instantons and anti-instantons that break the $`U(1)_T`$ symmetry. The on-shell relations (5.3) imply that in this case both the anomaly $`𝒯^{(+4)}`$ and the hyper-Kähler potential $`𝒦^{(+4)}`$ can be expressed in terms of a real analytic superfield $`T^{++++}`$ satisfying the constraints $$D^{++}T^{++++}=0\mathrm{and}\stackrel{}{\overline{T}}{}_{}{}^{++++}=T^{++++},$$ $`(5.10)`$ which are the same as the off-shell constraints (2.6) defining an $`O(4)`$ projective superfield $`T^{ijkl}`$, in the ordinary N=2 superspace, $`T^{++++}(\zeta ,u)=u_i^+u_j^+u_k^+u_l^+T^{ijkl}(x,\theta ,\overline{\theta })`$. Unlike the $`O(2)`$ tensor multiplet defined by eq. (2.2), the $`O(4)`$ multiplet does not have a conserved vector amongst its field components. Hence, the $`U(1)`$ isometry, if any, in the N=2 NLSM to be constructed in terms of $`T^{++++}`$, is no longer tri-holomorphic (or translational). The Taub-NUT NLSM (5.6) arises in the limit $`T^{++++}(L^{++})^2`$. The two-parametric family of the hyper-Kähler potentials (5.2) describes the (hyper-Kähler and $`SU(2)_R`$-invariant) deformations of the Atiyah-Hitchin (AH) metric . The N=2 (projective) superspace description of the AH metric in terms of an $`O(4)`$ projective supermultiplet (2.6) was found in ref. . The AH metric is known to be the only regular metric in the family . The ‘difference’ between the AH and Taub-NUT metrics, being considered as the metrics in the quantum moduli space of an N=2 gauge theory (in the region where quantum perturbation theory applies), can be interpreted as the (exponentially small) instanton contributions . Similar remarks are valid in the more general case (5.1) . Another simple example of the N=2 superconformal anomaly, which is still under control, is possible in the case of the unbroken $`U(1)_T`$ symmetry (5.5). It implies that the function $`𝒯^{(+4)}(q^+,\stackrel{}{\overline{q}}{}_{}{}^{+};u)`$ be a function of the invariant product $`(q^+\stackrel{}{\overline{q}}{}_{}{}^{+})`$ and harmonics $`u_i^{}`$ only. The most general ‘Ansatz’ reads $$𝒯^{(+4)}(q^+\stackrel{}{\overline{q}}{}_{}{}^{+};u)=\underset{l=0}{\overset{\mathrm{}}{}}\xi ^{(2l)}\frac{(\stackrel{}{\overline{q}}{}_{}{}^{+}q_{}^{+})^{l+2}}{l+2},$$ $`(5.11)`$ whose harmonic-dependent ‘coefficients’ $`\xi ^{(2l)}`$ are given by $$\xi ^{(2l)}=\xi ^{(i_1\mathrm{}i_{2l})}u_{i_1}^{}\mathrm{}u_{i_{2l}}^{},l=1,2,\mathrm{},$$ $`(5.12)`$ and obey the reality condition $$\stackrel{}{\overline{\xi }}{}_{}{}^{(2l)}=(1)^l\xi ^{(2l)}.$$ $`(5.13)`$ The associated solution to eq. (4.14) takes the similar form (5.11), while it appears to be the hyper-Kähler potential describing the multi-Taub-NUT metrics . The multi-Taub-NUT metrics are known to describe static configurations of several (BPS) monopoles . This is not surprising from the viewpoint of the brane engineering of the effective N=2 supersymmetric gauge field theories in M-theory , where the N=2 field theory hypermultiplets are associated with (parallel) D6-branes whose configurations in M-theory are just described by the multi-Taub-NUT metrics . The brane technology/M-theory also suggest a possibility of yet another generalization, by adding a (parallel) orientifold $`O6^{}`$ to the D6-branes . In geometrical terms, this means replacing the $`U(1)_T`$ translational isometry by a rotational $`U(1)_R`$ isometry in the N=2 NLSM target space. Though the associated N=2 NLSM in HSS were recently constructed in terms of an $`O(4)`$ projective superfield , we are unaware about their explicit reformulation in terms of the FS superfields. We conjecture that this should give rise to the $`D_k`$ series of the asymptotically locally flat (self-dual) metrics in the four-dimensional target space of N=2 NLSM. ## 6 Conclusion In the preceeding section we gave a few non-trivial examples of the N=2 superconformal anomaly $`𝒯^{(+4)}`$ that saturates the hyper-Kähler potential of the effective hypermultilet LEEA given by an N=2 NLSM. To get more (non-perturbative) solutions to our main eq. (4.14), it may be better to find first the organizing principle behind all those solutions. In the Seiberg-Witten case (2.12), it was the underlying Riemann surface or an integrable system. It is natural to expect a similar hidden curve behind the hypermultiplet LEEA too . The $`SU(2)_R`$-invariant effective NLSM metrics in our examples (sect. 5) coincide with the standard metrics in the (BPS) monopole moduli space of the classical $`SU(2)`$-Yang-Mills-Higgs system, with magnetic charge $`n=1`$ (Taub-NUT) or $`n=2`$ (Atiyah-Hitchin). In the N=2 NLSM, whose target space metric is given by the monopole moduli space metric of higher magnetic charge $`n>2`$, the $`SU(2)_R`$ symmetry is necessarily broken . As was shown in ref. , a BPS monopole of magnetic charge $`n>1`$ can always be described by the spectral curve (Riemann surface) of genus $`(n1)^2`$. It is therefore, conceivable that more general solutions to the anomalous N=2 superconformal Ward identity (4.14) are also encoded in terms of the spectral curve, like the Seiberg-Witten-type solutions to eq. (2.12). After a dimensional reduction to three spacetime dimensions, our results support the conjectured mirror symmetry between the Coulomb and Higgs branches . Quantum breaking of conformal symmetry in a (classically) superconformal field theory is related to the appearance of dynamically generated scales. The superconformal compensators can be naturally interpreted as the superfield extensions of the scale parameters. From this physical point of view, the (first) N=2 chiral compensator is apparently related to the N=2 superfield extension of the renormalization scale squared, whereas the second N=2 compensator appears to be the N=2 superfield extension of the induced N=2 NLSM coupling constant that is (roughly) proportional to the inverse central charge squared (sect. 5). It would be interesting to understand better the physical significance of the second N=2 compensator from the viewpoint of the underlying (non-abelian) N=2 gauge theory and brane technology. ## Acknowledgements I am indebted to L. Alvarez-Gaumé, J. Ambjorn, S. Cherkis, S. Gates Jr., A. Gorsky and E. Ivanov for useful discussions. I am also grateful to Niels Bohr Institute in Copenhagen for hospitality extended to me during the completion of this work. ## Appendix: N=2 supergravity in HSS All (rigid) N=2 supersymmetric field theories can be naturally defined in HSS, in terms of unconstrained analytic superfields, with manifest N=2 supersymmetry . N=2 matter (FS) hypermultiplets are described by complex analytic superfields $`q^+`$ of $`U(1)`$ charge $`+1`$, whereas their coupling to N=2 super-Yang-Mills fields is described via the (Lie algebra-valued) extension of the FS hypermultiplet kinetic operator $`D^{++}`$ by a gauge connection $`V^{++}`$, so that the new gauge-covariant operator $`𝒟^{++}=D^{++}+V^{++}`$ preserves analyticity. The rigid N=2 superconformal transformations also preserve the analytic subspace of flat N=2 HSS (sect. 3), so it is natural to define N=2 conformal supergravity (SG) in HSS along the similar lines : by preserving analyticity (of N=2 matter & gauge fields), the analytic conjugation (= the product of usual complex conjugation and Weyl reflection of the sphere $`S^2`$), and unimodularity (of harmonics). In this Appendix we briefly review the HSS formulation of N=2 SG along the lines of ref. — see refs. for more details. Let $`\{\zeta ^M,u,\theta ^{\widehat{\alpha }}\}`$ be the coordinates of the full N=2 HSS in the analytic basis, $`\widehat{\alpha }=(\alpha ,\stackrel{_{{}_{}{}^{}}}{\alpha })`$. The conformal N=2 SG transformations are naturally defined in HSS as the analyticity-preserving diffeomorphisms , \| $`\delta \zeta ^M=`$ \| $`\lambda ^M(\zeta ,u),`$ \| \| --- \| --- \| \| $`\delta u^{i+}=`$ \| $`\lambda ^{++}(\zeta ,u)u^i,\delta u^i=0,`$ \| \| $`\delta \theta ^{\widehat{\alpha }}=`$ \| $`\lambda ^{\widehat{\alpha }}(\zeta ,u,\theta ^{}).`$ \| $`(A.1)`$ Accordingly, the covariant derivative $`𝒟^{++}`$ in N=2 conformal SG can be put into the form $$𝒟^{++}=D^{++}+H^{M++}D_M+H^{(+4)}D^{}+H^{\widehat{\alpha }+}D_{\widehat{\alpha }}^+,$$ $`(A.2)`$ where the SG vielbeine $`(H^{M++},H^{(+4)}H^{++++},H^{\widehat{\alpha }+})`$ have been introduced in front of the flat N=2 HSS covariant derivatives $`D_M=(_{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }},D_{\widehat{\alpha }}^{})`$, $`D^{}`$ and $`D_{\widehat{\alpha }}^+`$. The N=2 conformal SG parameters (A.1) can be similarly organized into the one-forms $$\lambda =\mathrm{\Lambda }^MD_M+\mathrm{\Lambda }^{++}D^{}\mathrm{and}\rho =\rho ^{\widehat{\alpha }}D_{\widehat{\alpha }}^+.$$ $`(A.3)`$ Since the SG derivative $`𝒟^{++}`$ is supposed to preserve analyticity, $`D_{\widehat{\alpha }}^+𝒟^{++}\mathrm{\Phi }^{(+p)}=0`$, of the analytic superfield $`\mathrm{\Phi }^{(+p)}`$ of a positive $`U(1)`$ charge $`p`$, $`D_{\widehat{\alpha }}^+\mathrm{\Phi }^{(+p)}=0`$, the vielbeine of eq. (A.2) have to obey certain linear constraints, whose solution is given by $$H^{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }++}=iD^{\alpha +}\overline{D}^{\stackrel{_{{}_{}{}^{}}}{\alpha }+}𝒢,H^{\alpha +++}=\frac{1}{8}D^{\alpha +}(\overline{D}^+)^2𝒢,H^{(+4)}=(D^+)^4𝒢,$$ $`(A.4)`$ where $`(D^+)^4=\frac{1}{16}(D^+)^2(\overline{D}^+)^2`$, and the real unconstrained (general HSS superfield) N=2 SG pre-potential $`𝒢(\zeta ,u,\theta ^{})`$ is subject to the pre-gauge transformations $$\delta 𝒢=\frac{1}{4}(D^+)^2\mathrm{\Omega }^{}+\frac{1}{4}(\overline{D}^+)^2\stackrel{}{\overline{\mathrm{\Omega }}}^{}$$ $`(A.5)`$ with the complex unconstrained HSS parameter $`\mathrm{\Omega }^{}(\zeta ,u,\theta ^{})`$. Accordingly, the gauge HSS superfield parameters $`(\mathrm{\Lambda }^M,\mathrm{\Lambda }^{++})`$ are to be expressed in terms of a single real unconstrained HSS superfield $`l^{}(\zeta ,u,\theta ^{})`$ , $$\mathrm{\Lambda }^{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}=iD^{\alpha +}\overline{D}^{\stackrel{_{{}_{}{}^{}}}{\alpha }+}l^{},\mathrm{\Lambda }^{\alpha +}=\frac{1}{8}D^{\alpha +}(\overline{D}^+)^2l^{},\mathrm{\Lambda }^{++}=(D^+)^4l^{}.$$ $`(A.6)`$ It is easy to see that $`H^{\widehat{\alpha }+}`$ is pure gauge, so that we can simply ignore it, $`H^{\widehat{\alpha }+}=0`$. Then the transformation law of $`𝒢`$ under the remaining gauge transformations with the parameters $`(\lambda ,l^{})`$ is given by a simple formula, $$\delta 𝒢=\lambda 𝒢+𝒟^{++}l^{}.$$ $`(A.7)`$ In the absence of superconformal anomalies, there exist a (WZ-type) gauge, where the HSS superfield $`𝒢`$ is independent upon harmonics, being subject to the constraints $$D_{ij}𝒢(x,\theta ,\overline{\theta })=\overline{D}_{ij}𝒢(x,\theta ,\overline{\theta })=0.$$ $`(A.8)`$ This equation coincides with eq. (2.7) defining the N=2 conformal supercurrent, so that the independent field components of $`𝒢`$ (in the WZ gauge) are in one-to-one correspondence with the field content of the off-shell N=2 conformal SG, as it should. To construct N=2 matter and gauge couplings in N=2 Poincaré (non-conformal or Einstein) SG, one has to compensate ‘truly’ N=2 superconformal gauge transformations, namely, dilatations, special conformal transformations, $`U(1)`$ chiral and $`SU(2)_{\mathrm{conf}.}`$ rotations, and N=2 special supersymmetry . Some of them are compensated by a real analytic gauge superfield $`V_5^{++}(\zeta ,u)`$ as the first compensator, subject to abelian gauge transformations in HSS, $$\delta V_5^{++}=𝒟^{++}\mathrm{\Lambda }_5,$$ $`(A.9)`$ with the real analytic parameter $`\mathrm{\Lambda }_5(\zeta ,u)`$. In the context of gauging the N=2 supersymmetry algebra, the $`V_5^{++}`$ superfield is associated with the central charge generator $`\widehat{Z}`$. Though the N=2 conformal SG itself has the vanishing central charge, N=2 matter hypermultiplets may have non-vanishing central charges. It is natural to incorporate the central charge generator into the HSS covariant derivatives by redefining $`𝒟^{++}`$ to $`𝒟_Z^{++}=𝒟^{++}+V_5^{++}\widehat{Z}`$, etc. At the component level the N=2 supermultiplet $`V_5^{++}`$ (in a WZ gauge) adds extra $`8_B+8_F`$ off-shell (field) degrees of freedom to the N=2 Weyl multiplet, thus forming together the so-called minimal off-shell N=2 SG multiplet with $`32_B+32_F`$ field components . The extended HSS covariant derivative $`𝒟_Z^{++}`$ also has to preserve analyticity, which implies a linear constraint on $`V_5^{++}`$. A solution to the constraint reads $$V_5^{++}=\frac{i}{4}(D^+)^2𝒢\frac{i}{4}(\overline{D}^+)^2𝒢+v_5^{++},$$ $`(A.10)`$ where $`v_5^{++}`$ is real analytic. The pre-gauge transformations (A.5) are to be appended by $$\delta v_5^{++}=i(D^+)^4(\mathrm{\Omega }^{}\stackrel{}{\overline{\mathrm{\Omega }}}{}_{}{}^{}).$$ $`(A.11)`$ The related restrictions on the HSS parameter $`\mathrm{\Lambda }_5`$ in eq. (A.9) are $$\mathrm{\Lambda }_5=\frac{i}{4}(D^+)^2l^{}\frac{i}{4}(\overline{D}^+)^2l^{}+\lambda _5,$$ $`(A.12)`$ where $`\lambda _5`$ is real analytic. One easily finds $$\delta v_5^{++}=\lambda v_5^{++}+𝒟^{++}\lambda _5.$$ $`(A.13)`$ To compensate the remaining unwanted gauge symmetries (e.g., $`SU(2)_{\mathrm{conf}.}`$), one needs a second compensator that may have either a finite or the infinite number of the auxiliary field components. The three minimal formulations of N=2 Poincaré SG, each having $`40_B+40_F`$ off-shell (field) degrees of freedom, were described in ref. . Unfortunately, all of them impose some restrictions on allowed N=2 matter couplings, which makes them of limited use in the context of generic N=2 NLSM. The most universal choice is given by a real analytic compensator $`\omega `$ that compensates the N=2 local supersymmetry transformations of the analytic measure, $$\omega ^{}(\zeta ^{},u^{})=Ber^1\left(\frac{(\zeta ^{},u^{})}{(\zeta ,u)}\right)\omega (\zeta ,u).$$ $`(A.14)`$ It allows one to accommodate any N=2 NLSM in N=2 SG via the ‘covariantization’ of the (rigid) N=2 NLSM action in HSS, by using the invariant analytic measure $`d\zeta ^{(4)}du\omega `$. It is the absence of an analytic density that is responsible for the restricted N=2 matter couplings in the more conventional formulations of N=2 SG . Coupling to N=2 SG deformes hyper-Kähler geometry of the N=2 NLSM target space into quaternionic geometry . The actions of quaternionic NLSM coupled to N=2 SG in HSS were recently investigated in ref. , where the density $`\omega `$ was constructed in terms of a FS hypermultiplet superfield $`q^+`$ as $`\omega (u_a^{}q^{a+})^2`$. When going in the opposite direction, a rigidly N=2 superconformal (special hyper-Kähler) NLSM arises from the quaternionic one after putting all the N=2 conformal SG fields to zero, $`𝒢=0`$, together with the vanishing Maxwell multiplet, $`v_5^{++}=0`$, and $`\omega =1`$.
warning/0003/astro-ph0003441.html	ar5iv	text	# The violent past of Cygnus X–2 ## 1 Introduction The $`9.84`$ d period neutron–star binary Cygnus X–2 has long been regarded as an archetypal long–period low–mass X–ray binary (Cowley et al. 1979, Webbink et al. 1983) where nuclear expansion of a Hayashi line giant drives mass transfer. Yet recent optical photometry and high–resolution spectroscopy, while confirming that the donor has a low mass (Casares et al. 1998; Orosz & Kuulkers 1999), unambiguously showed that the spectral type of the optical counterpart is A$`5\pm 2`$ (Casares et al. 1998), significantly too hot for a Hayashi line donor. King & Ritter (1999; hereafter KR) and Podsiadlowski & Rappaport (2000; PR) argued that the system must be the descendent of an intermediate–mass X–ray binary (IMXB), with $`34\mathrm{M}_{}`$ as the likely initial donor mass. Such systems undergo a rapid mass transfer phase, previously regarded as fatal, with transfer rates exceeding the Eddington value by several orders of magnitude. For a narrow range of initial separations the systems never reach the Hayashi line, or evolve away from it, during the subsequent phase with slower mass transfer. KR suggested that Cygnus X–2 is the product of a “case B” mass transfer sequence, where the donor star was already expanding towards the giant branch when mass transfer began (Kippenhahn & Weigert ). By contrast, PR preferred an evolution where mass transfer started while the donor was still on the main sequence, with core hydrogen burning terminating during the mass transfer phase. KR’s considerations are semi–analytical and entirely based on generalised main–sequences (Giannone et al. 1968), while PR performed detailed binary sequences with a full stellar code. In Sect. 2 and 3 we use our evolutionary code to reexamine critically the case B solution rejected by PR, and resolve the discrepancy between KR and PR. To see how common a Cygnus X–2–like evolution is we consider IMXBs with still higher initial donor masses in Sect. 4. A further peculiarity of Cygnus X–2, addressed in Sect. 5, is its dynamical state. The system has a measured line-of-sight velocity of -208.6 km/s (Casares et al. 1998) and a Galactic latitude of -11.32 and longitude of 87.33. At a distance of 11.6 kpc from the sun (Smale 1998), this places Cygnus X–2 at a Galactocentric distance of 14.2 kpc, and a distance from the Galactic plane of 2.28 kpc. Integrating the equations of motion in the Galactic potential we use Monte Carlo techniques to investigate possible trajectories of Cygnus X–2. ## 2 Constraining the prehistory of Cygnus X–2 The observed location of the donor star of Cygnus X–2 in the middle of the Hertzsprung gap in the HR diagram, and its small mass ($`0.40.7\mathrm{M}_{}`$), suggest that the mass of the hydrogen–rich envelope remaining above the donor’s helium core is very small ($`M_\mathrm{H}/M_{\mathrm{He}}0.05`$). Stars with an initial mass $`34.5\mathrm{M}_{}`$ leave the core hydrogen burning phase with a helium core mass $`M_{\mathrm{He}}`$ in this range (e.g. Bressan et al. 1993). Mass transfer from such an intermediate–mass star on to a less massive neutron star is thermally unstable and involves an initial phase of rapid mass transfer with a highly super–Eddington transfer rate. This rapid phase lasts roughly until the mass ratio is reversed and the donor’s Roche lobe expands upon further mass transfer. The subsequent slower transfer phase proceeds on the donor’s thermal timescale, which is essentially given by the Kelvin–Helmholtz time of the donor when it left the main sequence. PR identified two possible routes from the post–supernova binary, i.e. after the formation of the neutron star, to the present system configuration. One route is via a genuine case B evolution, the other via an evolution they labelled ‘case AB’. In the former the donor star has already left the main sequence when it fills its Roche lobe for the first time. The mass transfer phase is short–lived, with most of the mass transferred at a highly super–Eddington rate. Hence the neutron star mass increase is negligible, even if it accretes at the Eddington rate during the whole evolution. In the case AB solution discussed by PR mass transfer already starts during the donor’s main–sequence phase. The system briefly detaches when core hydrogen burning terminates, but mass transfer resumes when shell–burning is well established. As the rapid mass transfer phase terminates before the donor leaves the main sequence the subsequent slow phase proceeds on a much longer timescale than in the genuine case B solution — the thermal time of the now less massive donor at the terminal main–sequence. The transfer rate is sub–Eddington for some time and the neutron star grows in mass. Given this we explore three different prescriptions to project the evolution of Cygnus X–2 backwards in time: (Model 1) To represent a case B evolution we assume that $`M_1=`$const., and that any material lost from the system carries the specific orbital angular momentum of the neutron star. Then the initial orbital separation $`a_i`$ is $$a_i=a\frac{M_1+M_2}{M_1+M_{2i}}\left(\frac{M_2}{M_{2i}}\right)^2\mathrm{exp}\left(\frac{2(M_{2i}M_2)}{M_1}\right),$$ (1) where $`M_{2i}`$, $`M_2`$ denote the initial and present donor mass, and $`a`$ the present separation (e.g. KR). This solution is valid only if $`a_i>a_{\mathrm{TMS}}`$, the orbital separation of a binary with a Roche–lobe filling terminal main–sequence (TMS) star of mass $`M_{2,i}`$. (Model 2) The same applies to a case AB evolution, except that we require $`a_i<a_{\mathrm{TMS}}`$. (Model 3) To account for a possible prolonged phase with sub–Eddington mass transfer in a case AB sequence we allow mass transfer to be conservative (total binary mass and orbital angular momentum is constant) for the last part of the evolution. We assume that the neutron star mass at birth was $`1.4\mathrm{M}_{}`$, and that the present neutron star mass is $`>1.4\mathrm{M}_{}`$. Hence the evolution backwards in time consists of two branches. During the conservative phase, characterised by $`aM_1^2M_2^2=`$constant, $`M_1`$ reduces from its present value to $`1.4\mathrm{M}_{}`$. The second phase is calculated with constant neutron star mass ($`1.4\mathrm{M}_{}`$), as in models 1 and 2 above. The present system parameters of Cygnus X–2 are: orbital period $`P_0=9.844`$ d, mass ratio $`q=0.34\pm 0.04`$ (Casares et al. 1998), and neutron star mass $`M_1=1.41.8\mathrm{M}_{}`$. This mass range accommodates the canonical neutron star mass at birth, as well as the estimate by Orosz & Kuulkers (1999) based on modelling the observed ellipsoidal variations. If we adopt a specific value $`M_1`$ for the present neutron star mass, then the present donor mass is constrained to the range $`0.30M_1M_20.38M_1`$ by the observed mass ratio. This translates into a range of allowed initial separations, as shown by the hatched regions in Fig. 1, for case B (model 1, wide spacing) and case AB (model 2, narrow spacing) evolution. The parameter space for case AB solutions with model 3 assumptions is tiny and hence not shown in the figure. For conservative mass transfer (model 3) the orbital separation of Cygnus X–2 today increases less steeply with decreasing donor mass than for the isotropic wind case (model 2). Hence the separation for model 3 sequences describing the past evolution of Cygnus X–2 is generally larger than for the corresponding model 2 sequences. In particular, model 3 solutions require larger initial separations. Most of them are in conflict with the limit $`a_i<a_{\mathrm{TMS}}`$, i.e. inconsistent with the assumption that the donor was on the main sequence when mass transfer started. More generally, if the mass lost from the system in sequences with $`M_1=`$ const. carries more (less) specific angular momentum than that of the neutron star, the required initial separation is larger (smaller) than estimated by (1). A somewhat higher loss seems more likely than a smaller loss, hence this reduces the parameter space available for case AB solutions. Even if Fig. 1 indicates a viable solution it is not clear if the corresponding evolutionary sequence reproduces Cygnus X–2. Although the donor will have the observed mass (and hence radius) at the observed orbital period $`P_0`$, the effective temperature is of course not constrained by the above considerations. To check this we need a detailed binary sequence with full stellar models, as presented in the next section. In particular, it appears that the case AB sequences described by PR require the donor star to be already fairly close to the end of core hydrogen burning at turn–on of mass transfer. This is likely to narrow the parameter space for case AB sequences even further. We note that the limits shown in Fig. 1 depend via $`a_{\mathrm{TMS}}`$ somewhat on the stellar input physics, in particular on the amount of convective core overshooting (cf. the discussion in the next section). ## 3 Model calculations We calculated several detailed early massive case B binary evolution sequences, using Mazzitelli’s stellar code in its 1989 version (see Mazzitelli 1989, and references therein, for a summary of the input physics). Mass transfer was treated as in Kolb (1998). Table 1 summarises initial and final system parameters. The neutron star mass was $`1.4\mathrm{M}_{}=`$ const. in each case, with angular momentum loss treated as in model 1 above. Specifically, we present in detail the evolution of a system with initial donor mass $`M_2=3.5\mathrm{M}_{}`$ and orbital separation $`a_i=13.3R_{}`$ (Sequence I) and $`11.1R_{}`$ (sequence II). The figures 2-4 confirm the general behaviour of case B evolutionary sequences as described above. Sequence I gets very close to the observed state of Cygnus X–2, although it fails to fit all observed parameters simultaneously with high accuracy. The evolutionary track in the HR diagram passes through the Cygnus X–2 error box, but at a period slightly longer than the observed value $`P_0=9.844`$ d. This can be seen in Fig. 4, which shows that at $`P_0`$ the sequence I donor is somewhat too cool, while in sequence II it is somewhat too hot. Hence the initial separation for a simultaneous fit of $`P`$, $`T_{\mathrm{eff}}`$ (and $`L`$) lies between $`11.1R_{}`$ and $`13.3R_{}`$. The transfer rate at $`P_0`$ for sequence I is $`5\times 10^8\mathrm{M}_{}\mathrm{yr}^1`$, a factor $`23`$ higher than the Eddington rate. Observational estimates place the actual accretion rate in Cygnus X–2 close to the Eddington limit. With no evidence of significant mass loss from the system we expect that the transfer rate is also of the order of the Eddington rate. As the transfer rate in the model sequence is set by the thermal time of the progenitor star, a somewhat smaller initial secondary mass (e.g. $`3.0\mathrm{M}_{}`$) would give a lower rate; this has already been noted by PR. We note the following main differences between our sequence I and the case B sequence calculated by PR: (1) PR chose $`21.3R_{}`$ as the initial separation. This leads to the combination $`1.4+0.66\mathrm{M}_{}`$ for the component masses at $`P=P_0`$, i.e. to a mass ratio $`0.47`$, significantly larger than observed. Hence the PR case B sequence represents a poor fit for Cygnus X–2. (2) More importantly, the mass transfer rate in the slow phase of sequence I decreases from $`10^7\mathrm{M}_{}\mathrm{yr}^1`$ to $`10^8\mathrm{M}_{}\mathrm{yr}^1`$ over a period of 2.5 Myr, while in PR’s model the transfer rate is always larger than $`10^7\mathrm{M}_{}\mathrm{yr}^1`$ (which is inconsistent with observations). A similar calculation by Tauris et al. (2000) with essentially the same stellar code as the one used by PR gives a similarly high rate in the slow phase. The initial separation for our sequence II is essentially the same as the one used by PR for their case AB evolution ($`11.5R_{}`$, corresponding to a stellar radius $`5.6R_{}`$ of the donor at turn–on of mass transfer). The most important difference between the two calculations is the degree of convective overshooting assumed during the main–sequence phase: none in our models, a very strong one in PR’s models. Hence PR’s main–sequence band is much wider, with $`6.6R_{}`$ as the maximum radius $`R_{2,\mathrm{max}}`$ that a $`3.5\mathrm{M}_{}`$ star reaches during core hydrogen burning, compared to $`3.6R_{}`$ in our case. A moderate extent of core overshooting is favoured in the literature, giving $`R_{2,\mathrm{max}}=5.0R_{}`$ (Schaller et al. 1992), $`5.7R_{}`$ (Bressan et al. 1993) and $`4.8R_{}`$ (Dominguez et al. 1999) for a $`3.5\mathrm{M}_{}`$ star with solar composition. Schröder et al. (1997) — who use the same code and input physics as PR — prefer a rather efficient overshooting leading to values up to $`R_{2,\mathrm{max}}=7.1R_{}`$, while in a subsequent paper (Pols et al. 1997) the same authors concluded that $`R_{2,\mathrm{max}}=5.9R_{}`$ (which corresponds to their overshooting parameter $`\delta =0.12`$) provides the best overall fit to observations. This is only barely larger than the donor’s radius at turn–on of mass transfer in PR’s case AB solution. Hence standard assumptions about overshooting favour a case B solution for Cygnus X–2 over a case AB solution. We conclude that a case B solution is a viable fit for the evolutionary history of Cygnus X–2. Given the parameter space limitations for a case AB evolution it does appear as the more likely solution for Cygnus X–2, even though this implies that the presently observed state of Cygnus X–2 is rather short–lived, of order several million years. PR pointed out that the surface composition of the donor in Cygnus X–2 would be significantly hydrogen–depleted ($`X0.1`$) and show signs of CNO–processing if it had undergone a case AB evolution. Unfortunately, this does not distinguish unambiguously between a case AB and case B evolution, as the same is in principle true for a donor that had undergone a case B mass transfer. The models in our sequence I close to the position of Cygnus X–2 have a surface hydogen mass fraction of $`X=0.29`$, while C/N and O/N are close to the equilibrium values for CNO burning at $`2\times 10^7`$ K. ## 4 Sequences with initially more massive donor stars The calculations discussed so far represent early case B mass transfer solutions for systems with a moderate initial mass ratio, $`q=M_2/M_12.5`$. Systems with smaller mass ratio ($`q1`$) have no initial rapid mass transfer phase and can be understood semi–analytically (see e.g. Kolb 1998, Ritter 1999), while the fate of systems with yet larger initial mass ratio $`3`$ is unclear. Hjellming (1989) pointed out that sustained thermal–timescale mass transfer can lead to a delayed transition to dynamical–timescale mass transfer. The reason for this is that the adiabatic mass–radius index of initially radiative stars decreases significantly when the radiative envelope is stripped rapidly. This causes the Roche lobe to shrink faster than the star. Estimates for the critical initial mass ratio which just avoids the delayed dynamical instability give $`q34`$ (Hjellming 1989; Kalogera & Webbink 1996), although none of these are based on self–consistent mass transfer calculations. This critical value is relevant for identifying possible descendants of neutron–star systems undergoing early massive case B mass transfer. These should appear as binary millisecond pulsars lying significantly below the orbital period–white dwarf mass relation found by Rappaport et al. (1995; see also Tauris & Savonije 1999 for calculations with updated input physics) for systems descending from Hayashi line low–mass X–ray binaries. Obvious candidates are systems with a fairly high–mass white dwarf but short orbital period. The most discrepant system is B0655+64 ($`P=1.03`$ d, $`M_{\mathrm{WD}}0.7`$; cf. KR). The fairly massive white dwarf implies a donor mass $`5\mathrm{M}_{}`$ in the progenitor binary. We performed test calculations similar to our sequences I and II, but with initial donor mass $`3.75\mathrm{M}_{}`$, $`4.0\mathrm{M}_{}`$ and $`5.0\mathrm{M}_{}`$. While the $`3.75\mathrm{M}_{}`$ sequence was stable throughout, the $`4.0\mathrm{M}_{}`$ and $`5.0\mathrm{M}_{}`$ sequences encountered runaway mass transfer (where we stopped the calculations) rather early in the rapid mass transfer phase (see Tab.1). An additional test sequence with a $`3\mathrm{M}_{}`$ donor near the end of core hydrogen burning and initial mass ratio 4 (i.e. primary mass $`0.75\mathrm{M}_{}`$, assumed constant, as above) encountered runaway mass transfer at a donor mass $`2.6\mathrm{M}_{}`$. This is in perfect agreement with Hjellming’s prediction, based on his Fig. IV.1 and a Roche lobe curve corresponding to $`q=4`$. Unless a neutron star binary can survive even such a dynamical–timescale mass transfer, the phase space available for a Cygnus X–2–like evolution is severely limited by this upper limit on the initial donor mass. Our calculations suggest that if most neutron stars form with $`1.4\mathrm{M}_{}`$ the maximum white dwarf mass in an endproduct of early massive case B evolution is $`0.55\mathrm{M}_{}`$. If this is true, neither B0655+64 nor the recently discovered system J1453-58 ($`P=12.422`$ d, $`M_{\mathrm{WD}}0.88\mathrm{M}_{}`$; cf. Manchester et al. 1999) can have formed in this way. The occurrence of the delayed dynamical instability is intimately linked to the internal structure of the donor star. Therefore it is not surprising that the maximum initial donor mass for early case B mass transfer, just avoiding this instability, depends on the stellar input physics. Using an updated version of Eggleton’s stellar code (see e.g. Tauris & Savonije 1999), Tauris et al. (2000) found $`5\mathrm{M}_{}`$ for this limit, with a corresponding maximum white dwarf mass $`0.8\mathrm{M}_{}`$ in millisecond pulsar binaries formed in this way. ## 5 The trajectory of Cygnus X–2 As described in the introduction, Cygnus X–2 is in an unusual dynamical state. In this section we investigate possible trajectories of Cygnus X–2, given the observed line–of–sight velocity and the constraints on the evolutionary state as derived in Secs. 2 and 3. In Fig. 5 we plot the cumulative velocity distribution for systems with initial separations of $`d=15R_{}`$, i.e. immediately after circularisation of the binary orbit following the SN explosion producing the neutron star. The secondary is taken to have a mass $`M_2=3.5`$ M while the mass of the helium star prior to the supernova is $`M_1=5.0`$ M. We consider two different distributions for the kick velocity imparted to the neutron star; namely those by Hansen & Phinney (1997) and Fryer (1999). Throughout the following section, when we refer to kick velocities we mean the kick imparted on the binary allowing for the effects of mass–loss and an asymmetric supernova, and not simply the kick imparted on the neutron star from the latter. Given its present position and velocity, we integrate the trajectory of Cygnus X–2 backwards to the birth of the neutron star using a model for the Galactic potential suggested by Paczyński (1990) (reviewed in the Appendix). From sections 2 and 3 it is clear that the age $`t`$ of the neutron star is essentially the main–sequence lifetime of the progenitor star of the present donor. To account for the allowed range of the initial donor mass and for uncertainties from the width of the main–sequence we choose $`150<t/\mathrm{Myr}<275`$. As only the line-of-sight velocity, $`𝐯_{los}`$, is known, we considered a set of trajectories where the current velocity was given by $$𝐯=𝐯_{}+𝐯_{los}+\alpha 𝐯_1+\beta 𝐯_2$$ (2) where $`𝐯_{los}`$, $`𝐯_1`$, and $`𝐯_2`$ are mutually orthogonal, and $`𝐯_1`$ has a zero component in the $`z`$ direction. When computing trajectories, it is most convenient to work in units related to the mass and size of the galaxy. As a natural unit of velocity we use 207 km/s, the Kepler velocity at 1 kpc distance from a point mass $`10^{10}`$ M. In the above equation $`𝐯_1`$, and $`𝐯_2`$ are given unit lengths in our code units. We integrated the equations (A9) for a range of values of $`\alpha `$ and $`\beta `$ and noted the radius, $`r`$, and time, $`t`$, whenever the trajectory cut the Galactic plane. We also computed the kick velocity, $`V_\mathrm{k}`$ the system must have had if the neutron star has formed at that time, assuming the system had previously been in a circular orbit in the Galactic plane. In Fig. 6 we plot the values of $`r`$ and $`V_\mathrm{k}`$ for all values of $`\alpha `$ and $`\beta `$. We note that for $`r<10`$ kpc, 140 km/s $`<V_\mathrm{k}<`$ 600 km/s. This range of required values of $`V_\mathrm{k}`$ should be compared with that expected given the neutron star kick distributions of Hansen & Phinney (1997) and Fryer (1999) in Fig. 5. It is clear from Fig. 5 that the expected value of $`V_\mathrm{k}`$ is far lower than that required for many of the trajectories shown in Fig. 6. In other words, in a large fraction of systems formed at a radius $`r<10`$ kpc, the trajectory of the binary will confine the system to the inner regions of the Galaxy; Cygnus X–2 must have been relatively unusual in reaching out beyond the solar circle. We plot the time of formation, $`t`$, as a function of $`r`$ for all values of $`\alpha `$ and $`\beta `$ in Fig. 7. As in Fig. 6, the data points where $`150`$ Myr $`<t<275`$ Myr and $`V_k200`$ km/s are plotted as larger dots. From the work of Sections 2 and 3 these are the values most likely to be applicable to a progenitor of Cygnus X–2. We note from Fig. 7 that the binary is unlikely to have originated within 5 kpc of the Galactic centre as trajectories cutting the Galactic plane at such small radii do so at the wrong time. Fig. 6 also shows they would require unreasonably large kicks. In Fig. 8 we plot the values of $`\alpha `$ and $`\beta `$ for those trajectories which cut the Galactic plane $`t`$ Myr ago, where $`150<t<275`$, with system kick velocities required to take the binary from a circular orbit restricted by $`V_k200`$ km/s. Values of $`\alpha `$ and $`\beta `$ besides those shown in Fig. 8 were considered, but none satisfied the above conditions. We see that only a small number of possible trajectories are credible with these restrictions applied. Here we note the effect of the Galactic rotation. For $`\alpha <0.4`$ the trajectories are in the opposite direction to the Galactic rotation. In other words the kick the system receives on formation will oppose its initial velocity in a circular orbit, and will therefore have to be larger than for a similar trajectory in the same direction as the Galactic rotation. Here the effect of the kick is boosted by the initial velocity of the circular orbit. For illustration, in Fig. 9 we plot the trajectory ($`R`$ and $`z`$) for $`\alpha =0.75`$ and $`\beta =0.2`$. In this case, the trajectory cuts the Galactic plane on two occasions within the last 250 Myr, the first occasion being 95 Myr ago, at a radius, $`R_{\mathrm{cut}}=14.1`$ kpc, and the second occasion being 172 Myr ago, at a radius , $`R_{\mathrm{cut}}=5.9`$ kpc. If the system was formed in the disc at this latter time, and was initially travelling on a circular orbit, it would have had to receive a kick, $`V_\mathrm{k}130`$ km/s, which sits comfortably within the range given in Fig. 5. Thus far we have demonstrated that there are viable trajectories for Cygnus X–2 compatible with the observed line-of-sight velocity. We now ask the reverse question: what fraction of binaries originating at some radius will produce systems resembling Cygnus X–2? Given that we expect systems to have originated within the solar circle, as most massive stars are located within 10 kpc of the Galactic centre, we might suspect that Cygnus X–2 is relatively unusual in being at a larger distance from the Galactic centre and significantly away from the Galactic plane. By Monte Carlo simulation we were able to produce a large number of systems at various initial radii, assumed to be initially on circular orbits, with kick velocities drawn from the distributions given in Fig. 5, and follow their trajectories in the Galactic potential. Investigation demonstrated that in determining whether a particular trajectory would produce a Cygnus X–2 like system, the maximum radius reached by the system, $`R_{\mathrm{max}}`$ was a good diagnostic. This is illustrated in Fig. 10 where we plot the radii of 100 binaries as a function of time where the initial trajectories are drawn randomly, from the kick distribution of Hansen and Phinney (see Fig. 5), and choosing the initial radius to be between 2 kpc and 10 kpc. The trajectories fall into two categories: those where the binary remains at a radius similar to that at which it was located initially, and those where the binary is ejected significantly into the Galactic halo, reaching maximum radii of $`>10`$ kpc. Cygnus X–2 clearly belongs to the latter category. In order for a system to be at a radius today similar to that of Cygnus X–2, we require $`15<R_{\mathrm{max}}<20`$ kpc, providing the system originated in the Galactic disc somewhere within 10 kpc of the Galactic centre. A number of systems will be ejected to even larger radii. In such cases a binary of age similar to Cygnus X–2 would still be on the outward bound portion of its orbit. In Fig. 11 we plot $`R_{\mathrm{max}}`$ as a function of initial radius $`R`$. In this case the initial radius was chosen randomly from 2 kpc to 10 kpc, and the kick velocity drawn from the distribution given by Hansen and Phinney (see Fig. 5). From this figure we note that Cygnus X–2–like objects may be produced when $`r>5`$ kpc, and that the relative frequency is relatively independent of formation radius. Systems having larger values of $`R_{\mathrm{max}}`$ will also be produced, the frequency increasing with radius of formation. The relative frequency for Cygnus X–2–like systems and those on longer period orbits is plotted as a function of $`R`$ in Fig. 12. Assuming that the formation rate of binaries scales as the surface density of stars, which is given by $`\mathrm{\Sigma }\mathrm{\Sigma }_0e^{r/3\mathrm{k}\mathrm{p}\mathrm{c}}`$ (see Binney & Merrifield 1999), we computed the relative number of Cygnus X–2–like systems which would be produced by integrating over the entire disc within the solar circle. The fraction of systems resembling Cygnus X–2 ($`f_{\mathrm{cyg}}`$) and the total fraction of systems which will travel significantly further out than their formation radius ($`f_{\mathrm{out}}`$) are listed in Table 2 as a function of the initial separation within the post-supernova binary, $`d`$, the primary helium star mass prior to the supernova, $`M_1`$, and the mass of the secondary, $`M_2`$. We found that $`7`$% of all binaries will produce Cygnus X–2–like binaries on trajectories that would place them at a radius similar to Cygnus X–2 today. A further $`515`$% will be further out than Cygnus X–2, whilst the remainder will be located closer to the Galactic centre. These results apply equally to both velocity distributions plotted in Fig. 5. ## 6 Summary Our binary evolution calculations with full stellar models have verified that Cygnus X–2 can be understood as the descendant of an intermediate–mass X–ray binary (IMXB). The most likely evolutionary solution is an early massive case B sequence, starting from a donor with mass $`3.5\mathrm{M}_{}`$. This implies that the presently observed state is rather short–lived, of order 3 Myr. This in turn points to a large IMXB formation rate of order $`10^710^6`$ yr<sup>-1</sup> if Cygnus X–2 is the only such system in the Galaxy. It is likely that there are more as yet unrecognized IMXBs, so the Galactic IMXB formation rate could be even higher. Hercules X–1 seems to be in the early stages of a similar case A or case AB evolution (e.g. van den Heuvel 1981). The prehistory of Cygnus X–2 is sensitive to the width of the main–sequence band in the HR diagram, i.e. to convective overshooting in that phase. The alternative evolutionary solution for Cygnus X–2 suggested by PR, a case A mass transfer followed by a case B phase, is viable only if overshooting is very effective, i.e. more effective than hitherto assumed in the literature. Using Mazzitelli’s stellar code (Mazzitelli 1989) we found that neutron star IMXBs that start case B mass transfer with initial donor mass $`4\mathrm{M}_{}`$ will encounter a delayed dynamical instability. (Note that the maximum donor mass that just avoids this instability depends on stellar input physics). The components are likely to merge and perhaps form a low–mass black hole. Given the high formation rate of Cygnus X–2–like objects the formation rate of such black holes could be rather substantial. We have shown that the large Galactocentric distance of Cygnus X–2 and its high negative radial velocity do not require unusual circumstances at birth of the neutron star. There are viable trajectories for Cygnus X–2 that firstly cut the Galactic plane inside the solar circle at a time consistent with the evolutionary age of Cygnus X–2, and secondly which require only a moderate neutron star kick velocity ($`<200`$km/s). We estimate the fraction of systems resembling Cygnus X–2, i.e. of systems on orbits that reach Galactocentric distances $`>15`$ kpc, as $`1020\%`$ of the entire IMXB population of systems. ### Acknowledgements MBD gratefully acknowledges the support of a URF from the Royal Society. ARK thanks the UK Particle Physics & Astronomy Research Council for a Senior Fellowship. This work was partially supported by a PPARC short–term visitors grant. We thank the anonymous referee for a careful reading of the manuscript and for comments that helped to improve the paper. APPENDIX: THE GALACTIC POTENTIAL The Galactic potential can be modelled as the sum of three potentials. The spheroid and disc components are given by $$\mathrm{\Phi }_s(R,z)=\frac{GM_s}{\left(R^2+[a_s+(z^2+b_s^2)^{1/2}]^2\right)^{1/2}}$$ (3) $$\mathrm{\Phi }_d(R,z)=\frac{GM_d}{\left(R^2+[a_d+(z^2+b_d^2)^{1/2}]^2\right)^{1/2}}$$ (4) where $`R^2=x^2+y^2`$. The component from the Galactic halo can be derived assuming a halo density distribution, $`\rho _\mathrm{h}`$, given by $$\rho _h=\frac{\rho _c}{1+(r/r_c)^2}$$ (5) where $`r^2=x^2+y^2+z^2`$. The above density distribution yields the potential $$\mathrm{\Phi }_h=\frac{GM_c}{r_c}\left[\frac{1}{2}\mathrm{ln}\left(1+\frac{r^2}{r_c^2}\right)+\frac{r_c}{r}\mathrm{atan}\left(\frac{r}{r_c}\right)\right]$$ (6) where $`M_c=4\pi \rho _cr_c^3`$. The total Galactic potential is the sum $$\mathrm{\Phi }=\mathrm{\Phi }_s+\mathrm{\Phi }_d+\mathrm{\Phi }_h$$ (7) Following Paczyński (1990), we use the following choice of parameters: $$a_s=0,b_s=0.277\mathrm{kpc},M_s=1.12\times 10^{10}M_{},$$ (8) $$a_d=3.7\mathrm{kpc},b_d=0.20\mathrm{kpc},M_d=8.07\times 10^{10}M_{},$$ (9) $$r_c=6.0\mathrm{kpc},M_c=5.0\times 10^{10}M_{},$$ (10) Because of the cylindrical symmetry of the potential, the integration of the trajectories can be simplified to consider the evolution of the $`z`$ and $`R`$ components only, as given below $`{\displaystyle \frac{dR}{dt}}`$ $`=`$ $`v_R,{\displaystyle \frac{dz}{dt}}=v_z,`$ $`{\displaystyle \frac{dv_R}{dt}}`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{\Phi }}{R}}\right)_z+{\displaystyle \frac{j_z^2}{R^3}},{\displaystyle \frac{dv_z}{dt}}=\left({\displaystyle \frac{\mathrm{\Phi }}{z}}\right)_R`$ (11) where the $`z`$ component of the angular momentum, $`j_z=Rv_\varphi `$.
warning/0003/cond-mat0003270.html	ar5iv	text	# Moments of the ARPES spectral function of an undoped Mott insulator. ## I The Calculation ### A Half-filled Hubbard model at large U The large U Hubbard model at half filling is the simplest model that produces the large charge gap and the antiferromagnetic tendencies observed experimentally in CuO planes. The Hamiltonian, $$\widehat{H}=\underset{\sigma ,i,j}{}t_{ij}C_{\stackrel{}{i},\sigma }^{}C_{\stackrel{}{j},\sigma }+U\underset{\stackrel{}{j}}{}n_\stackrel{}{j}^\sigma n_\stackrel{}{j}^\sigma $$ (1) contains a strong onsite repulsion term $`U_\stackrel{}{j}n_\stackrel{}{j}^\sigma n_\stackrel{}{j}^\sigma `$ that freezes (together with the lattice potential) charge motion at half-filling provided that the kinetic energy term is small enough to be treated as a perturbation, $`\|t_{ij}\|U`$. When $`t_{ij}/U=0`$ the ground state consists of singly occupied sites, has energy zero, is $`2^N`$-fold degenerate where N is the size of the system, and is completely characterized by the spin configuration, $`\{S_i\}`$. The perturbative effects of the kinetic energy modify all of the above statements except the last one: each of the perturbed states, though it contains an admixture of doubly occupied sites, can still be labeled by the spin configuration of the unperturbed state from which it has evolved. Thus the expectation value of any operator in the ground state manifold can always be expressed in terms of spin variables. Formally this is accomplished by computing perturbatively in powers of $`t_{ij}/U`$ the unitary transformation $`\mathrm{exp}[i\widehat{X}]`$ that expresses the evolution of the low energy unpeturbed states $`\|\{S_i\}>`$ as a function of increasing kinetic energy: $`\stackrel{~}{\|\{S_i\}>}`$ $`=`$ $`e^{i\widehat{X}}\|\{S_i\}>`$ (2) $`<\widehat{𝒪}>`$ $`=`$ $`<\{S_i\}\|e^{i\widehat{X}}\widehat{𝒪}e^{i\widehat{X}}\|\{S_i\}>,`$ (3) where $`𝒪`$ is any observable. In the familiar fashion, this transformation maps the low energy physics of the Hubbard Hamiltonian into an effective Heisenberg antiferromagnet whose leading term is: $`\widehat{H}_{eff}`$ $`=`$ $`e^{i\widehat{X}}\widehat{H}e^{i\widehat{X}}`$ (4) $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j}{}}{\displaystyle \frac{4t_{ij}^2}{U}}({\displaystyle \frac{1}{4}}\stackrel{}{S_i}\stackrel{}{S_j})+𝒪(t_{ij}^4/U^3).`$ (5) ### B The spectral function The emission spectral function is $`A(\stackrel{}{k},w)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma ,m,n}{}}e^{\beta E_n}\|<m\|C_{\stackrel{}{k},\sigma }^{}\|n>\|^2\delta (\omega +E_mE_n)`$ (6) $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma }{}}{\displaystyle 𝑑te^{iwt}<C_{\stackrel{}{k},\sigma }^{}(t)C_{\stackrel{}{k},\sigma }(0)>}.`$ (7) Frequency moments of $`A(\stackrel{}{k},w)`$ correspond to ground state (or thermodynamic) averages of the following operators: $`n_\stackrel{}{k}`$ $``$ $`{\displaystyle \frac{dw}{2\pi }A(k,w)}={\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma }{}}<C_{\stackrel{}{k},\sigma }^{}C_{\stackrel{}{k},\sigma }>`$ (8) $`A_1(\stackrel{}{k})`$ $``$ $`{\displaystyle \frac{dw}{2\pi }\omega A(\stackrel{}{k},w)}`$ (9) $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma }{}}<C_{\stackrel{}{k},\sigma }^{}[\widehat{H},C_{\stackrel{}{k},\sigma }]>`$ (10) $`A_2(\stackrel{}{k})`$ $``$ $`{\displaystyle \frac{dw}{2\pi }\omega ^2A(\stackrel{}{k},w)}`$ (11) $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma }{}}<C_{\stackrel{}{k},\sigma }^{}[\widehat{H},[\widehat{H},C_{\stackrel{}{k},\sigma }]]>`$ (12) A systematic evaluation of these averages, using the perturbative expression for $`\widehat{X}`$ in powers of $`t_{ij}/U`$, can be found in the Appendix. ## II The spectral weight, $`n_\stackrel{}{k}`$ We begin with a discussion of the spectral weight because it has an immediate physical interpretation as the occupation probability. As shown in the Appendix, $$n_\stackrel{}{k}=\frac{1}{2}[1\frac{4\stackrel{~}{ϵ_k}}{U}+O(t^3/U^3)]$$ (13) where $$\stackrel{~}{ϵ_k}=\underset{j}{}S_{0i}t_{0j}e^{i\stackrel{}{k}\stackrel{}{R}_j}$$ (14) and $$S_{ij}=<\frac{1}{4}\stackrel{}{S_i}\stackrel{}{S_j}>$$ (15) is the equilibrium spin correlation between spins i and j. $`\stackrel{~}{ϵ_k}`$ is a sort of renormalized band energy in which each hopping matrix element $`t_{ij}`$ is renormalized by a factor of $`S_{ij}`$. However, this renormalized energy does not correspond in any simple way to the energy of any elementary excitation of the system. Note that the non-interacting free electron band is given by $`ϵ_k=_jt_{0j}e^{i\stackrel{}{k}\stackrel{}{R}_j}`$. This expression for $`n_\stackrel{}{k}`$ can be derived in a different, simpler manner, which is readily generalizible to more complicated situations, such as the three band Cu-O or Emery model. From the Hellman-Feynman theorem, it follows that $$\underset{\sigma }{}<[c_{i\sigma }^{}c_{j\sigma }+\mathrm{H}.\mathrm{C}.]>=E/t_{ij}=[J_{ij}/t_{ij}]S_{ij}$$ (16) where $`E`$ is the internal energy, which can be computed using the effective Hamiltonian $`\widehat{H}_{eff}`$. From the expression for $`J_{ij}`$ in terms of $`t_{ij}`$, the result in Eq.13 follows immediately. ### A The remnant fermi surface At low temperatures, the short-range spin correlation functions are essentially temperature independent, and equal to their value in the ground state. To be concrete, let us consider the Hubbard model with nearest, second, and third neighbor hopping, $`t`$, $`t^{}`$ and $`t^{\prime \prime }`$, respectively; this sort of model was used in the numerical studies to fit the dispersions seen in ARPES. The zero temperature spin correlations of the corresponding spin 1/2 Heisenberg model have been computed fairly accurately in numerical studies. For the Heisenberg model with only nearest-neighbor exchange coupling, the spin correlation functions are $`S_{01}7/12`$ and $`S_{03}S_{02}\frac{1}{20}`$ between nearest, next-nearest, and third nearest neighbor sites, respectively. These correlations are, moreover, found to be relatively insenstitive to the inclusion of a modest amount of further neighbor exchange couplings, which anyway are expected to be quite small since $`J^{}/J=[t^{}/t]^2`$. In computing $`\stackrel{~}{ϵ}_k`$ the antiferromagnetic correlations between neighboring spins imply a factor of 1/2 renormalization of $`t`$, compared to a factor of 1/20 renormalization of $`t^{}`$ and $`t^{\prime \prime }`$. Since in most cases of physical interest, $`\|t^{}\|,\|t^{\prime \prime }\|10\|t\|`$, even when $`t^{}`$ and $`t^{\prime \prime }`$ are large enough to make signicant shifts in the original Fermi surface defined by $`ϵ_k`$, the occupation probability is well approximated (Figure 1) as $$n_\stackrel{}{k}\frac{1}{2}[1+\frac{7t(cos(k_x)+cos(k_y))}{3U}]$$ (17) It is both $`\frac{1}{2}`$ and has the steepest slope along the Fermi surface of the non-interacting electrons with only the nearest neighbor hopping. In fact, ARPES experiments that observe such momentum dependence of $`n_\stackrel{}{k}`$ have been interpreted as an indication that there is a “remnant Fermi surface” in the undoped Mott insulator. By contrast, our result suggests that the observed $`n_\stackrel{}{k}`$ is reflective of the the spin physics of the strongly correlated Neel state rather than a vestige of the original Fermi surface. ### B Specific heat of the Neel transition The natural connection we find between the spectral weight and the Neel state can be exploited further. Since, as we already pointed out, the Heisenberg Hamiltonian is dominated by the nearest neighbor term, and assuming that $`J_{n.n.}`$ is only weakly (if at all) temperature dependent, the specific heat is $$C(T)=J\frac{}{T}<\frac{1}{4}\stackrel{}{S_1}\stackrel{}{S_0}>=J\frac{S_1}{T}.$$ (18) Since $$n_\stackrel{}{k}\frac{1}{2}[1+\frac{8t(cos(k_x)+cos(k_y))}{U}S_1],$$ (19) by measuring the temperature dependence of $`S_1`$ (as extracted from $`n_\stackrel{}{k}`$) and differentiating it one gets the magnetic contribution to the specific heat. The specific heat extracted in this manner contains only the contribution from the spin fluctuations. ## III Higher Moments Higher spectral moments correspond to the derivatives $`<C_{\stackrel{}{k},}^{}(t)C_{\stackrel{}{k},}(0)>`$ at $`t=0`$ and thus provide further insight into the problem of a hole in a Mott insulator. Our results for first and second moments, $`A_1(\stackrel{}{k})`$ and $`A_2(\stackrel{}{k})`$, are presented in a table below. The moments are written as sums over momentum space Fourier harmonics, each corresponding to a summation over sites a given near neighbor distance away on the square lattice. $`A_n(\stackrel{}{k})`$ $`=`$ $`{\displaystyle \underset{j}{}}A_n^j\gamma _j(\stackrel{}{k}),`$ (20) $`\gamma _j(\stackrel{}{k})`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{R}=j^{th}n.n.}{}}e^{i\stackrel{}{k}\stackrel{}{R_j}}`$ (21) \| i \| $`A_1^j`$ \| $`A_2^j`$ \| \| --- \| --- \| --- \| \| 1 \| $`t(\frac{1}{2}S_1)`$ \| $`(2tt^{}+tt^{\prime \prime })(S_11)`$ \| \| \| $`\frac{2t^{}t^{\prime \prime }}{U}(S_2+2S_1)`$ \| \| \| \| $`\frac{tt^{\prime \prime }}{U}(S_3+2S_1)`$ \| \| \| 2 \| $`t^{}(\frac{1}{2}S_2)`$ \| $`(2t^{}t^{\prime \prime }+t^2)(S_21)`$ \| \| \| $`\frac{2t^{}t^{\prime \prime }}{U}(S_3+2S_2)`$ \| \| \| \| $`\frac{t^2}{U}(S_2+2S_1)`$ \| \| \| 3 \| $`t^{\prime \prime }(\frac{1}{2}S_1)`$ \| $`((t^{})^2+\frac{t^2}{2})(S_31)`$ \| \| \| $`\frac{(t^{})^2}{U}(S_3+2S_2)`$ \| \| \| \| $`\frac{t^2}{2U}(S_3+2S_1)`$ \| \| \| 4 \| $`\frac{tt^{}}{U}(S_4+S_2+S_1)`$ \| $`(t^{}t^{\prime \prime }+tt^{\prime \prime })(S_41)`$ \| \| \| $`\frac{tt^{\prime \prime }}{U}(S_4+S_3+S_1)`$ \| \| \| 5 \| $`\frac{(t^{})^2}{U}(S_5+2S_2)`$ \| $`(\frac{(t^{})^2}{2}+(t^{\prime \prime })^2)(S_51)`$ \| \| \| $`\frac{(t^{\prime \prime })^2}{U}(S_5+2S_3)`$ \| \| \| 6 \| $`\frac{tt^{\prime \prime }}{U}(S_6+S_3+S_1)`$ \| $`tt^{\prime \prime }(S_61)`$ \| \| 7 \| $`\frac{t^{}t^{\prime \prime }}{U}(S_7+S_3+S_2)`$ \| $`t^{}t^{\prime \prime }(S_71)`$ \| \| 10 \| $`\frac{(t^{\prime \prime })^2}{2U}(S_{10}+2S_3)`$ \| $`\frac{(t^{\prime \prime })^2}{2}(S_{10}1)`$ \| One notices that all computed moments are finite. Although the existance of the moment expansion is a general requirement of any physical system, it is more rule than exception that approximations lead to divergencies past some finite order. Though we haven’t constructed an explicit proof, there are indications that the expansion is well behaved for the Hubbard model. In principle, the moment expansion can be used to study the QP dispersion ($`E_\stackrel{}{k}`$) directly: a coherent oscillation results in $`E_\stackrel{}{k}^n`$ contribution to $`A_n(\stackrel{}{k})`$. The crudest (single mode) approximation of this sort identifies $`\overline{E}_\stackrel{}{k}=A_1(\stackrel{}{k})/n_\stackrel{}{k}`$. We found it to be in a surprising agreement with previously obtained results for the momentum dependence of the QP energy in t-t’-t”-J as well as t-J models (one needs to assume that term proportional to t does not contribute to QP dispersion). Since $`\overline{E}_\stackrel{}{k}`$ rather seriously overestimates the overall bandwidth, it isn’t clear if one is justified claiming to have obtained even an approximate QP energy yet. Another likely use of our results (or rather their extension to higher moments and orders in perturbation theory) can be in comparing with spectral functions obtained by other means (either numerics, self-consistent Born approximation or other). In conclusion, we have outlined a well controlled method for analysing the spectral moments of a hole in a Mott insulator. We find the occupation probability, $`n_\stackrel{}{k}`$, is in agreement with the well established experimental result, which as our calculation suggests is strongly constrained by the presence of AF order. We further propose that the temperature dependence of $`n_\stackrel{}{k}`$ can be used to study the specific heat of the Neel transition. The implications of our results for higher moments are yet to be understood properly. I would like to thank Steven Kivelson for innumerable illuminating conversations and suggesting this problem in the first place. I would also like to greatfully acknowledge M.Z.Hasan, A. Lanzara, F.Ronning and Z.X.Shen for useful discussions and the hospitality of the Physics and Applied Physics Departments at Stanford University where this work was carried out. This work was supported by NSF grant number DMR98-08685. ## A Evaluation of moments To compute the moments we first express them in terms of real space electronic correlations $`n_k`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma }{}}<C_{\stackrel{}{k},\sigma }^{}C_{\stackrel{}{k},\sigma }>`$ (A1) $`=`$ $`{\displaystyle \frac{1}{2}}\{1+{\displaystyle \underset{i,\sigma }{}}<C_{\stackrel{}{i},\sigma }^{}C_{\stackrel{}{0},\sigma }>\gamma _i(k)\}`$ (A2) $`A_1(\stackrel{}{k})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma }{}}<C_{\stackrel{}{k},\sigma }^{}[\widehat{H},C_{\stackrel{}{k},\sigma }]>`$ (A3) $`=`$ $`ϵ_kn_k+{\displaystyle \frac{U}{2}}{\displaystyle \underset{i,\sigma }{}}\gamma _i(k)<C_{\stackrel{}{i},\sigma }^{}C_{\stackrel{}{0},\sigma }n_{0,\sigma }>`$ (A4) $`A_2(\stackrel{}{k})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma }{}}<C_{\stackrel{}{k},\sigma }^{}[\widehat{H},[\widehat{H},C_{\stackrel{}{k},\sigma }]]>`$ (A5) $`=`$ $`ϵ_kn_\stackrel{}{k}\overline{E}_\stackrel{}{k}+{\displaystyle \frac{U^2}{2}}{\displaystyle \underset{i,\sigma }{}}\gamma _i(k)<C_{\stackrel{}{i},\sigma }^{}C_{\stackrel{}{0},\sigma }n_{0,\sigma }>`$ (A6) $`+`$ $`{\displaystyle \frac{U}{2}}{\displaystyle \underset{i,j,\sigma }{}}\gamma _i(k)t_{j,0}<C_{\stackrel{}{i},\sigma }^{}(C_{\stackrel{}{j},\sigma }n_{0,\sigma }`$ (A7) $`+`$ $`C_{\stackrel{}{j},\sigma }C_{\stackrel{}{0},\sigma }^{}C_{\stackrel{}{0},\sigma }+C_{\stackrel{}{j},\sigma }^{}C_{\stackrel{}{0},\sigma }C_{\stackrel{}{0},\sigma }>`$ (A8) As before $`\gamma _i(\stackrel{}{k})=_{R_i}e^{i\stackrel{}{k}\stackrel{}{R_i}}`$,where the sum is over i’th nearest neighbors. Next, these electronic correlations are evaluated perturbatively ($`S_{ij}=<\frac{1}{4}\stackrel{}{S}_i\stackrel{}{S}_j>`$): $`{\displaystyle \underset{\sigma }{}}<C_{\stackrel{}{i},\sigma }^{}C_{\stackrel{}{0},\sigma }>`$ $`=`$ $`{\displaystyle \frac{4t_{i0}}{U}}S_{i0}+O(t^3/U^3)`$ (A9) $`{\displaystyle \underset{\sigma }{}}<C_{\stackrel{}{i},\sigma }^{}C_{\stackrel{}{0},\sigma }n_{0,\sigma }>`$ $`=`$ $`{\displaystyle \frac{2t_{i0}}{U}}S_{i0}`$ (A10) $`+{\displaystyle \underset{j,R_j}{}}{\displaystyle \frac{t_{iR_j}t_{R_j0}}{U^2}}\{3S_{iR_j}S_{R_j0}`$ $``$ $`S_{i0}+3i<\stackrel{}{S}_i(\stackrel{}{S}_{R_j}\times \stackrel{}{S}_0)>\}`$ (A11) On a square lattice and in a state that doesn’t break time reflection invariance ($`<\stackrel{}{S}_i(\stackrel{}{S}_{R_j}\times \stackrel{}{S}_0)>=0`$) the second sum is (for different i): \| 1 \| $`4tt^{}S_{20}+2tt^{\prime \prime }S_{30}`$ \| \| --- \| --- \| \| 2 \| $`2(2S_1S_2)t^2+4S_3t^{}t^{\prime \prime }`$ \| \| 3 \| $`(2S_1S_3)t^2+2(2S_2S_3)t^2`$ \| \| 4 \| 2$`(S_1+S_2S_4)tt^{}+2(S_1+S_3S_4)tt^{\prime \prime }`$ \| \| 5 \| $`(2S_2S_5)t^2+2(2S_3S_5)t^{\prime \prime 2}`$ \| \| 6 \| 2$`(S_1+S_3S_6)tt^{\prime \prime }`$ \| \| 7 \| 2$`(S_2+S_3S_7)t^{}t^{\prime \prime }`$ \| \| 10 \| $`(2S_3S_{10})t^{\prime \prime 2}`$ \| Substituting these terms into A1,A2,A3 yields results of Section III.
warning/0003/astro-ph0003321.html	ar5iv	text	# 1 Introduction ## 1 Introduction Gravitational collapse of a stellar core to a black hole has been studied since many years. Efforts have been also devoted to the calculation of the gravitational radiation emitted in this process. Most studies have been based on a perturbative approach,<sup>?</sup><sup>,</sup><sup>?</sup><sup>,</sup><sup>?</sup><sup>,</sup><sup>?</sup><sup>,</sup><sup>?</sup> others on the numerical solution of the full Einstein equations <sup>?</sup><sup>,</sup><sup>?</sup>. All these papers consider the collapse of a “naked” stellar core, described by a “dust” of particles or, at most, by a politropic equation of state, without taking into account the presence of the outer layers of the star, which are involved in the process. In particular, depending on the ratio between the energy released in the final explosion of a massive star and the binding energy of the ejected material, two different kinds of collapse have been outlined <sup>?</sup><sup>,</sup><sup>?</sup>: the prompt collapse, in which a large fraction of the star collapses on a dynamical time-scale forming a massive black hole, and the delayed collapse, in which a low mass black hole or, alternatively, a neutron star forms at the beginning and later accretes matter, due to fall-back. If a neutron star is the initial outcome of the collapse, fall-back pushes its mass above the critical mass and the formation of a black hole takes place. In both cases, this light black hole continues to slowly accrete matter until the final mass is reached. In the gravitational wave community the prompt collapse has been considered for a long time as representative of realistic collapse processes. The estimated mass (several solar masses) of the first black hole candidates (like $`CygnusX1`$) has led to the idea that black holes should often be born with a ”typical” mass of $`10M_{}`$. This assumption appears no more justified now: we know that the evolution of massive stars in binary systems (to which all observed black hole candidates belong) is different from that of single stars; in addition, more refined observational techniques have allowed to find black hole candidates of few solar masses. In the light of these results and of the recent numerical simulations which we will describe, the prompt birth of such massive black holes ($`10M_{}`$ or more) should be considered as a very rare event. In this paper, we will discuss the delayed collapse model, the expected mass distribution of isolated black holes and the detectability of the emitted signal, by forthcoming interferometric detectors. Our main aim is to understand how our perspectives of detection change with respect to the ”naive” prompt collapse. The plan of the paper is as follows. In Sec.2 we will shortly describe the delayed collapse scenario. In Sec.3 we will estimate the contribution of fall-back to the total gravitational emission in the delayed collapse discussing the consequences. In Sec.4 we will discuss the detectability of the emitted signal by forthcoming interferometric detectors, and compare with the prompt collapse which could happen for progenitor stars above about $`40M_{}`$, if stellar winds were much less important in the evolution of massive stars than it is currently believed. In Sec.5 we will derive the theoretical final mass distribution function for isolated black holes, using the results of recent simulations of the collapse of massive single stars. Finally, in Sec.6 the results and their implications will be discussed. ## 2 Star Collapse to a Black Hole The fate of a massive star ($`m_{prog}>9M_{}`$) is the core collapse with the formation of a compact object: a neutron star or a black hole. Numerical simulations show that the actual final product, and also the way in which it is formed, depend on many factors, among which the mass and the angular momentum of the progenitor, the explosion energy, the high density matter equation of state ($`EOS`$) and also the way in which the physics is implemented (regarding, for instance, neutrino physics or angular momentum transport). Moreover, results cannot be considered conclusive until fully relativistic 3D simulations will be performed. It has been shown that, for a wide range of progenitor masses and explosion energies, the shock cannot expell all the matter outside the collapsing core, so that part of the helium mantle and heavy elements may slow down below the escape velocity and be accreted by the just formed neutron star, with a timescale of minutes to hours. If the neutron star mass grows above a critical value, a black hole forms and we have a delayed collapse. On the other hand, if the explosion completely fails, or if it is too weak, a black hole immediately forms: this is the prompt collapse. According to recent simulations (starting from non-rotating progenitors) by Woosley $`\&`$ Weaver <sup>?</sup> and by Fryer $`\&`$ Kalogera <sup>?</sup>, typical collapses are always delayed, the prompt ones occurring only for very massive progenitor stars ($`M>40M_{}`$) if stellar winds are negligible, an assumption which appears rather unlikely. The $`EOS`$ of high density matter plays a basic role in the determination of neutron stars limiting mass, i.e. the critical mass for black hole formation, $`m_{min}`$. For conventional EOS the neutron star maximum mass ranges between $`1.7M_{}`$ and $`2.2M_{}`$ (e.g. <sup>?</sup><sup>,</sup> <sup>?</sup><sup>,</sup><sup>?</sup><sup>,</sup><sup>?</sup>), with a $`10÷20\%`$ increase if rotation is taken into account. On the other hand, if $`\pi `$ or $`K`$ condensation, or formation of quark matter, takes place at very high densities (e.g. <sup>?</sup><sup>,</sup><sup>?</sup><sup>,</sup><sup>?</sup><sup>,</sup><sup>?</sup>), the $`EOS`$ is softened and this reduces the maximum mass that the pressure of degenerate matter can sustain. In this case, the critical neutron star mass is $`1.5M_{}`$. The masses of 26 neutron stars, measured in pulsars, are compatible with a gaussian distribution with $`\overline{m}=1.35\pm 0.04`$, and then are consistent with that previous limit <sup>?</sup>. Moreover, the lack of evidence for neutron stars with mass near the maximum derived from conventional $`EOS`$ is reinforced by the lack of evidence for a pulsar as a remnant of the supernova $`SN1987A`$ whose progenitor had a mass of $`18M_{}`$, which should have left behind a remnant mass of about $`1.5M_{}`$ <sup>?</sup><sup>,</sup><sup>?</sup>. However, no definite conclusion can be still drawn<sup>?</sup><sup>,</sup><sup>?</sup>. Given the uncertainties in the neutron stars equation of state, in the following we will always refer to two different values of the black hole minimum mass: $`m_{min}=1.5M_{}`$, representative of soft equations of state (soft EOS), and $`m_{min}=2M_{}`$ for standard, with no phase transition, equations of state (conventional EOS). ## 3 Gravitational Radiation from Fall-Back In this section, we want to estimate the contribution of fall-back to the emission of gravitational waves in the delayed collapse of a massive star to a black hole. There are many studies on the capture by a black hole of finite-size shells of matter <sup>?</sup><sup>,</sup><sup>?</sup><sup>,</sup><sup>?</sup>. A general result is that the amount of gravitational radiation emitted is always smaller than that emitted by the capture of a pointlike particle with the same mass of the shell. This is a consequence of destructive interference of the radiation emitted by different parts of the infalling extended matter. These results are confirmed also by the recent, more realistic, calculations by Papadopoulos $`\&`$ Font <sup>?</sup>. In the case of an axisymmetric irrotational shell of matter with mass $`\mu `$, much smaller than the mass $`m`$ of the black hole, they find that the efficiency in the emission of gravitational waves decreases as a function of the radial width $`L`$ of the shell as follows: $$ϵ_s=\frac{\mathrm{\Delta }E}{mc^2}=810^3\left(\frac{\mu }{m}\right)^2\left(\frac{m}{M_{}}\right)^{2.4}\left(\frac{L}{1km}\right)^{2.4}$$ (1) We shall now extrapolate this equation for $`\mu >m`$ and compare to the efficiency we expect from core collapse which, for a maximally rotating core, is $`ϵ_c710^4`$ in the axisymmeric case <sup>?</sup>. For instance, assuming that the initial mass of the newly formed black hole is $`m=2M_{}`$, and that the mass of the shell is $`\mu =10M_{}`$, we find that the condition $`ϵ_s<0.1ϵ_c`$ requires $`L>100km`$. This condition is largely verified for massive stars. Then, we can conclude that the main burst of gravitational radiation is emitted at the formation of the black hole while subsequent accretion of matter gives no important contributions. To this respect, two possibilities can occur. First, the $`Fe`$ core of the star has a mass greater than $`m_{min}`$. In such a case the black hole initial mass is equal to the core mass. Second, the $`Fe`$ core is lighter than $`m_{min}`$. In this case, a neutron star is produced at the beginning. It then accretes matter until the critical mass is reached, so that the black hole initial mass is equal to the minimum one. From the core masses given by Woosley $`\&`$ Weaver <sup>?</sup> we see that for soft EOS ($`m_{min}=1.5M_{}`$) all stars with initial mass greater than $`18M_{}`$ produce cores that immediately collapse to a black hole of mass greater than $`m_{min}`$, but lower than $`1.8M_{}`$. The same happens according to Fryer <sup>?</sup>. For conventional EOS ($`m_{min}=2M_{}`$) black holes are formed, after fall-back on a neutron star, from the collapse of progenitors of mass $`m_{prog}>26M_{}`$ (following Fryer $`\&`$ Kalogera <sup>?</sup> this value is reduced to $`m_{prog}20M_{}`$). Then, all black holes form with an initial mass equal, or nearly equal, to the minimum one, presumably in the range $`1.5÷2.6M_{}`$$`^\text{e}`$footnotetext: $`^\text{e}`$The upper limit takes into account also the possible stabilizing effect of rotation., depending on the nuclear density matter equation of state and on the rotation rate. In the following we will assume, for simplicity of calculation, that all black holes form with the same mass $`m_{min}`$, so that their initial mass distribution function, for the delayed collapse scenario, can be written as a $`\delta `$-function: $`f(m)=\delta (mm_{min})`$. We will consider $`m_{min}=1.5M_{}`$ and $`m_{min}=2M_{}`$. ### 3.1 The rate of black hole formation Let us now estimate the expected event rate for collapses. We assume a Galactic rate for Supernovae of type $`II`$ given by $`R_{SNII}=0.02yr^1`$. Such rate is the sum of the rate of collapses to neutron stars and, in the delayed scenario, of the rate of collapses to black hole (a supernova explosion is produced by all delayed collapses, indipendently of $`m_{min}`$). For instance, for soft EOS, we have seen that a black hole is produced by progenitor stars with mass $`18M_{}m_{prog}40M_{}`$. Their formation rate, $`R_{bh}`$, is a fraction $`\lambda `$ of the rate of formation of neutron stars $`R_{ns}`$, that are generated by progenitors with mass $`9M_{}m_{prog}18M_{}`$. Of course, this ratio depends on the initial mass function of the progenitor stars: $`\lambda =0.43`$ for a Salpeter law with exponent $`\alpha =2.35`$, while $`\lambda =0.32`$ for a Scalo law with $`\alpha =2.7`$. Then, we can write $$R_{SNII}=R_{ns}+R_{bh}=R_{ns}(1+\lambda ),$$ (2) from which we can get $`R_{ns}`$ and $`R_{bh}`$. In Tab.(1) we tabulate the galactic rates of black hole formation in the different cases, considering both delayed ($`R_{bh}`$) and prompt ($`R_{bh,p}`$) collapses. Moreover, the total expected rate within the Virgo cluster, $`R_{\mathrm{@}20Mpc}`$, is given in the last column. From Tab.(1) we see that the fraction of delayed collapses producing a black hole is in the range $`10\%÷43\%`$ of those leading to a neutron star. These percentages increase to $`18\%÷60\%`$ if prompt collapse, for stars more massive than $`40M_{}`$, is also taken into account. ## 4 The Detectability Gravitational signals emitted in star collapse to a black hole are characterized by an initial part, emitted during the in-fall phase and bounce, followed by an oscillating tail, which can be described as a superposition of damped sinusoids, corresponding to the black hole quasi-normal modes. The frequency and the damping time of the quasi-normal modes are a function of the black hole parameters, mass and angular momentum. For axisymmetric collapses, the main contribution (typically $`90\%`$ of the whole energy emitted) is given by the $`l=2`$ mode, for which the frequencies and the damping times are given, as a function of the rotation parameter $`a=J/\left(\frac{Gm^2}{c}\right)`$, in Tab.(2). We use the waveforms calculated by Stark $`\&`$ Piran <sup>?</sup>, who have computed, in the framework of the full non-linear theory, the gravitational signals emitted in the axisymmetric collapse to a black hole of a rotating core, for different values of the angular momentum. The maximum efficiency reached in their simulations is $`ϵ_{max}710^4`$. From the waveforms we have computed the corresponding one-sided power spectrum $`f(\nu )`$, i.e. the flux of energy per unit frequency which is given by $$f(\nu )=\frac{\pi c^3\nu ^2}{2G}<h^2(\nu )>$$ (3) where $`<h^2>`$ denotes the average of the squared gravitational signal with respect to its angular dependence. These computed energy spectra allow to evaluate the signal to noise ratio ($`SNR`$), according to the well-known formula given by Eq.(4) below, which holds if the matched filter is applied to the data. This optimum filtering procedures implies that we know the exact waveform emitted and this is a rather optimistic hypothesis even in the case of star collapses to a black hole. We know that in the simplest cases (low rotation, no hydrodynamical effects) most of the energy is emitted in the phase of quasi-normal ringing, which can be described as a superposition of damped sinusoids. On the other hand, if the collapse has a high degree of rotation, or if strong hydrodynamical effects take place, the emitted waveform may be no longer dominated by the quasi-normal ringing and would be less predictable <sup>?</sup>. In this case, other, less optimum, data analysis procedures should be used <sup>?</sup>. Then, our results have to be considered as upper limits on the $`SNR`$. The average squared $`SNR`$ can be expressed as $$\overline{SNR^2}=\frac{8G}{5\pi c^3}_0^{\mathrm{}}\frac{f(\nu )}{\nu ^2S_h(\nu )}𝑑\nu .$$ (4) where $`S_h(\nu )`$ is the detector noise power spectrum. In Eq.(4) an average on the source-detector relative position and on the polarization of the gravitational waves emitted is also performed. As discussed in Sec.3, the delayed collapse produces black holes that, immediately after birth, have mass very close to the minimum mass. Thus, the emitted signal depends essentially on the black hole angular momentum and on the distance at which the collapse takes place. In prompt collapses, on the contrary, the $`SNR`$ clearly depends on the mass of the black hole and, due to the typical sensitivity curve of ground-based interferometers, the radiation emitted by more massive black holes is more easily detectable, so that the delayed collapse process is definitely less favourable than the “naive” prompt one. In Tab.(3) we give the $`SNR=\sqrt{\overline{SNR^2}}`$ for delayed collapses taking place in the Galaxy (we fix $`r=10kpc`$), for different values of the minimum black hole mass and of the rotation parameter. The values of the $`SNR`$ increase with the angular momentum of the black hole. This is the consequence of two effects: first, the efficiency of emission increases as $`a^4`$; second, for very high angular momentum (say, $`a>0.8`$) the bounce of the collapsing star produces a lower frequency component in the signal energy spectrum which fits better to the sensitivity curve of interferometers. On the other hand, the expected rate of detectable events is low: collapses to a black hole are detectable, essentially, only within the Local Group, with a total rate, strongly dominated by the Milky Way, which is less than $`1`$ event per century, see Tab.(1). We stress, however, that higher $`SNR`$ could be obtained if collapses had a higher degree of asymmetry and there are some observative indications supporting this hypothesis <sup>?</sup><sup>,</sup><sup>?</sup><sup>,</sup><sup>?</sup><sup>,</sup><sup>?</sup>. Moreover, the initial core collapse to a neutron star which takes place, in the delayed case, if $`m_{min}=2M_{}`$, could be a promising process. If it is highly asymmetric, as observations seem to indicate, a large amount of gravitational radiation could be emitted in the range of frequencies where interferometric detectors reach their best sensitivity$`^\text{f}`$footnotetext: $`^\text{f}`$It should be noted, however, that according to the newtonian simulations by Rampp et al. <sup>?</sup>, the amount of gravitational radiation emitted in non-axisymmetric collapses is comparable to that of the axisymmetric case.. We have repeated the calculation of the $`SNR`$ for the advanced $`LIGO`$ detector, using an approximation to its sensitivity curve given by Flanagan $`\&`$ Huges <sup>?</sup>. Results are given, for delayed collapses, in Tab.(4) where a distance of $`r=20Mpc`$ has been assumed. Within this distance, corresponding approximately to the Virgo Cluster, the expected black hole formation rate is $`1÷5yr^1`$, but we see that the detection perspectives are not much better because the $`SNR`$ is much smaller than one. The situation is different for prompt collapses. In Fig.(4) we have plotted the signal-to-noise ratio for prompt collapses as a function of the newborn black hole mass and it appears that the emitted gravitational radiation could be detectable if the rotation rate is high enough. ## 5 Black hole final mass distribution function In this section we calculate the black hole final mass distribution function. We do not consider the chemical evolution of the galaxy, which affects the mass distribution of compact remnants (see, e.g., Timmes et al. <sup>?</sup>). Such a theoretical distribution could be useful in statistical studies of recently formed black holes. The black hole final mass distribution depends on the mass distribution of progenitor stars and on the relation between the mass of progenitors and that of the final remnants produced after the collapse. For the first one we assume a power law $`f(m_{prog})m_{prog}^\alpha `$, with $`\alpha =2.35`$ (Salpeter’s law) and $`\alpha =2.7`$ (Scalo’s law). Regarding the relation between the mass of progenitor stars, $`m_{prog}`$, and that of remnants, $`m`$, we use (after conversion from baryonic to gravitational mass) the results of Woosley $`\&`$ Weaver <sup>?</sup> and of Fryer $`\&`$ Kalogera <sup>?</sup>, who have performed systematic calculations of the evolution and explosion of non rotating massive stars. It must be stressed again that such evolutionary calculations are still subject to many uncertainties and this reflects in some differences between their results, obtained using different assumptions. The remnant masses they report are calculated a long time after the collapse, so that the possible fall-back is included. The distribution function we are searching for is simply given by $$f(m)=f(m_{prog}(m))\frac{dm_{prog}}{dm}$$ (5) In Fig.(5) we plot the function $`f(m)`$ obtained using a fit of the progenitor mass-remnant mass relation found by Woosley $`\&`$ Weaver <sup>?</sup> (see their table 3), considering different values of $`m_{min}`$ and of the exponent $`\alpha `$. The maximum progenitor mass considered by Woosley $`\&`$ Weaver is $`m_{prog}=40M_{}`$. According to them, for masses greater than this, the effect of stellar wind becomes relevant and strongly modifies the star evolutionary path: a smooth convergence of remnant masses to $`1.5M_{}`$, corresponding to neutron stars or low mass black holes, is expected for progenitors of mass greater than about $`40M_{}`$. As a consequence, black holes with mass greater than $`10.3M_{}`$ \- which is what is predicted by Woosley $`\&`$ Weaver for a $`40M_{}`$ progenitor star if stellar winds are neglected - should never be formed, for progenitors of solar metallicity.$`^\text{a}`$footnotetext: $`^\text{a}`$Black holes of mass greater than that value could be produced if the metallicity is lower. Also Fryer $`\&`$ Kalogera <sup>?</sup> show that low mass remnants are produced, for progenitors more massive than about $`40M_{}`$, if stellar winds are taken into account. On the other hand, if stellar winds are neglected, they find that, for masses above about $`40M_{}`$, the star collapse takes place in the prompt way, i.e. a black hole of mass nearly equal to the mass of the progenitor star is produced in a dynamical timescale. In this case the black hole mass distribution function is simply proportional to $`m^\alpha `$. The final mass distribution function, resulting from the progenitor mass-remnant mass relation found by Fryer $`\&`$ Kalogera, is plotted in Fig.(5) (we refer to their “most likely” model, in which $`50\%`$ of the explosion energy goes in unbinding the star). There are clear differences between the distributions functions plotted in Figs.(5,5). First of all, the expected range of masses for black holes is larger in the second case, reaching $`14.7M_{}`$. Moreover, in this case the function $`f(m)`$ is monotonically decreasing, while in Fig.(5) there is a minimum around $`3M_{}`$ followed by a relative maximum at $`4.5M_{}`$. These differences are a consequence, obviously, of the different relations for progenitor vs. remnant masses$`^\text{b}`$footnotetext: $`^\text{b}`$In particular, the minimum in Fig.(5) is due to an increase of the slope of $`m_{prog}(m)`$ for $`m_{prog}30M_{}`$ while the following maximum is produced by the first term in the right-hand side of Eq.(5), which is dominant for large $`m_{prog}`$. and are due, mainly, to two reasons. First, explosion simulations are carried on by Woosley $`\&`$ Weaver in one-dimension, while Fryer $`\&`$ Kalogera simulations are two-dimensional. Second, Woosley $`\&`$ Weaver constrain the kinetic energy of ejected material at infinity$`^\text{c}`$footnotetext: $`^\text{c}`$Defined as the difference between the explosion energy and the binding energy of the ejecta to a constant value, about $`1.210^{51}erg`$, while in Fryer $`\&`$ Kalogera this quantity decreases for increasing progenitor mass. This explains why their remnant masses are greater. The theoretical distribution functions we have found cannot be compared with observations, which refer to black candidates belonging to binary systems. It is widely accepted that the evolution of high mass stars in binary systems is quite different from that of single stars of comparable mass. A detailed discussion on massive binary systems evolution and on the possible mechanisms and bias effects which can explain the observed black hole candidates mass distribution can be found in refs. <sup>?</sup> and <sup>?</sup>. ## 6 Conclusions In this paper we have considered the collapse of a massive star to a black hole, exploring the models of prompt and delayed collapse. According to recent simulations, the formation of stellar mass black holes should take place through a delayed collapse. The prompt collapse could happen for very massive progenitors (mass greater than $`40M_{}`$) only if stellar winds were negligible, an assumption which appears to be not very reliable. In the delayed collapse all black holes were born with a mass equal to its minimum value, or just a little greater, and then slowly accrete matter up to their final mass. We have shown that the main burst of gravitational radiation is emitted when the black hole forms so that the gravitational energy spectrum is peaked in the range $`4.5÷9kHz`$, depending on the initial mass and the angular momentum of the black hole. Such frequencies do not match very well with the sensitivity curve of ground-based interferometers which reach their best sensitivity in the band $`60÷1000Hz`$. We have estimated the detectability of the emitted gravitational radiation in the delayed case, the most reliable scenario, and in the prompt case, for progenitor stars more massive than $`40M_{}`$. For each collapse model, we have considered different values of black hole minimum mass and angular momentum and different laws for progenitors mass distribution. We have also derived the theoretical black hole final mass distribution function. Delayed collapses are detectable only inside the Local Group of galaxies by interferometers of the first generation. Obviously, delayed collapses are less detectable than the prompt ones, due to the lower matching of their characteristic frequencies to the sensitivity curve of the detectors. On the other hand, the initial stage of the delayed collapse, with the formation of a neutron star following a strongly asymmetric explosion (if $`m_{min}=2M_{}`$), could be a promising source of gravitational waves. Detection perspectives of delayed collapses are not so much better for adavnced interferometers because, if distances up to the Virgo Cluster are considered (where we expect $`1÷5ev/yr`$), the signal-to-noise ratio is much lower than one, unless a very high degree of asimmetry is produced. On the contrary, prompt collapses could be detected with large enough $`SNR`$. Delayed collapses to a black hole belong to a class of high frequency sources of gravitational waves, which comprises various processes involving compact objects, as for instance the excitation of neutron star w-modes <sup>?</sup>, the coalescence of two neutron stars with the formation of a light black hole <sup>?</sup>, dynamical instabilities in neutron stars <sup>?</sup>, and some kinds of secular instability in neutron stars <sup>?</sup>. Such sources cannot be efficiently detected by present resonant detectors and forthcoming interferometers because they are expected to emit at frequencies higher than those at which both interferometric and resonant detectors have their best sensitivity. In past years local arrays of small resonant detectors, which are particularly suited for the detection of high frequency gravitational radiation, have been proposed<sup>?</sup><sup>,</sup><sup>?</sup>. A detailed study of their detection performances, considering different geometries, dimensions and materials, has been done <sup>?</sup>. The non-continuous background of gravitational waves produced by the ensemble of the star collapses to a black hole, which occurred at a higher rate in the early phases of the Universe, has recently been calculated <sup>?</sup>. It would be interesting to repeat the calculation in the case of delayed collapse. In such a case we have the superposition of energy spectra nearly with the same shape and all peaked at nearly the same frequency, in the range $`4.5÷9kHz`$ (this holds also for star metallicity $`Z=0`$). Roughly speaking, as most of the collapses take place at redshift $`z2`$, we expect to have, at the detector, a background spectrum strongly peaked somewhere in the band $`1.5÷3kHz`$. Aknowledgements I want to thank S. Frasca and M. A. Papa for the useful discussions and suggestions, T. Piran for his encouragement, V. Ferrari for her careful reading of the manuscript and the anonymous referees for their comments aiming to an improvement of this paper. References
warning/0003/gr-qc0003068.html	ar5iv	text	# On the strength of the Kerr singularity and cosmic censorship ## I Introduction The cosmic censorship hypothesis of Penrose states that a physically realistic gravitational collapse can never result in a “naked” singularity—that is, all singularities arising in such situations should always be enclosed within an event horizon and hence invisible to distant observers. This hypothesis plays a fundamental role in the theory of black holes. Unfortunately, in spite of much sustained effort, no complete proof (or convincing counterexample) to cosmic censorship has been found as yet. Many partial results, however, have been established which considerably restrict the class of possible naked singularities. A cosmic censorship theorem of this type has previously been proved by one of the present authors . The aim of this paper is to answer a certain essential question closely related to that result. There exist many exact solutions of the Einstein equations in which naked singularities do occur. Among these solutions the most important one, due to its direct astrophysical applications, is undoubtedly that of Kerr . This solution depends on two parameters $`m`$ and $`a`$ (see below), and represents the exterior gravitational field of a rotating body with mass $`m`$ and angular momentum $`ma`$, as measured from infinity in geometrized units with $`c=G=1`$. As is well known , the maximal analytic extension of the Kerr solution with $`m>0`$ and $`a0`$ contains the ring curvature singularity that may be interpreted as the outcome of collapse of a rotating object. When $`\|a\|m`$, this singularity is always hidden behind an event horizon. But if $`\|a\|>m`$, there is no event horizon and the singularity is visible for all observers; moreover, there are closed timelike curves through every point of the spacetime . Clearly, the Kerr solution is highly idealized and the singularity with $`\|a\|>m`$ cannot be a counterexample to Penrose’s hypothesis. However, it is not unlikely that the pathologies occurring in the case $`\|a\|>m`$ could also arise in more general scenarios of the collapse of a rapidly rotating star. The censorship theorem of Ref. shows that, under certain reasonable assumptions, a physically realistic collapse developing from a regular initial state cannot lead to the formation of a final state resembling the Kerr solution with $`\|a\|>m`$ —i.e. of a naked singularity accompanied by closed timelike curves. An important role in this result plays a certain inextendibility condition, which is assumed to hold for all (achronal) null geodesics terminating at the Kerr-like naked singularity under consideration. This condition characterizes the curvature strength of the singularity; roughly speaking, it holds for a given null geodesic $`\lambda `$ if the curvature near the singularity is strong enough so that at least one irrotational congruence of Jacobi fields along $`\lambda `$ is forced to refocus as $`\lambda `$ approaches the singularity (for more details on this condition, see Refs. ). It is well known that the inextendibility condition will always hold for achronal null geodesics terminating at the so-called strong curvature singularities defined by Tipler (see below). These singularities have the property that all objects approaching them are crushed to zero volume. One often assumes that all singularities arising in physically realistic collapse should be of the strong curvature type (see, e.g., Refs. \[9–11\]). However, Israel has suggested that the Kerr singularity fails to be strong in the sense of Tipler’s definition, for it tends to cause repulsive effects. Since the inextendibility condition is similar in spirit to Tipler’s definition (in both cases some Jacobi fields are refocused), the question one immediately asks is whether this condition can still be expected to hold for null geodesics terminating at the Kerr singularity—i.e. whether the censorship theorem of Ref. can be applied to singularities of this type. In this paper we will obtain a result that provides some positive answer to this question. Namely, we will show that, in general, all null geodesics reaching the Kerr singularity with $`m>0`$ and $`a0`$ do in fact terminate, contrary to the suggestion of Israel, at Tipler’s strong curvature singularity. ## II Preliminaries To begin with, we need to recall some of the basic facts on the Kerr solution. In Boyer and Lindquist coordinates $`(r,\theta ,\varphi ,t)`$ the Kerr metric is given by (cf. Ref. , p. 161): $$ds^2=\rho ^2\left(\frac{dr^2}{\mathrm{\Delta }}+d\theta ^2\right)+(r^2+a^2)\mathrm{sin}^2\theta d\varphi ^2dt^2+\frac{2mr}{\rho ^2}(a\mathrm{sin}^2\theta d\varphi dt)^2,$$ (1) where $`\rho ^2r^2+a^2\mathrm{cos}^2\theta `$ and $`\mathrm{\Delta }r^22mr+a^2`$. As mentioned earlier, the constant $`m`$ represents the mass of the metric source while $`a`$ is its angular momentum per unit mass. The Kerr spacetime is stationary and axisymmetric, with Killing vector fields $`\xi ^a(/t)^a`$ and $`\omega ^a(/\varphi )^a`$. Moreover this spacetime is Ricci flat and is of Petrov type D. The ring singularity is located in the equatorial plane $`\theta =\pi /2`$ at points where $`r=0`$. As shown by Carter , the only null geodesics which can reach this singularity are those lying strictly in the equatorial plane on the positive $`r`$ side. The equations of motion for these geodesics are (cf. Ref. , p. 328): $$u^r\frac{dr}{ds}=\pm \left[E^2+\frac{2m}{r^3}\left(LaE\right)^2\frac{1}{r^2}\left(L^2a^2E^2\right)\right]^{1/2},$$ (2) $$u^\varphi \frac{d\varphi }{ds}=\frac{1}{\mathrm{\Delta }}\left[\frac{2ma}{r}E+\left(1\frac{2m}{r}\right)L\right],$$ (3) $$u^t\frac{dt}{ds}=\frac{1}{\mathrm{\Delta }}\left[\left(r^2+a^2+\frac{2a^2m}{r}\right)E\frac{2ma}{r}L\right],$$ (4) where $`u^r`$, $`u^\varphi `$ and $`u^t`$ are the coordinate basis components of the tangent vector, $`𝐮`$, to a given null geodesic $`\lambda (s)`$ parametrized by an affine parameter $`s`$. (Note that the corresponding component $`u^\theta d\theta /ds`$ of $`𝐮`$ identically vanishes as $`\lambda (s)`$ lies in the equatorial plane.) The quantities $`E`$ and $`L`$ are the constants of the motion associated with the Killing vectors $`\xi ^a`$ and $`\omega ^a`$; they are defined as follows: $`E\xi _au^a`$ and $`L\omega _au^a`$. Physically, $`E`$ and $`L`$ can be interpreted, respectively, as the energy at infinity and the angular momentum about the symmetry axis, $`\theta =0`$, of a photon moving along $`\lambda (s)`$. It is also worth noting here that if $`L=aE`$, then $`\lambda (s)`$ must belong to one of the two principal null congruences associated with the algebraic type D of the solution (see Ref. , p. 329). The $`\pm `$ signs in Eq. (2) correspond to outgoing and ingoing geodesics, respectively. Let us also recall that, according to Tipler’s definition , an affinely parametrized null geodesic $`\lambda (s)`$ is said to terminate in a strong curvature singularity at affine parameter value $`s_0`$ if the following holds (cf. Ref. , p. 160): Let $`\mu (s)`$ be a 2-form on the normal space to the tangent vector to $`\lambda (s)`$ determined by two linearly independent vorticity-free Jacobi fields $`𝐙_1(s)`$ and $`𝐙_2(s)`$ along $`\lambda (s)`$, i.e. $`\mu (s)𝐙_1𝐙_2`$. For all $`\mu (s)`$ that vanish for at most finitely many $`s`$ in some neighbourhood $`(s_0,s_1]`$ of $`s_0`$, we have $`lim_{ss_0}\mu (s)=0`$. Very useful criteria for determining whether a given null geodesic terminates in a strong curvature singularity have been found by Clarke and Królak . One of these criteria, which will be used in proving our result, can be formulated as the following Proposition 2.1. Let $`\lambda (s)`$ $`(0<ss_1)`$ be an affinely parametrized null geodesic and let $`\{𝐄_i\}`$ $`(i=1,2,3,4)`$ be a pseudo-orthonormal tetrad parallely propagated along $`\lambda (s)`$, with $`𝐄_1𝐄_1=𝐄_2𝐄_2=𝐄_3𝐄_4=𝐄_4𝐄_3=1`$, all other scalar products vanishing and $`𝐄_4=𝐮`$, where $`𝐮`$ is the tangent vector to $`\lambda (s)`$. Let $`C^m_{4n4}`$ $`(m,n\{1,2\})`$ be some component of the Weyl tensor with respect to the tetrad $`\{𝐄_i\}`$. If there exists some affine parameter value $`s_0(0,s_1]`$ such that $`C^m{}_{4n4}{}^{}Ks^2`$ on $`(0,s_0]`$, where $`K`$ is some positive constant, then $`\lambda (s)`$ terminates in Tipler’s strong curvature singularity as $`s0`$. Proof. The proof follows immediately from Proposition 7 of Ref. . ## III The main result We are now in a position to state and prove our main result. Theorem 3.1. Let $`\lambda (s)`$ $`(s>0)`$ be an affinely parametrized null geodesic in the Kerr spacetime with $`m>0`$ and $`a0`$. Suppose that $`\lambda (s)`$ approaches the ring singularity as $`s0`$. If $`\lambda (s)`$ does not belong to either of the two principal null congruences, then $`\lambda (s)`$ terminates in Tipler’s strong curvature singularity as $`s0`$. In very brief outline, the proof of our result runs as follows. We first construct a certain pseudo-orthonormal tetrad $`\{𝐄_i\}`$ parallely propagated along the geodesic $`\lambda (s)`$. We next find one of the components of the Weyl tensor with respect to the tetrad $`\{𝐄_i\}`$. It turns out that if $`\lambda (s)`$ is not a member of the principal null congruences, then this component must grow along $`\lambda (s)`$ as fast as $`r^5`$, where $`r`$ is the Boyer-Lindquist radial coordinate on $`\lambda (s)`$. Using Eq. (2), we then show that this tetrad component of the Weyl tensor must diverge along $`\lambda (s)`$ at least as fast as $`s^2`$, where $`s`$ is the affine parameter. This implies, by Proposition 2.1, that $`\lambda (s)`$ must then terminate in Tipler’s strong curvature singularity. The rigorous proof is as follows. Proof. Let us assume that the spacetime is parametrized by the Boyer-Lindquist coordinates $`(r,\theta ,\varphi ,t)`$. Since $`\lambda (s)`$ reaches the ring singularity at $`r=0`$, it must lie entirely in the equatorial plane $`\theta =\pi /2`$ on the positive $`r`$ side. The coordinate $`r`$ will change along $`\lambda (s)`$ according to Eq. (2). From this equation it follows easily that there may exist at most two positive values of $`r`$ on $`\lambda (s)`$ for which $`dr/ds=0`$, and so we will have $`dr/ds0`$ along $`\lambda (s)`$ for all $`r`$ in a sufficiently small interval about 0. It is thus clear that there must exist some affine parameter value $`s_1>0`$ such that $`dr/ds>0`$ along $`\lambda (s)`$ for all $`s(0,s_1]`$, because $`r>0`$ and $`s>0`$ on $`\lambda (s)`$, and $`r0`$ along $`\lambda (s)`$ as $`s0`$. Let us now fix such a parameter value $`s_1`$. According to Proposition 2.1, in order to prove the theorem, it suffices to show that the rate of growth of the curvature along $`\lambda (s)`$ is strong enough for $`s`$ in some small interval about 0. For convenience, and without loss of generality, we can thus assume that the affine parameter on $`\lambda (s)`$ ranges over the interval $`(0,s_1]`$. Let $`𝐮`$ be the tangent vector to $`\lambda (s)`$. The coordinate basis components $`u^r`$, $`u^\varphi `$ and $`u^t`$ of $`𝐮`$ are given by Eqs. (2)-(4), with the $`+`$ sign in Eq. (2) as $`dr/ds>0`$ along $`\lambda (s)`$. The corresponding component $`u^\theta `$ of $`𝐮`$ identically vanishes as $`\lambda (s)`$ lies entirely in the equatorial plane. Given any point $`p\lambda (s)`$, we define $`𝐤`$ to be the vector in the tangent space $`T_p`$ with the following coordinate basis components: $`k^r=k^\varphi =k^t=0`$ and $`k^\theta =r^1`$, where $`r`$ is the radial coordinate of $`p`$. One can readily verify that, with respect to the scalar product given by the Kerr metric (1), $`𝐤`$ is a unit spacelike vector orthogonal to $`𝐮`$, i.e. $`𝐤𝐤=1`$ and $`𝐤𝐮=0`$. In addition, $`𝐤`$ is parallely transported along $`\lambda (s)`$ (see Appendix A). The pair $`\{𝐤,𝐮\}`$ can easily be extended to the pseudo-orthonormal tetrad $`\{𝐄_i\}`$ mentioned in Proposition 2.1. To see this, let us first fix some point $`q\lambda (s)`$. Let us now take a unit spacelike vector orthogonal to the plane spanned by the vectors $`𝐮`$ and $`𝐤`$ in the tangent space $`T_q`$; we will denote this vector by $`𝐧`$. (Clearly, such a vector can always be found since $`𝐮`$ is null and $`𝐤`$ is spacelike.) Now let $`\{𝐞_i\}`$ $`(i=1,2,3,4)`$ be some orthonormal basis in the space $`T_q`$, with the spacelike vectors $`𝐞_2`$ and $`𝐞_3`$ chosen so that $`𝐞_2=𝐤`$ and $`𝐞_3=𝐧`$. Using the spacelike vector $`𝐞_1`$ and the timelike vector $`𝐞_4`$, we now define $`𝐯_1\alpha (𝐞_4+𝐞_1)`$ and $`𝐯_2\alpha (𝐞_4𝐞_1)`$, where $`\alpha 0`$ is some constant. Evidently, these new vectors are null and each of them is orthogonal to both $`𝐤`$ and $`𝐧`$. Moreover, as $`𝐯_1𝐯_2=2\alpha ^20`$, at least one of them, say $`𝐯_1`$, is not parallel to the vector $`𝐮`$, i.e. we must have $`𝐮𝐯_10`$. By suitable choice of the constant $`\alpha `$ in the definition of $`𝐯_1`$ one can always normalize $`𝐯_1`$ so that $`𝐮𝐯_1=1`$. The pseudo-orthonormal tetrad $`\{𝐄_i\}`$ in the space $`T_q`$ can now be chosen as follows: $`𝐄_1𝐧`$, $`𝐄_2𝐤`$, $`𝐄_3𝐯_\mathrm{𝟏}`$ and $`𝐄_4𝐮`$. By parallely transporting this tetrad along $`\lambda (s)`$ one obtains the pseudo-orthonormal basis at each point of $`\lambda (s)`$. We shall now find one of the components $`C^i_{jkl}`$ of the Weyl tensor with respect to the tetrad $`\{𝐄_i\}`$. The components $`C_{mjkl}`$ of the Weyl tensor with respect to $`\{𝐄_i\}`$ can be obtained from the coordinate components $`C_{abcd}`$ of the Weyl tensor according to $$C_{mjkl}=C_{abcd}E_m^aE_j^bE_k^cE_l^d,$$ (5) where $`E_m^a`$ (resp., $`E_j^b`$, $`E_k^c`$ and $`E_l^d`$) is the $`a`$th (resp., $`b`$th, $`c`$th and $`d`$th) coordinate basis component of the vector $`𝐄_m`$ (resp., $`𝐄_j`$, $`𝐄_k`$ and $`𝐄_l`$) of the tetrad $`\{𝐄_i\}`$. Having the tetrad components $`C_{mjkl}`$, we can now easily find the tetrad components $`C^i_{jkl}`$ of the Weyl tensor: $$C^i{}_{jkl}{}^{}=\eta ^{im}C_{mjkl},$$ (6) where $`\eta ^{im}`$ is the inverse of the matrix $`\eta _{im}𝐄_i𝐄_m`$; that is, we have $`\eta ^{11}=\eta ^{22}=\eta ^{34}=\eta ^{43}=1`$ and $`\eta ^{im}=0`$ in all other cases. Since the tetrad $`\{𝐄_i\}`$ is chosen so that $`𝐄_2=𝐤`$ and $`𝐄_4=𝐮`$, and the vectors $`𝐤`$ and $`𝐮`$ are given in the explicit form, we can find, applying (5) and (6), an explicit expression for the component $`C^2_{424}`$ of the Weyl tensor with respect to $`\{𝐄_i\}`$. This is done in Appendix B; the result is $$C^2{}_{424}{}^{}=\frac{3m(LaE)^2}{r^5}.$$ (7) The task is now to show that if $`\lambda (s)`$ is not a member of the principal null congruences, then there exists some affine parameter interval $`(0,s_0]`$ of $`\lambda (s)`$ on which $`C^2{}_{424}{}^{}Ks^2`$, where $`K`$ is some positive constant. To do this, let us first recall that the affine parameter $`s`$ and the radial coordinate $`r`$ on $`\lambda (s)`$ are related by Eq. (2), with the + sign as $`dr/ds>0`$ on $`(0,s_1]`$. Let us now rewrite this equation in the form $$\frac{ds}{dr}=\frac{r^{3/2}}{F(r)},$$ (8) where $`F(r)\left[r^3E^2r(L^2a^2E^2)+2m(LaE)^2\right]^{1/2}`$. As $`dr/ds>0`$ on $`(0,s_1]`$, it is obvious that $`ds/dr>0`$ on $`(0,r_1]`$, where $`r_1`$ denotes the value of the coordinate $`r`$ of the point $`\lambda (s_1)`$. By (8) it is thus clear that $`F(r)`$ is strictly positive on $`(0,r_1]`$. It is also clear, as $`lim_{r0}F(r)=\|LaE\|\sqrt{2m}`$, that $`F(r)`$ is bounded on $`(0,r_1]`$. So there exists a positive number $`F_0sup\{F(r)\|0<rr_1\}`$. Consider now the function $`y(r)2r^{5/2}(5F_0)^1`$ defined on $`(0,r_1]`$. Since $`dy/dr=r^{3/2}/F_0`$ and $`F_0F(r)>0`$ on $`(0,r_1]`$, by (8) we have $`ds/dr\|_r^{}dy/dr\|_r^{}>0`$ for all $`r^{}(0,r_1]`$. From this it follows easily, as $`lim_{r0}y(r)=lim_{r0}s(r)=0`$, that $`s(r)y(r)`$ for all $`r(0,r_1]`$. Evidently, both $`s(r)`$ and $`y(r)`$ are strictly increasing on $`(0,r_1]`$, and so there exist their inverse functions. Let $`r:(0,s_1](0,r_1]`$ be the inverse function of $`s(r)`$ and let $`z:(0,s_0](0,r_1]`$ be the inverse function of $`y(r)`$ (note that $`s_0s_1`$ since $`s(r)y(r)`$ on $`(0,r_1]`$). From the fact that $`s(r)y(r)`$ for all $`r(0,r_1]`$, it follows immediately that $`r(s)z(s)`$ for all $`s(0,s_0]`$. Note also that $`z(s)=(5F_0s/2)^{2/5}`$, which is clear from the definition of $`y(r)`$. We thus have $`r(s)(5F_0s/2)^{2/5}`$ on $`(0,s_0]`$. Combining this inequality with (7), and taking into account the fact that $`m>0`$, we obtain $$C^2{}_{424}{}^{}=\frac{3m(LaE)^2}{r^5(s)}Ks^2$$ (9) for all $`s(0,s_0]`$, where $`K12m(LaE)^2(5F_0)^2`$. Suppose now that $`\lambda (s)`$ does not belong to either of the two principal null congruences; then we have $`LaE`$ (see Sec. 2), and hence $`K>0`$. By (9) and Proposition 2.1 it is thus clear that $`\lambda (s)`$ must terminate in Tipler’s strong curvature singularity as $`s0`$, which is the desired conclusion. ## IV Concluding remarks We have examined the curvature strength of the Kerr singularity with $`m>0`$ and $`a0`$. We have shown that every null geodesic reaching this singularity, with the exception of those belonging to the principal null congruences, must in fact terminate in Tipler’s strong curvature singularity. This is a typical property of null geodesics approaching the Kerr singularity because the principal null geodesics are very special and can be considered to form a set of “measure zero” in the family of all null geodesics reaching the singularity (Ref. , pp. 328, 329). The existence of these special geodesics is due to the high symmetry of the solution and one can expect that such geodesics will not occur in more general spacetimes. Thus if one attempts to define the curvature strength of a Kerr-like naked singularity, which could possibly arise in a generic collapse, one way of doing this is to assume that all null geodesics approaching such a singularity will behave much the same as a typical null geodesic approaching the Kerr singularity with $`\|a\|>m`$ —i.e. will also terminate in a singularity of the strong curvature type. Using this assumption, one may then attempt to formulate and prove a theorem which would constraint or prohibit the occurrence of Kerr-like naked singularities in generic collapse situations. A theorem of this type was established in Ref. . In this context, it is worth recalling that the inextendibility condition assumed in the theorem of Ref. may in fact hold for a much more general class of possible singularities than only those of the strong curvature type, for the curvature need not necessarily diverge along geodesics satisfying this condition . However, we have checked that this condition fails to hold for the principal null geodesics approaching the Kerr singularity (in fact, it will always fail to hold for principal null geodesics in any Ricci-flat spacetime). Thus the theorem of Ref. does not exclude the possibility that a naked Kerr singularity accompanied by closed timelike curves could develop from some non-singular initial data. However, from the proof of that theorem it may be concluded that this can happen only if the inextendibility condition fails to hold for all null geodesics terminating in the past at the naked Kerr singularity and generating the future Cauchy horizon due to the formation of this singularity. According to our result, this is possible only if all these geodesics belong to the principal null congruences. This would be a very special case. ## Acknowledgments We wish to thank the referee for helpful comments. This work was supported by the Polish Committee for Scientific Research (KBN) under grant No. 2 P03B 073 15. ## A In this appendix, we demonstrate that the vector $`𝐤`$ is parallely transported along the null geodesic $`\lambda (s)`$ with the tangent vector $`𝐮`$. To this end we only need to show that $$\left(\frac{k^c}{x^b}+\mathrm{\Gamma }^c{}_{ab}{}^{}k_{}^{a}\right)u^b=0,$$ (A1) where $`\mathrm{\Gamma }^c_{ab}`$ are the connection coefficients, which can be obtained from the metric tensor $`g_{ab}`$ according to $$\mathrm{\Gamma }^c{}_{ab}{}^{}=\frac{1}{2}g^{cd}(\frac{g_{bd}}{x^a}+\frac{g_{da}}{x^b}\frac{g_{ab}}{x^d}).$$ (A2) Since $`u^\theta =k^r=k^\varphi =k^t=0`$ and $`k^\theta =r^1`$, the components of Eq. (A1) in the $`(r,\theta ,\varphi ,t)`$ coordinate system take the form for $`c=r`$: $$\left(\mathrm{\Gamma }^r{}_{\theta r}{}^{}u_{}^{r}+\mathrm{\Gamma }^r{}_{\theta \varphi }{}^{}u_{}^{\varphi }+\mathrm{\Gamma }^r{}_{\theta t}{}^{}u_{}^{t}\right)r^1=0,$$ (A3) for $`c=\theta `$: $$r^2u^r+\left(\mathrm{\Gamma }^\theta {}_{\theta r}{}^{}u_{}^{r}+\mathrm{\Gamma }^\theta {}_{\theta \varphi }{}^{}u_{}^{\varphi }+\mathrm{\Gamma }^\theta {}_{\theta t}{}^{}u_{}^{t}\right)r^1=0,$$ (A4) for $`c=\varphi `$: $$\left(\mathrm{\Gamma }^\varphi {}_{\theta r}{}^{}u_{}^{r}+\mathrm{\Gamma }^\varphi {}_{\theta \varphi }{}^{}u_{}^{\varphi }+\mathrm{\Gamma }^\varphi {}_{\theta t}{}^{}u_{}^{t}\right)r^1=0,$$ (A5) for $`c=t`$: $$\left(\mathrm{\Gamma }^t{}_{\theta t}{}^{}u_{}^{r}+\mathrm{\Gamma }^t{}_{\theta \varphi }{}^{}u_{}^{\varphi }+\mathrm{\Gamma }^t{}_{\theta t}{}^{}u_{}^{t}\right)r^1=0.$$ (A6) Applying (A2) to the Kerr metric (1), we can now calculate the connection coefficients appearing in Eqs. (A3)-(A6); the result is $$\mathrm{\Gamma }^r{}_{\theta r}{}^{}=\frac{a^2\mathrm{sin}\theta \mathrm{cos}\theta }{\rho ^2},\mathrm{\Gamma }^\theta {}_{\theta r}{}^{}=\frac{r}{\rho ^2},\mathrm{\Gamma }^\varphi {}_{\theta \varphi }{}^{}=\frac{(\rho ^2+2a^2mr\mathrm{sin}^2\theta )\mathrm{cos}\theta }{\rho ^4\mathrm{sin}\theta },$$ $$\mathrm{\Gamma }^\varphi {}_{\theta t}{}^{}=\frac{2mar\mathrm{cos}\theta }{\rho ^4\mathrm{sin}\theta },\mathrm{\Gamma }^t{}_{\theta \varphi }{}^{}=\frac{2ma^3r\mathrm{sin}^3\theta \mathrm{cos}\theta }{\rho ^4},\mathrm{\Gamma }^t{}_{\theta t}{}^{}=\frac{2ma^2r\mathrm{sin}\theta \mathrm{cos}\theta }{\rho ^4},$$ $$\mathrm{\Gamma }^r{}_{\theta \varphi }{}^{}=\mathrm{\Gamma }^r{}_{\theta t}{}^{}=\mathrm{\Gamma }^\theta {}_{\theta \varphi }{}^{}=\mathrm{\Gamma }^\theta {}_{\theta t}{}^{}=\mathrm{\Gamma }^t{}_{\theta r}{}^{}=\mathrm{\Gamma }^r{}_{\theta \varphi }{}^{}=0.$$ Substituting these coefficients into Eqs. (A3)-(A6), and putting $`\theta =\pi /2`$ as $`\lambda (s)`$ lies in the equatorial plane, we can now readily see that Eqs. (A3)-(A6) are satisfied, as desired. ## B We give in this appendix some details of the computation of the component $`C^2_{424}`$ of the Weyl tensor with respect to the tetrad $`\{𝐄_i\}`$. Since $`𝐄_2=𝐤`$ and $`𝐄_4=𝐮`$, and $`u^\theta =k^r=k^\varphi =k^t=0`$ and $`k^\theta =r^1`$, the expression (5) for $`C_{2424}`$ takes the form $$C_{2424}=r^2\left[C_{\theta r\theta r}(u^r)^2+C_{\theta \varphi \theta \varphi }(u^\varphi )^2+C_{\theta t\theta t}(u^t)^2\right]$$ $$+2r^2\left(C_{\theta r\theta \varphi }u^ru^\varphi +C_{\theta r\theta t}u^ru^t+C_{\theta \varphi \theta t}u^\varphi u^t\right),$$ (B1) where the components $`u^r`$, $`u^\varphi `$ and $`u^t`$ of $`𝐮`$ are given by Eqs. (2)-(4), with the + sign in Eq. (2). We recall that the Weyl tensor $`C_{abcd}`$ is defined by $$C_{abcd}=R_{abcd}+g_{a[d}R_{c]b}+g_{b[c}R_{d]a}+\frac{1}{3}Rg_{a[c}g_{d]b},$$ where $`R_{abcd}`$, $`R_{ab}`$ and $`R`$ denote the curvature tensor, the Ricci tensor and the curvature scalar, respectively. Since the Kerr spacetime is Ricci flat, $`R_{ab}`$ and $`R`$ will vanish, and so $`C_{abcd}=R_{abcd}`$. In this case the coordinate components of the Weyl tensor can be obtained from the metric components $`g_{ab}`$ according to $$C_{abcd}=\frac{1}{2}(\frac{^2g_{ac}}{x^dx^b}+\frac{^2g_{bd}}{x^cx^a}\frac{^2g_{bc}}{x^dx^a}\frac{^2g_{ad}}{x^cx^b})+g_{ef}(\mathrm{\Gamma }^e{}_{ca}{}^{}\mathrm{\Gamma }_{}^{f}{}_{bd}{}^{}\mathrm{\Gamma }^e{}_{da}{}^{}\mathrm{\Gamma }_{}^{f}{}_{bc}{}^{}),$$ where the connection coefficients are given by (A2). Applying this formula to the Kerr metric (1), we can now calculate the coordinate components of the Weyl tensor appearing in (B1); the result is $$C_{\theta r\theta r}=\frac{mr(3a^2\mathrm{cos}^2\theta r^2)}{\rho ^2\mathrm{\Delta }},$$ $$C_{\theta \varphi \theta \varphi }=\frac{mr(3a^2\mathrm{cos}^2\theta r^2)[a^2\mathrm{\Delta }\mathrm{sin}^2\theta +2(r^2+a^2)^2]\mathrm{sin}^2\theta }{\rho ^6},$$ $$C_{\theta t\theta t}=\frac{mr(3a^2\mathrm{cos}^2\theta r^2)(2a^2\mathrm{sin}^2\theta +\mathrm{\Delta })}{\rho ^6},$$ $$C_{\theta \varphi \theta t}=\frac{mra(3a^2\mathrm{cos}^2\theta r^2)[\mathrm{\Delta }+2(r^2+a^2)]\mathrm{sin}^2\theta }{\rho ^6},$$ $$C_{\theta r\theta \varphi }=C_{\theta r\theta t}=0.$$ Inserting these expressions in (B1), and putting $`\theta =\pi /2`$ as $`\lambda (s)`$ lies in the equatorial plane, we get $$C_{2424}=\frac{3m(LaE)^2}{r^5}.$$ Finally, we note that $`C_{2424}=C^2_{424}`$, which is clear from (6).
warning/0003/astro-ph0003177.html	ar5iv	text	# The DENIS Point Source Catalogue towards the Magellanic Clouds Based on observations collected at the European Southern Observatory ## 1 Introduction The DENIS project aims to survey the entire southern hemisphere simultaneously in three photometric bands, $`I`$ (Gunn–i, $`0.8\mu m`$), $`J`$ ($`1.25\mu m`$) and $`K_s`$ ($`2.15\mu m`$) with a spatial resolution of $`1\mathrm{}`$ in $`I`$ and $`3\mathrm{}`$ in the $`J`$ and $`K_s`$ bands, and limiting magnitudes of $`I=18`$, $`J=16`$, $`K_s=14`$. See Epchtein et al. (epal (1999)) for the first general release of DENIS data. Here we present a catalogue of DENIS Point Sources towards the Magellanic Clouds, requiring that objects are detected in at least two of the three photometric bands. At the distance of the Magellanic Clouds, $`(mM)=18.45\pm 0.1`$ for the LMC and $`(mM)=19.0\pm 0.1`$ for the SMC according to Westerlund (west (1997)), our catalogue contains: (1) all Asymptotic Giant Branch stars (AGB), except those with shells optically thick at $`2\mu m`$ and the faintest stars at the very beginning of the Early AGB branch (E-AGB), (2) upper Red Giant Branch stars (RGB), (3) most of the super-giants except those brighter than $`I=10.5`$, $`J=8.0`$, $`K_s=6.5`$ because they saturate the detectors, (4) relatively bright post-AGB stars. The catalogue will thus be a major tool for statistical studies of the post-main sequence stellar populations of the Magellanic Clouds. Dwarfs and giants are the main galactic sources seen in front of the Magellanic Clouds (Ruphy et al. ruph (1997)). Compared to earlier spectroscopic and photometric surveys of the Magellanic Clouds for red giants and super-giants, and for stars on the AGB, probably we find a few hundreds times more sources, for several reasons: (1) previous surveys were not sensitive enough (Westerlund west0 (1960), west1 (1961); Sanduleak & Philip saph (1977); Westerlund et al. weal1 (1978), weal2 (1981); Rebeirot et al. reb0 (1983)), (2) they were spatially limited (see e.g. Blanco, McCarthy, and Blanco bmc (1980); Blanco & McCarthy blamc (1990)), (3) they were restricted to a peculiar type of objects (e.g. Hughes hug (1989) in his search for Miras variables, Rebeirot et al. reb (1993) in their search for carbon stars). About $`1/4`$ of the sources discovered in these surveys were later observed in the $`JHK(L)`$ infrared photometric bands (e.g. Hughes & Wood hugwo (1990), Costa & Frogel cofrog (1996)). DENIS provides simultaneous $`IJK_s`$ observations of the entire Clouds, with a good sensitivity, and connecting for the first time the traditional optical and infrared wavelengths domains by simultaneous observations. Sect. $`2`$ describes the instrument characteristics and the observing technique. Sect. $`3`$ describes the data reduction procedure in the two “data analysis centers“ with particular attention to: flat and bias subtraction, point spread function, and astrometric and photometric calibration. Sect. $`4`$ discusses the quality of the data with regard to the selection criteria applied and to the completeness reached. Sect. $`4.2`$ discusses in particular the foreground sources belonging to our Galaxy. Finally, Sect. $`5`$ describes the content of the catalogue and Sect. $`6`$ gives conclusive remarks. The catalogue is available through the Strasbourg Astronomical Data Center (CDS); it carries the number II/228. ## 2 Observations The DENIS instrument is mounted at the Cassegrain focus of the 1–m ESO telescope (La Silla - Chile). It contains three cameras: a Tektronix CCD with $`1024\times 1024`$ pixels and two NICMOS infrared detectors with $`256\times 256`$ pixels. The array of the camera has four quadrants to reduce the read-out time, and each quadrant, especially in the $`I`$ band, presents different image characteristics and must be treated separately. The pixel sizes are $`1\mathrm{}`$ in $`I`$ and $`3\mathrm{}`$ in $`J`$ and $`K_s`$, respectively. The total integration time is $`9`$ secs for each image. The sampling of the image is $`1\mathrm{}`$ in all three wave bands. The $`J`$ and $`K_s`$ images are dithered to a $`1\mathrm{}`$ pseudo–resolution, using a microscanning mirror. They consist of a set of $`9`$ frames each obtained in $`1`$ sec integration time, shifted by $`\pm 1/3`$ pixel in right ascension (RA) plus $`\pm 7/3`$ pixel in declination (DEC). The DENIS strategy is to divide the sky into three declination zones and scan each in strips of $`30\mathrm{°}`$ in DEC and $`12\mathrm{}`$ in RA. The overlap in RA between consecutive strips is $`2\mathrm{}`$. Each strip consists of 180 images of $`12\mathrm{}\times 12\mathrm{}`$ with an overlap of $`2\mathrm{}`$ between each image. The observation of a photometric standard star consists of $`8`$ sub-images shifted according to a circular pattern in order to have the star always at a different position on the chip. One $`I`$ standard is observed before the observation of a strip and one $`J`$ and $`K_s`$ standard afterwards. On average $`6`$ to $`8`$ strips per night are observed. Data on the MCs were taken during observing seasons from August to March, the first centered on December 1995, and the last one on December 1998. The two clouds, LMC and SMC, were covered by $`119`$ and $`88`$ strips, respectively. ## 3 Data Reduction Data reduction took place in two centers: the Paris Data Analysis Center (PDAC) and the Leiden Data Analysis Center (LDAC). PDAC pre–processed the raw data and LDAC extracted and parametrized objects ranging from point sources to small extended sources. The reduction at LDAC by the first author led to the detection of numerous technical problems (astrometry for example), which had escaped the checks of the automatic pipeline. Compared to the first relase of DENIS data (Epchtein et al. epal (1999)) this catalogue differs in terms of flags, astrometric reference catalogue, association criteria and photometric calibration (use of the overlapping region between adjacent strips). Besides, our catalogue covers a portion of the sky not overlapping with the first DENIS data release. The source list (table 7) and the photometric information (table 8) are electronically available at CDS via this catalogue. With further DENIS data releases data on a strip by strip basis, therefore not merged into a single catalogue without treating the overlapping regions in terms of photometry and astrometry, will also be available. This means that all the strips covering the same region of the sky regardless of their quality will be available; multiple entries for a single objects could then be retrieved. We show in section 3.2.3 the consistency of our calibration within each cloud. The work of Fouqué et al. fou (1999) focuses on the absolute photometric calibration of the DENIS data. This calibration, once completed after the termination of the survery, might induce a systematic shift on the photometry of the present catalogue. ### 3.1 PDAC: Paris Data Analysis Center At the PDAC the images were corrected for sensitivity differences and atmospheric and instrumental effects. Their optical quality was judged on the basis of the parameters that describe the point spread function (see Sect. $`\mathrm{3.1.2}`$). #### 3.1.1 Flat and Bias The received intensity from the target image also contains background contribution from the telescope and the atmospheric radiation. Besides, the sensitivity varies across the array of the camera. The true signal ($`TS`$) for the pixel $`i,j`$ is obtained from: $$TS_{i,j}=I_{i,j}(t)F_{i,j}\times b(t)B_{i,j},$$ (1) where $`I_{i,j}`$ is the measured intensity after subtraction of dark currency, $`F_{i,j}`$ is the flat field after dark subtraction, $`b(t)`$ is the background and $`B_{i,j}`$ is the bias level. The background is estimated per image, at time $`t`$, with $$b(t)=\frac{_{ij}^N[I_{i,j}(t)/F_{i,j}]}{N},$$ (2) where $`N`$ is the total number of pixels per image. At this stage we use the flat field and the dark current values estimated from the previous night are used. Points outside the $`3\sigma `$ level are rejected in the sum. The four $`I`$ quadrants are treated separately. This is also done for each of the $`9`$ sub-images in the $`J`$ and $`K_s`$ bands. As a second step we select low–background images from the sunrise sequence ($`180`$ in a normal strip) with a low background value. To identify and avoid crowded fields and fields affected by saturated stars we combine measurements taken during different nights. We then determine the flat $`F_{i,j}`$ and the bias $`B_{i,j}`$ by minimizing the expression: $$\underset{t=1}{\overset{N}{}}[I_{i,j}(t)F_{i,j}b(t)B_{i,j}]^2,$$ (3) where $`N`$ is the number of selected images ($`180`$). In the third step, we applied the new values of the flat and the bias to the set of selected images to obtain a new estimate of the background value, more appropriate to the particular night. The quality of the determination of the parameters involved is improved by iteration of the above procedure. The bias so far determined is a mean value for the night. Because it varies during the night, its value for a given strip is estimated to be: $$B_{i,j}=\frac{_{t=1}^N[I_{i,j}(t)F_{i,j}b(t)]}{N},$$ (4) where $`N`$ is the total number of images per strip ($`180`$). After dark subtraction, the bias contains only the contribution of the instrumental and atmospheric emission which does not affect the $`I`$ band, but does affect strongly the $`K_s`$ band and a higher number of iterations is sometimes necessary. The large number of available flat/bias-images ($`180`$) gives a quite high degree of statistical confidence to both determinations. This is not true in case of calibration sequences that involve only $`8`$ images. In this case, the bias determined for the strip nearest in time is applied. #### 3.1.2 Point Spread Function The pixel size of the $`J`$ and $`K_s`$ channels is $`3\mathrm{}`$ and the sampling is $`1\mathrm{}`$ in both directions. The real width of any point source is therefore potentially narrower than the pixel ($`3\mathrm{}`$). In terms of signal processing the sources are not under-sampled, but the width of the filter is broader than the sampling. To estimate the width of the signal the convolution of the signal profile (assumed to be elliptical) and the pixel size has been taken into account. The method of least squares has been applied to the projection of sources onto RA, DEC and diagonal axes (Borsenberger, bors (1997)). In the $`I`$ and the $`J`$ bands there are enough sources to build a model that describes the behaviour of the projected widths in each image. In the $`K_s`$ band several images were stacked together prior to the determination. We refer to http://www-denis.iap.fr/docs/tenerife.html for more details on the PDAC data reduction. ### 3.2 LDAC: Leiden Data Analysis Center LDAC extracts point sources from the images delivered by PDAC. From these sources, it derives and then applies astrometric and photometric calibration to obtain a homogeneous point source catalogue. The astrometric reference catalogue is the USNO–A2.0 (Monet monet (1998)) that provides on average $`100`$ “stars” per DENIS image. The photometric DENIS standard stars belong to different photometric systems of which the major ones are: Landolt (lan (1992)), Graham (gra (1982)), Stobie et al. (sgr (1985)) and Menzies et al. (mcbl (1989)) in $`I`$; Casali & Hawarden (cashaw (1992)), Carter (car (1990)) and Carter & Meadows (carmea (1995)) in $`J`$ and $`K_s`$. An absolute calibration, together with a definition of DENIS photometric bands is given by Fouqué et al. (fou (1999)). #### 3.2.1 Source Extraction The first LDAC task is to reduce the information from each image into an object list. This is done using the SExtractor program (Bertin and Arnouts emb (1996)) version $`\mathrm{2.0.15}`$. #### 3.2.2 Astrometric Calibration Positions are determined through pairing information among frames, channels and with the reference catalogue. The astrometric solution makes use of the fact that each map has an area of overlap with neighboring maps, and that objects in the overlapping region have been observed many times. The projected position of the multiply observed sources, in terms of their pixel positions, contains information on the telescope pointing and the plate deformations. The plate deformation is derived through a triangulation technique, matching bright extracted objects with astrometric reference objects. The resulting global solution for each strip takes into account possible variations along the strip. The plate offsets are determined using all but the faintest extracted objects, matching among channels (wave bands) and in overlap. A least square fitting technique is then applied to the functional description of the detector deformation and its variation to obtain the full solution on the basis of the pairing information. Thereafter, the celestial position, its error and the geometric parameters of each object are calculated. The standard position accuracy derived is RMS $`0.001`$ arc sec with maximum excursions of $`1.32`$ arc sec. This error is in addition to the RMS of $`0.3`$ arc sec of the astrometric reference catalogue. #### 3.2.3 Photometric Calibration Magnitudes are estimated within a circular aperture of $`7\mathrm{}`$ in diameter after a de–blending process, that determines which pixels are within the aperture, and what fraction they contribute to each individual source (Bertin & Arnouts, emb (1996)). This aperture collects $`95\%`$ of the light when considering a seeing of $`1.5\mathrm{}`$ and the pixel size of $`3\mathrm{}`$ for the infrared wave bands. For homogeneity we used the same aperture also for the $`I`$ band. The source magnitude ($`m`$) corresponding to the wavelength $`\lambda `$ is defined as: $$m_\lambda =2.5\mathrm{log}(S_\lambda )+m_{\lambda 0},$$ (5) where $`S_\lambda `$ is the observed flux and $`m_{\lambda 0}`$ defines the zero–point of the magnitude scale at the wavelength $`\lambda `$. The determination of the instrumental quantity $`m_{\lambda 0}`$ to correct the stellar magnitude for atmospheric effects is done on a nightly basis. First, standard star measurements are matched with the information stored in the standard star catalogue. Second the instrumental zero-point ($`m_{\lambda 0}`$) is derived for each of the eight measurements of the standard star assuming a fixed extinction coefficient, $`ϵ`$ (Eq. 6). The adopted values of $`ϵ`$ are $`0.05`$ for the $`I`$ band and $`0.1`$ for both the $`J`$ and $`K_s`$ bands. These values have been determined from the photometric measurements performed during calibration nights (nights where only standard stars were observed). $$m_{\lambda 0}=2.5\mathrm{log}(S_\lambda )+m_{\lambda \mathrm{ref}}+ϵ\times z$$ (6) $`m_{\mathrm{ref}}`$ is the magnitude of the standard star from the standard star catalogue and $`S_\lambda `$ is the flux as measured at a given air mass ($`z`$). Standard stars were selected near the airmass limit of the strips and to be roughly of the same spectral type; this simplifies the Taylor expression used to describe the extinction law because colour terms (Guglielmo et al. gugli (1996)) and the non–linear terms are of minor importance; in the infrared the dependence of the extinction on $`z`$ is almost linear for $`z<2`$. In principle, both $`m_{\lambda 0}`$ and $`ϵ`$ can be determined simultaneously and the non–linear terms can be incorporated as well if a sufficient number of star measurements are available, but for a single night there are not enough, in fact the use of the approximated law (Eq. 6) gives a systematic offset between the magnitude of the source in the overlap of two strips of comparable, but different, photometric conditions. After a considerable investigation it turned out that this offset could be greatly reduced if a fixed extinction coefficient is used. Some differences are left when the observations have been performed in different photometric conditions or when too few standard star measurements were done. Fig. 1 shows the computed differences between the magnitudes of the sources detected in the overlapping region of two strips observed under comparable photometric conditions ((a), (b), (c)) and of two strips observed with different photometric conditions ((d), (e), (f)) in the $`I`$, $`J`$ and $`K_s`$ bands, respectively. Faint sources give rise to a larger dispersion. The systematic shift is clearly visible in Fig. 1d–f. The final nightly value of $`m_{\lambda 0}`$, for each wave band, is calculated by averaging the single determinations for each standard star and among all the standard stars observed during that night, after removal of flagged (Sect. 4) measurements (this reduces on average the number of measurements per star from $`8`$ to $`6`$). The flagged measurements have a non-zero value for at least one of the types of flag considered in the pipeline reduction. Only standards fainter than $`I=10.5`$, $`J=8.0`$ and $`K_s=6.5`$ mag are used. The instrumental $`m_{\lambda 0}`$ and its standard deviation are listed for each strip in the quality table (table 8). Mean values ($`\pm 1\sigma `$) are: $`23.42\pm 0.07`$ ($`I`$), $`21.11\pm 0.13`$ ($`J`$) and $`19.12\pm 0.16`$ ($`K_s`$). Using the overlapping regions of adjacent strips to correct for remaining differences we performed a general photometric calibration, separately for the LMC and for the SMC. We calculated the magnitude difference of cross-identified sources between two adjacent strips of sources detected in three wave bands. The histogram of these differences in magnitude shows when a systematic shift is present between the two strips (Fig. 1). In only a few cases is the average magnitude affected by more than 0.1 mag. If necessary we applied a systematic shift (table 8). Experience showed that if a strip is poorly calibrated the magnitude difference in the overlap with the previous strip has a sign opposite to the difference found in the overlap with the next strip. Note that observing a strip in good photometric conditions but having too few standard star measurements to perform the calibration may induce alone this offset; the equal number of detected objects as a function of magnitude per band in both strips indicates, as just mentioned, a minor difference in the photometric conditions under which each strip was observed, increasing the confidence we place in the correction procedure. Only $`9`$ strips out of $`108`$ for LMC observations and $`3`$ strips out of $`81`$ for SMC observations show this behaviour. The maximum observed shift amounts to 2.45 mag in the $`J`$ band. Sources with corrected magnitude are easily recognized from their strip number associated to each detected band (Table 7). Table 8 reports the amount of the applied shift as a function of strip number. In some cases, the difference shows a dependence on declination, but the effect on the averaged magnitude, in the area of the Magellanic Clouds, is not significant (less than $`0.1`$ mag), and can be ignored. The internal statistical RMS error is between $`0.001`$ and $`0.4`$ mag at the detection limit, faint sources have larger errors. For completeness we included in the catalogue sources detected above and below the reference saturation and detection limits, their photometric errors (larger than $`0.4`$ mag) show the confidence of the detection. The standard deviation on $`m_{\lambda 0}`$ is in most of the strips below $`0.05`$ mag, but spreads from $`0.01`$ to $`0.2`$ mag. Larger values are detected in the strips where a photometric shift was also applied, therefore the resulting accuracy is, for these few cases, not better than $`0.1`$ mag. In all other cases the resulting accuracy has an RMS error better than $`0.05`$ mag. #### 3.2.4 Association All extracted objects are matched on the basis of their geometrical information assuming an elliptical shape (RA, DEC, $`a`$: semi–major axis, $`b`$: semi–minor axis, $`\theta `$: inclination angle) within one wave band, among the three wave bands within a strip and among different strips. The geometrical parameters of each object are evaluated at the $`3\sigma `$ level of the row image; $`a`$ and $`b`$ are the second order moments of the pixel distribution within the size of the photometric aperture. Typical values are $`1.8\mathrm{}`$, $`1.0\mathrm{}`$ and $`0.5\mathrm{}`$ for the $`I`$, $`J`$ and $`K_s`$ band, respectively, differences among the three wave bands mainly depend on the differences in sensitivity; the second order moments characterize the PSF. The effective area used during the association procedure is $`1.5`$ times (tolerance) the area defined by the $`a`$ and $`b`$ values of both object, when the association is performed within each band of a strip. Sources previously de–blended are not associated. When the association is done among different bands the tolerance value increases to $`2.5`$. We associate two objects when the center position of one of them is within the bounds of the ellipse of the other, even if the center of the second is outside the ellipse of the first one, and vice versa. For the coordinates, we always used a weighted average (based on the signal to noise ratio and detection conditions as derived from the source extraction program and the astrometric calibration). For the magnitudes we decided not to average or to combine magnitudes from different epochs (strips) because of the possible variability of a large fraction of the detected objects. Objects associated within the same strip are given with the average of the magnitudes. When the association involves overlapping strips we distinguish the following cases: (1) for objects detected in all three wave bands in both strips we choose the entry from the strip with the lowest value of $`_{i=1}^N\sqrt{a_ib_i}`$, where N is the number of sources detected in the overlap; (2) for objects detected in an unequal number of wave bands, we chose the entry from the strip with the highest number of detected wave bands; (3) for objects detected in two different wave bands we choose the entry from the strip with the lowest $`_{i=1}^N\sqrt{a_ib_i}`$, including the third magnitude from the other strip. When the strip numbers of the detected wave bands differ the observations refer to different epochs. The criteria given conserve the major property of the DENIS data: simultaneousness. We refer to ftp.strw.leidenuniv.nl /pub/ldac/pipeline.ps for more details on the LDAC data reduction. ## 4 Data Quality For each image PDAC flags problems of different kinds (see table 1). The flags are used by LDAC to identify image defects and consequently flag the extracted objects if necessary. During the source extraction process, LDAC produces more flag information (tables 2 and 3). Artifact flags are not present in the output parameter list because they have been used as a primary selection criteria to filter the catalogue; most of the cosmic rays, glitches and optical ghosts will have been eliminated. When dust is present on the mirrors of the telescope and on the lenses of the instrument spurious objects are created. Most are bright and easily recognizable. During the pipeline reduction photometric fluxes are calculated in $`7\mathrm{}`$ and in $`15\mathrm{}`$ apertures. “Dusty–like objects” give a negative flux in the larger aperture (and its value is set to $`99`$). In their proximity the flat value for the pixels is dominated by their continuous presence in all images along the strip, therefore we end up with an area with negative flux next to the ’dusty–like objects’. This area is not always in the same position because of bending of the telescope during the observation of a strip. To eliminate these spurious detections we required that both aperture fluxes were positive. This selection also allowed the removal of glitches not previously flagged, sources too close to the image borders or too close to broad dead pixel regions and dummy sources with photometric errors greater than $`0.2`$ mag. An additional filtering criterion is based on the diagram of the isophotal area of one object at the $`1\sigma `$ level of the raw image (Isophotal area–pixels) versus the peak intensity (ln(MaxVal), Peak intensity–ADU); see Fig. 2. Area (3) of point sources (stars) is clearly identified: the objects have a Gaussian intensity energy distribution. Area (2): galaxies are extended objects and, relative to stars, their area increases faster for increasing intensity–ADU. The broadening of the locus of stars is due to the variation of the PSF over the field and of the seeing. Areas (1) and (4) contain cosmic rays and electronic glitches and are easily distinguishable. We accepted only sources in areas (2) and (3); the same cut between stars/galaxies/glitches–cosmic rays was applied to all strips. Finally we eliminated sources for which the object PSF could not match the instrumental PSF. This led to the loss of a few percent of $`K_s`$ detections; this effect does not depend on the source brightness and arises as a consequence of image defects. Filtering based on the flags, on dust on the detector and on the previous diagram were applied before we made the cross identification between the different wave bands. ### 4.1 Completeness A few strips had to be rejected during the reduction phase because of poor quality. These strips have been re–observed, but the data reduction is not yet started and the strips have not been included in the catalogue. Table 4 lists the right ascensions of the absent strips (each covers $`1^m20^s`$ in RA). Considering the overlap with adjacent strips we have missed $`7.7`$ % and $`6.3`$ % of the LMC and SMC regions, respectively. We now consider the completeness of the catalogue under two different aspects: completeness of objects detected in only two wave bands or in all three wave bands. Fig. 3 displays histograms of the number of sources in the catalogue in $`0.05`$ mag bins. Fig. 3a–d refer to the LMC and Fig. 3e–h to the SMC. Table 5 contains the magnitude of the maxima in the various histograms. A full discussion of these histograms will be given elsewhere (M.R. Cioni, H.J. Habing, M. Messino, in preparation). We limit ourselves to a few comments. (1) Comparing Fig. 3a and Fig. 3b, and similarly Fig. 3e and Fig. 3f suggests that (3b) and (3f) contain sources similar to (3a) and (3e), but they are below the detection limit in the $`K_s`$ band. Fig. 3b and 3f contain many more sources than Fig. 3a and 3e, respectively. (2) The $`I`$ and $`K_s`$ histograms of (3d) and (3h) are approximately scaled down versions of the $`I`$ and $`K_s`$ histograms in (3a) and (3e). This suggests that they contain the same kind of sources, and that the sources in (3d) and (3h) have not been detected in the $`I`$ band, i.e. the detection rate in the $`J`$ band is never $`100`$%, although it will be very close. (3) The nature of the sources in (3c) and (3g) remain unspecified for the moment. (4) The magnitudes of the maximum count as given in table 5 show that the magnitudes referring to the SMC are about $`0.25`$ mag fainter – this reflects the larger distance to the SMC. This conclusion is not true for the counts of sources detected only in $`I`$ and $`J`$. These counts may contain a large foreground component. Fig. 4 displays the cumulative distributions of the sources in the catalogue. From the overlap of adjacent strips, in the same wave band, we estimate a $`5`$% difference in the number of detected sources. This difference is partly due to regions of insensitive pixels on the frame borders, especially in the $`K_s`$ band. ### 4.2 Galactic Foreground Sources Galactic sources in the foreground have not been removed from the catalogue. Therefore, we now discuss the probability that any given source belongs to the Magellanic Clouds or to the Milky Way Galaxy. Fig. 5a shows that the count of sources detected in all three wave bands has a strong maximum inside the LMC area i.e. $`69\mathrm{°}>\delta >71\mathrm{°}`$. Outside of this area the count falls down to a plateau at an average value of $`50`$ sources per $`0.5`$ degrees in declination; this plateau represents the foreground contribution. In Fig. 5b we show the colour–colour diagram of all sources within the peak area of the LMC, and in Fig. 5c for all sources outside of the LMC. The foreground sources in (5c) are probably ordinary dwarf stars and red giants, for which we expect colours ($`0.5`$,$`0.5`$) and ($`1.0`$,$`1.0`$), respectively (Bessell & Brett besbret (1988)). The area outside the LMC is about $`7`$ times the area used in Fig. 5b and this explains why the total number of objects within ($`JK_s<1`$) and ($`IJ<1`$) is much larger in Fig. 5c than in Fig. 5b: the fraction of foreground sources in Fig. 5b is very small indeed. Fig. 5d, (e) and (f) refer to sources detected in three wave bands plus sources detected only in $`I`$ and $`J`$. The comparison between Fig. 5d (in the LMC) and Fig. 5e (outside the LMC) shows again what sources may be galactic and what sources are not. Sources in Fig. 5d with $`I<16`$ and $`IJ>1.2`$ are almost all LMC objects. The same is true for sources with $`IJ<0.4`$; these are probably early type main–sequence stars in the LMC. Sources with $`I>16`$ and $`IJ>1`$ are foreground objects. Fig. 5f shows the histogram obtained by adding up all sources in Fig. 5d (full drawn line) and in Fig. 5e (dashed line) irrespective of the value of $`I`$. The difference in distribution of points between Fig. 5d and Fig. 5e is obvious. From all strips and all colours we conclude that, on average, $`30`$% of the sources in the catalogue belong to the Galaxy rather than to the Magellanic Clouds. See also Cioni et al. (cioni (1998)) for the separation of foreground and Magellanic stars within DENIS data. A more elaborate discussion will be presented later (M-R. Cioni and H.J. Habing, in preparation). ### 4.3 Confusion When the source density is too high sources will blend with other sources, a process usually called confusion. A critical value is $`1`$ source per $`50`$ detection elements (IRAS explanatory supplement, vol. 1, VIII–4): if the source density is higher confusion becomes statistically probable. The typical size of a detected source does not exceed $`2\mathrm{}`$; see Sect. 3.2.4. Fig. 6 and Fig. 7 contain contour diagrams of source density in bins of constant right ascension. The maximum values is $`500`$ sources in $`0.25\times 0.1`$ square degrees in the LMC at $`\delta =70.5`$ which implies $`1`$ source per $`200`$ arcsec<sup>2</sup>. This is well below the confusion limit. Note that the confusion is not set by the size of the photometric aperture because the area of the aperture is independent of the detection process. Within the aperture there might be two de–blended sources, each pixel belongs to one or the other source or is shared between the two; the size of the aperture represents the contour limit where this pixel association process has to stop. ## 5 Contents of the catalogue The present version of the DENIS point source catalogue of the Magellanic Clouds contains sources detected in at least two wave bands within the area $`4^h08^m00^s<\alpha <6^h46^m40^s`$, $`61\mathrm{°}>\delta >77\mathrm{°}`$ for the LMC and $`0^h05^m20^s<\alpha <2^h00^m00^s`$, $`68\mathrm{°}>\delta >78\mathrm{°}`$ for the SMC of which the source density is shown in Fig. 6 and Fig. 7, respetively. The fraction of non–real objects is negligible as most glitches are present only in one wave band. Fig. 6 and Fig. 7 show the density maps of all detected sources towards the LMC and the SMC, respectively. The strips not yet included are explicitely indicated in the upper horizontal axis. Fig. 6 contains about $`\mathrm{1\hspace{0.17em}300\hspace{0.17em}000}`$ sources and Fig. 7 contains about $`\mathrm{300\hspace{0.17em}000}`$. Table 6 reports the approximate number of sources detected in three or two wave bands. The two parts of the catalogue are ordered by increasing RA. Two tables define the meaning of the columns in each part: a table that contains the detected sources (table 7) and a table that describes the quality of the detections on a strip by strip basis (table 8). ### 5.1 Nomenclature #### 5.1.1 Data table (table 7) The star name is composed by the acronym DCMC (DENIS Catalogue of the Magellanic Clouds) and by the coordinates of the source at the epoch 2000 (Columns 1 and 2) following the IAU convention ‘(Dubois et al. dub (1994)). For example, a star with $`(\alpha ,\delta )=(1^h25^m20.15^s,73\mathrm{°}30\mathrm{}15.67\mathrm{})`$ would have the designation DCMC $`J012520.15733015.6`$; right ascension is truncated at the second decimal and the declination at the first. Columns 3, 4, 5, 6, 7 and 8 give the right ascension ($`h`$, $`m`$, $`s`$) and the declination ($`\mathrm{°}`$, $`\mathrm{}`$, $`\mathrm{}`$) at the epoch J2000. Column 9 gives the positional error, i.e. the statistical error calculated during the data processing. Columns 10 and 11 give the pixel coordinates in the corresponding image (Column 12) of the corresponding strip (Column 13), in the $`I`$ band, where the source is detected. The same quantities are given in Columns 18, 19, 20 and 21 for the $`J`$ band and in Columns 26, 27, 28 and 29 for the $`K_s`$ band. Columns 14, 22 and 30 give the source magnitudes and Columns 15, 23 and 31 give the associated statistical errors in the $`I`$, $`J`$ and $`K_s`$ bands, respectively. Columns 16, 24 and 32 give the extraction flag (table 2) in the three wave bands. Columns 17, 25 and 33 give the image flag (table 1) in the three wave bands. Columns 34 and 35 give the $`B`$ and $`R`$ magnitudes from the cross-identification with the USNO–A2.0 catalogue. #### 5.1.2 Strip quality table (table 8) Column 1 gives the strip number. Column 2 gives the date of observation of the given strip. Columns 3 and 4 give the nightly zero–point and its standard deviation for the $`I`$ band. The same quantities are given in Columns 5 and 6 for the $`J`$ band and in Columns 7 and 8 for the $`K_s`$ band. Column 9 notes information peculiar to the strip in question, for example the photometric shift (Sect. $`\mathrm{3.2.3}`$). #### 5.1.3 Access The DENIS Point Source Catalogue towards the Magellanic Clouds, data (table 7) and quality (table 8) tables, are electronically available from CDS via http://cdsweb.u-strasbg.fr/denis.html. ## 6 Concluding remarks The catalogue is a suitable tool for the study of late–type stars in the Magellanic Clouds. These studies may involve the statistical separation of various species of stars, i.e. RGB and AGB (both O–rich and C–rich); the characterization of the mass loss properties of these stars, when combined with measurements in the mid and far–IR; the relations of infrared colours and magnitudes with variability, when combined with measurements of light curves (EROS, MACHO) or comparable photometric data (2MASS); the interpretation of the Hertzsprung–Russel diagram through theoretical evolutionary models; the investigation of metallicity effects inside the Magellanic Clouds and in comparison with our own Galaxy; the study of the history of star formation. ###### Acknowledgements. The authors kindly thank Ian Glass for his useful comments and J.L. Chevassut, F. Tanguy and K. Weestra for their technical support, the whole DENIS staff and all the DENIS observers who collected the data. The DENIS project is supported by the SCIENCE and the Human Capital and Mobility plans of the European Commission under grants CT920791 and CT940627 in France, by l’Institut National des Sciences de l’Univers, the Ministère de l’Education Nationale and the Centre National de la Recherche Scientifique (CNRS) in France, by the State of Baden–Württemberg in Germany, by the DGICYT in Spain, by the Sterrewacht Leiden in Holland, by the Consiglio Nazionale delle Ricerche (CNR) in Italy, by the Fonds zur Förderung der wissenschaftlichen Forschung and Bundesministerium für Wissenschaft und Forschung in Austria, by the Fundation for the development of Scientific Research of the State São Paulo (FAPESP) in Brazil, by the OKTA grants F–4239 and F–013990 in Hungary, and by the ESO C & EE grant A–04–046.
warning/0003/hep-ph0003321.html	ar5iv	text	# Measuring 𝛾 Cleanly with CP-Tagged 𝐵_𝑠 and 𝐵_𝑑 Decays \[ ## Abstract We propose a new method for measuring the CKM phase $`\gamma `$ using the partial rates for CP-tagged $`B_s`$ decays. Such an experiment could be performed at a very high luminosity symmetric $`e^+e^{}`$ collider operating at the $`\mathrm{{\rm Y}}(5S)`$ resonance, where the $`B_s\overline{B}_s`$ state is produced in a state of definite CP. We also discuss CP-tagging in the $`B_d`$ system at the $`\mathrm{{\rm Y}}(4S)`$, where a time-dependent analysis is required to compensate for the anticipated large CP violation in $`B_d\overline{B}_d`$ mixing. \] The accurate determination of the Unitarity Triangle of the Cabibbo-Kobayashi-Maskawa (CKM) quark mixing matrix is one of the most important problems of experimental $`B`$ physics. Reasonably precise measurements of some of its parameters – three sides, determined by $`\|V_{ub}\|`$, $`\|V_{cb}\|`$ and $`\|V_{td}\|`$, and one angle $`\mathrm{sin}2\beta `$ – will soon be performed at the $`B`$ Factories operating at the $`\mathrm{{\rm Y}}(4S)`$ and at Run II of the Tevatron. Although this is enough to fix the triangle up to discrete ambiguities, one really wants to overconstrain the system and thereby be sensitive to deviations from the CKM description of flavor-changing processes. In view of this goal, it is important to measure the angles $`\alpha `$ and $`\gamma `$ as well. The situation with these other angles is more problematic. A number of methods have been proposed to measure or constrain $`\alpha `$ and $`\gamma `$, but unfortunately they each suffer to some degree from either theoretical or experimental difficulties . In what follows we investigate a new proposal to constrain $`\gamma `$ in the decays of two $`B_s`$ mesons produced in a coherent state. This can be achieved if the pair comes from the decay of a $`b\overline{b}`$ meson such as the $`\mathrm{{\rm Y}}(5S)`$. In this case one has not only the option of tagging the flavor of the initial $`B_s`$, but the alternative of tagging it as an eigenstate of CP. The advantage of the $`B_s`$ in this regard is that CP violation in its mixing is small in the Standard Model. We will also discuss an analogous proposal for the $`B_d`$ system at the $`\mathrm{{\rm Y}}(4S)`$, where large CP violation in mixing complicates the situation. The CKM angle $`\gamma `$ is defined to be $$\gamma =\mathrm{arg}\left[\frac{V_{ud}V_{ub}^{}}{V_{cd}V_{cb}^{}}\right].$$ (1) Here we consider the possibility of extracting $`\gamma `$ from the $`B_s`$ decays to the final states $`D_s^\pm K^{}`$ (or the analogous $`D{}_{s}{}^{()}K_{}^{()}`$ combinations). The fact that $`D_s^{}K^+`$ and $`D_s^+K^{}`$ can be reached from both the $`B_s`$ and its CP conjugate $`\overline{B}_s`$ already has been exploited in a proposal to extract $`\gamma `$ from a time-dependent study . The two transition amplitudes have similar magnitudes, $`\|A(B_sD_s^{}K^+)\|\|A(\overline{B}_sD_s^{}K^+)\|\lambda ^3`$, where $`\lambda =\mathrm{sin}\theta _C0.22`$ is the small parameter which controls the hierarchy of the CKM matrix. Hence triangles built from these amplitudes need not suffer from being “squashed”. Given sufficient statistics, the time-dependent analysis eventually should yield $`\mathrm{sin}\gamma `$ at some level of accuracy. By contrast, our proposal allows one to measure $`\mathrm{sin}\gamma `$ using branching ratios only, with no need to determine the time at which the decay occurs. What will be necessary, instead, is to measure not only flavor-tagged but also CP-tagged $`B_s`$ decays. We begin by defining amplitudes for $`B_s`$ and $`\overline{B}_s`$ decay to the final state $`D_s^{}K^+`$, $`A_1`$ $`=`$ $`A(B_sD_s^{}K^+)=a_1e^{i\delta _1},`$ (2) $`A_2`$ $`=`$ $`A(\overline{B}_sD_s^{}K^+)=a_2e^{i\gamma }e^{i\delta _2}.`$ (3) The amplitude $`A_1`$ arises from the quark transition $`\overline{b}\overline{c}u\overline{s}`$ and is real (in the Wolfenstein parameterization), while $`A_2`$ arises from $`bu\overline{c}s`$ and carries the relative weak phase $`e^{i\gamma }`$. There are no penguin contributions. The amplitudes $`A_i`$ also have strong phases $`e^{i\delta _i}`$. The CP conjugated amplitudes are given by $`\overline{A}_1`$ $`=`$ $`A(\overline{B}_sD_s^+K^{})=a_1e^{i\delta _1}e^{2i\xi },`$ (4) $`\overline{A}_2`$ $`=`$ $`A(B_sD_s^+K^{})=a_2e^{i\gamma }e^{i\delta _2}e^{2i\xi },`$ (5) where the phase $`\xi `$ depends on the convention for CP transformations of the $`B_s`$ states, $$\mathrm{CP}\|B_s=e^{2i\xi }\|\overline{B}_s,\mathrm{CP}\|\overline{B}_s=e^{2i\xi }\|B_s.$$ (6) Any physical observable must be independent of $`\xi `$. We also define a set of amplitudes for the CP eigenstates of the $`B_s`$ meson, $$\|B_s^\pm =\frac{1}{\sqrt{2}}\left[\|B_s\pm e^{2i\xi }\|\overline{B}_s\right],$$ (7) to decay into the same $`D_sK`$ final states, $`A_\pm `$ $`=`$ $`A(B_s^\pm D_s^{}K^+),`$ (8) $`\overline{A}_\pm `$ $`=`$ $`A(B_s^\pm D_s^+K^{}).`$ (9) These amplitudes satisfy simple triangle relations, $`\sqrt{2}A_\pm `$ $`=`$ $`A_1\pm e^{2i\xi }A_2=(a_1\pm a_2e^{i\gamma }e^{i\delta })e^{i\delta _1},`$ (10) $`\sqrt{2}\overline{A}_\pm `$ $`=`$ $`\overline{A}_2\pm e^{2i\xi }\overline{A}_1=\pm (a_1\pm a_2e^{i\gamma }e^{i\delta })e^{i\delta _1},`$ (11) where $`\delta =\delta _2\delta _1+2\xi `$ is the convention-independent (and observable) strong phase difference. It is clear that any construction which is insensitive to $`\delta `$ will also be insensitive to the unphysical phase $`\xi `$. It is also clear that by changing $`\xi `$, it is possible to take $`B_s^\pm `$ to be a linear combination of $`B_s`$ and $`\overline{B}_s`$ with any relative phase. We will derive a relation for $`\mathrm{sin}2\gamma `$ involving the magnitudes of the amplitudes $`A_i`$ and $`\overline{A}_i`$. From the freedom to choose $`\xi `$ in Eq. (10), it is clear that the CP-even and CP-odd amplitudes will yield triangle relations which contain the same information about $`\gamma `$. For any CP eigenstate $`B_s^{\mathrm{CP}}`$, then, it is possible to choose $`\xi `$ so that $`A_{\mathrm{CP}}`$ $`=`$ $`A(B_s^{\mathrm{CP}}D_s^{}K^+)=(A_1+A_2)/\sqrt{2},`$ (12) $`\overline{A}_{\mathrm{CP}}`$ $`=`$ $`A(B_s^{\mathrm{CP}}D_s^+K^{})=(\overline{A}_1+\overline{A}_2)/\sqrt{2}.`$ (13) Without loss of generality, we also choose a convention in which $`\delta _1=0`$, in which case the triangle relations are very simple. They are illustrated graphically in Fig. 1, where the amplitudes may be interpreted as vectors in the complex plane. As drawn, the angle between $`A_2`$ and $`\overline{A}_2`$ is $`2\gamma `$. For an analytical solution, it is convenient to define $`\alpha `$ $`=`$ $`{\displaystyle \frac{2\|A_{\mathrm{CP}}\|^2\|A_1\|^2\|A_2\|^2}{2\|A_1\|\|A_2\|}},`$ (14) $`\overline{\alpha }`$ $`=`$ $`{\displaystyle \frac{2\|\overline{A}_{\mathrm{CP}}\|^2\|\overline{A}_1\|^2\|\overline{A}_2\|^2}{2\|\overline{A}_1\|\|\overline{A}_2\|}},`$ (15) in terms of which we find $$\mathrm{sin}2\gamma =\pm \left(\alpha \sqrt{1\overline{\alpha }^2}\overline{\alpha }\sqrt{1\alpha ^2}\right).$$ (16) The determination of $`\gamma `$ itself then has an eightfold ambiguity. This construction is quite analogous to that of Ref. , in which the extraction of $`\mathrm{sin}\gamma `$ from $`B_d`$ decay to a CP eigenstate $`D_{\mathrm{CP}}`$ was studied. (However, the triangles are “squashed” in that analysis.) Note that if $`\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s`$ is significant, the squared amplitudes $`\|A_i\|^2`$ and $`\|\overline{A}_i\|^2`$ are proportional to partial rates, e.g., $$\|A_1\|^2\mathrm{\Gamma }(B_sD_s^{}K^+),$$ (17) rather than to branching ratios. The measurement of $`\|A_{\mathrm{CP}}\|`$ requires that one tag the initial $`B_s`$ state as a CP eigenstate $`B_s^{\mathrm{CP}}`$. This is possible if the $`B_s`$ is produced in the decay $`\mathrm{{\rm Y}}(5S)B_s\overline{B}_s`$. Since the $`\mathrm{{\rm Y}}(5S)`$ is a CP even state and the $`B_s`$ and $`\overline{B}_s`$ are emitted in a relative $`p`$-wave, the CP eigenvalues of the $`B_s/\overline{B}_s`$ mixtures are anti-correlated. Hence if the “tagging” $`B_s`$ decays to a CP eigenstate such as $`D_s^+D_s^{}`$, then the other $`B_s`$ is constrained to be a CP eigenstate as well. It is crucial that the tagging decay be one in which direct CP violation is expected to be small. Tagging modes with spin one particles, such as $`B_s^{\mathrm{CP}}\psi \varphi `$ and $`D_s^{}\overline{D}_s^{}`$, can be used if an angular analysis is performed to select a final state of definite CP. So long as the amplitudes $`\|A_{\mathrm{CP}}\|`$ and $`\|\overline{A}_{\mathrm{CP}}\|`$ are measured with the same CP tagging mode on the opposite side, it is unimportant whether a CP even or CP odd tag is employed. This simple method of CP tagging relies on the Standard Model expectation that CP violation in $`B_s`$ mixing is not significant. As $`B_s`$ mixing is generated dominantly by $`tW`$ box diagrams, CP violating effects are proportional to $`\mathrm{sin}2\beta _s`$, where $$\beta _s=\mathrm{arg}\left[\frac{V_{ts}V_{tb}^{}}{V_{cs}V_{cb}^{}}\right]\lambda ^2$$ (18) is small. Furthermore, we assume that CP violation in $`B_s`$ mixing will have been constrained experimentally by the time the analysis proposed here is performed. The method is not appropriate to a hadronic production environment such as the Tevatron or the LHC, since in this case the $`B_s`$ and $`\overline{B}_s`$ do not arise from an initial CP eigenstate. Nor, of course, can this analysis be performed at the $`B`$ Factories as presently configured to operate at the $`\mathrm{{\rm Y}}(4S)`$. To our knowledge, this is the first proposal for a clean measurement of a CKM phase which is unique to an $`e^+e^{}`$ collider operating at the $`\mathrm{{\rm Y}}(5S)`$. We now make a crude estimate of the number of $`B_s\overline{B}_s`$ pairs required to measure $`\mathrm{sin}\gamma `$ with a precision of $`0.1`$, for which approximately $`10^2`$ reconstructed events would be needed. To be concrete, we take the tagging mode $`B_s^{\mathrm{CP}}D_s^+D_s^{}`$. With order-of-magnitude estimates of the relevant branching ratios, $`(B_s^{\mathrm{CP}}D_s^+D_s^{})10^2`$ and $`(B_s^{\mathrm{CP}}D_sK)2\times 10^4`$, and assuming that the $`D_s`$ can be reconstructed efficiently by combining a number of decay modes, we find a combined CP-tagged branching fraction $`_{\mathrm{tot}}10^6`$. Hence approximately $`10^8`$ $`B_s\overline{B}_s`$ events would be needed for this measurement. The decays of the $`\mathrm{{\rm Y}}(5S)`$ to $`B_s`$ flavored mesons produce primarily the combination $`B{}_{s}{}^{}\overline{B}_s^{}`$, as well as $`B{}_{s}{}^{}\overline{B}_{s}^{}`$, $`B_s\overline{B}_s^{}`$ and $`B_s\overline{B}_s`$. The relative rates have been computed in a variety of models, yielding the estimates $$\sigma (B_s\overline{B}_s)/\sigma (B{}_{s}{}^{}\overline{B}{}_{s}{}^{})0.10.2$$ (19) and $$\sigma (B{}_{s}{}^{}\overline{B}_{s}^{}+B_s\overline{B}{}_{s}{}^{})/\sigma (B{}_{s}{}^{}\overline{B}{}_{s}{}^{})0.050.5.$$ (20) A $`B_s^{}`$ produced in this way decays to a $`B_s`$ and a very soft photon, so the other combinations will also be seen as $`B_s\overline{B}_s`$. In fully reconstructed $`B_s`$ decays the combinations can be separated by measuring the boost of the $`B_s`$ . From the ratio $`\sigma (B{}_{s}{}^{}\overline{B}{}_{s}{}^{})/\sigma (\mathrm{{\rm Y}}(4S))0.1`$ of production cross sections and the fact that an $`e^+e^{}`$ collider with luminosity $`=10^{33}\mathrm{cm}^2\mathrm{sec}^1`$ produces $`3.6\times 10^7`$ $`\mathrm{{\rm Y}}(4S)`$ events per year , we see that $`10^2`$ CP-tagged $`B_s\overline{B}_s`$ events, decaying in this mode, are within the reach of a $`B`$ Factory upgraded to operate at $`10^{35}\mathrm{cm}^2\mathrm{sec}^1`$. Since it is not necessary to measure the time-dependence of the decay, the experiment could be performed in a future high luminosity run of the Cornell $`e^+e^{}`$ storage ring. Moreover, there are ways to increase substantially the number of usable events. First, one may repeat the analysis with the final states $`D_s^{}K`$, $`D_sK^{}`$ and $`D_s^{}K^{}`$. Note that no angular analysis is necessary here, since one is not isolating a CP eigenstate on the decay side. Second, one can add additional CP-tagging modes such as $`B_s^{\mathrm{CP}}\psi \varphi `$ or $`D_s^{}\overline{D}_s^{}`$. Although in this case an angular analysis would be required to separate final states of definite CP, studies in the $`B_d`$ system indicate that this can be done without a large cost in tagging efficiency . This gain in efficiency will be offset in part by the cost of fully reconstructing the $`D_s`$ states, a penalty which we have not explicitly included. Finally, it also may be possible to use the $`B{}_{s}{}^{}\overline{B}_{s}^{}`$ and $`B_s\overline{B}_s^{}`$ combinations for CP tagging. Parity conservation requires that the pair be produced in a relative $`p`$-wave. Therefore the initial state is of the form $$\frac{1}{\sqrt{2}}\left[B_s^+B_s^{}+B_s^{}B_s^+\right],$$ (21) where $`B_s^\pm `$ are the CP eigenstate mixtures of $`B_s^{}`$ and $`\overline{B}_s^{}`$, in analogy with the $`B_s`$ combinations (7). After the transition $`B{}_{s}{}^{}B_s\gamma `$, in which the magnetic photon carries $`\mathrm{CP}=1`$, the CP eigenvalues of the $`B_s/\overline{B}_s`$ mixtures on the two sides are correlated (rather than anti-correlated, as in direct $`\mathrm{{\rm Y}}(5S)B_s\overline{B}_s`$). As we have shown, our analysis is equivalent for correlated and anti-correlated states. Unfortunately, it is not possible to CP tag using the dominant $`B{}_{s}{}^{}\overline{B}_s^{}`$ combination, since the total spin quantum number is not fixed at production. Taken together, the use of additional modes on the tagging and decay sides and of the $`B{}_{s}{}^{}\overline{B}_{s}^{}`$ and $`B_s\overline{B}_s^{}`$ initial states should allow one to relax considerably the luminosity requirement estimated above. Alternatively, for a given integrated luminosity these enhancements would allow a statistically more precise measurement of $`\mathrm{sin}\gamma `$. Generally speaking, we believe that our proposal is feasible within many of the scenarios under discussion for future luminosity upgrades of the $`B`$ Factories now operating at the $`\mathrm{{\rm Y}}(4S)`$. As shown above, the measurement of $`\gamma `$ from this analysis is insensitive to the strong phase difference $`\delta `$ between the amplitudes $`A_1`$ and $`A_2`$. In fact, $`\delta `$ could be extracted simultaneously with $`\gamma `$ from the amplitude triangles shown in Fig. 1. Nevertheless, it is useful to have some idea of whether $`\delta `$ may be expected to be large. It is clear that elastic rescattering of the $`D_sK`$ final state will be the same for $`B_s`$ and $`\overline{B}_s`$ transitions and will not lead to nonzero $`\delta `$. Instead, what is needed to generate $`\delta 0`$ is rescattering through an intermediate state $`f`$ which is produced differently by $`B_s`$ and $`\overline{B}_s`$. Then we may write $`{\displaystyle \frac{A(B_sD_s^{}K^+)}{A(\overline{B}_sD_s^{}K^+)}}`$ $`=`$ $`{\displaystyle \frac{A_1^{\mathrm{dir}}+ϵ_fA_1^f+\mathrm{}}{A_2^{\mathrm{dir}}+ϵ_fA_2^f+\mathrm{}}}`$ (22) $``$ $`{\displaystyle \frac{A_1^{\mathrm{dir}}}{A_2^{\mathrm{dir}}}}{\displaystyle \frac{1+ϵ_fA_1^f/A_1^{\mathrm{dir}}}{1+ϵ_fA_2^f/A_2^{\mathrm{dir}}}},`$ (23) where $`A_i^{\mathrm{dir}}`$ are the amplitudes for the direct production of $`D_s^{}K^+`$, $`A_i^f`$ are the amplitudes for the production of the intermediate state $`f`$, and $`ϵ_f`$ is the amplitude for the rescattering $`fD_s^{}K^+`$. While $`A_1^{\mathrm{dir}}`$ and $`A_2^{\mathrm{dir}}`$ have the same strong phase, a strong phase difference can be generated if $`A_1^f/A_1^{\mathrm{dir}}A_2^f/A_2^{\mathrm{dir}}`$. For a similar discussion in the context of $`D`$ decays, see Ref. . To estimate the size of $`\delta `$ which could be generated, we consider a model in which $`f`$ is a two body intermediate state. We employ the formalism of Ref. , based on Regge phenomenology and naive factorization of the $`B_s`$ decay matrix elements. With $`f=D_s^{}K^{}`$ and rescattering to $`D_sK`$ via exchange of the Pomeron and $`\varphi `$ trajectories, we find $$\delta <5^{}.$$ (24) Although our estimate is extremely model-dependent, it does provide some evidence that this mechanism is unlikely to produce a large value of $`\delta `$. We note that a perturbative factorization formalism is not applicable to the decay $`B_sD_sK`$. Finally, we turn to the issue of CP tagging in $`B_d`$ decays. Here the situation is complicated by the fact that the Standard Model predicts large CP violation in $`B_d`$ mixing. Therefore a state which is tagged at time $`t=0`$ as being in a CP eigenstate will evolve by time $`t`$ into a linear combination $`B_d^\pm (t)`$ of the CP even ($`B_d^+`$) and CP odd ($`B_d^{}`$) states. The evolution is given by $$B_d^\pm (t)=e^{i(M_B+\mathrm{\Gamma }/2)t}\left(a_\pm (t)B_d^\pm +b_\pm (t)B_d^{}\right),$$ (25) where $`a_\pm (t)`$ $`=`$ $`\mathrm{cos}\left(\mathrm{\Delta }m_dt/2\right)\pm i\mathrm{cos}2\beta \mathrm{sin}\left(\mathrm{\Delta }m_dt/2\right),`$ (26) $`b_\pm (t)`$ $`=`$ $`\mathrm{sin}2\beta \mathrm{sin}\left(\mathrm{\Delta }m_dt/2\right),`$ (27) and $`\mathrm{\Delta }m_d`$ is the mass splitting between $`B_H`$ and $`B_L`$. Here $$\beta =\mathrm{arg}\left[\frac{V_{td}V_{tb}^{}}{V_{cd}V_{cb}^{}}\right]$$ (28) is an angle which is expected to be large. (If $`B_d`$ mixing receives a significant contribution from new physics, then the CP violating phase “$`\mathrm{sin}2\beta `$” extracted from the asymmetry in $`B_dJ/\psi K_S`$ is actually what governs the time evolution (26).) Note that in the CP conserving limit $`\mathrm{sin}2\beta 0`$, we have $`a_\pm (t)\mathrm{exp}(\pm i\mathrm{\Delta }m_dt/2)`$ and $`b_\pm (t)0`$, so the masses of the CP eigenstates are shifted to $`M_B\pm \frac{1}{2}\mathrm{\Delta }m_d`$ but the states do not mix. In analogy with the $`B_s`$ case, we define amplitudes for a $`B_d`$, tagged at $`t=0`$ as a CP eigenstate $`B_d^\pm `$, to decay into the final states $`D^\pm \pi ^{}`$ (or $`D^{}\pi ^{}`$ or $`D^{()}\rho ^{}`$) at time $`t`$, $`A_\pm (t)`$ $`=`$ $`A(B_d^\pm (t)D^{}\pi ^+),`$ (29) $`\overline{A}_\pm (t)`$ $`=`$ $`A(B_d^\pm (t)D^+\pi ^{}).`$ (30) The triangle relations analogous to Eq. (12) then take a form which depends on $`t`$, $`\sqrt{2}\left\|A_\pm (t)\right\|`$ $`=`$ $`\left\|r_{}(t)A_1+r_\pm (t)A_2\right\|,`$ (31) $`\sqrt{2}\left\|\overline{A}_\pm (t)\right\|`$ $`=`$ $`\left\|r_\pm (t)\overline{A}_1+r_{}(t)\overline{A}_2\right\|,`$ (32) where here the amplitudes $`A_i`$ and $`\overline{A}_i`$ are defined as in Eqs. (2) and (4) but for $`B_dD\pi `$, and $$r_\pm (t)=\left[1\pm \mathrm{sin}2\beta \mathrm{sin}\mathrm{\Delta }m_dt\right]^{1/2}.$$ (33) One may extract $`\gamma `$ by fixing a value of $`t`$ and then constructing the amplitude triangle with the sides scaled by $`r_\pm (t)`$ as in Eq. (31). The expressions for $`\alpha `$ and $`\overline{\alpha }`$ are modified to $`\alpha _\pm (t)`$ $`=`$ $`{\displaystyle \frac{2\|A_\pm (t)\|^2r_{}^2(t)\|A_1\|^2r_\pm ^2(t)\|A_2\|^2}{2r_+(t)r_{}(t)\|A_1\|\|A_2\|}},`$ (34) $`\overline{\alpha }_\pm (t)`$ $`=`$ $`{\displaystyle \frac{2\|\overline{A}_\pm (t)\|^2r_\pm ^2(t)\|\overline{A}_1\|^2r_{}^2(t)\|\overline{A}_2\|^2}{2r_+(t)r_{}(t)\|\overline{A}_1\|\|\overline{A}_2\|}}.`$ (35) The solution (16) for $`\mathrm{sin}2\gamma `$, written in terms of $`\alpha _\pm (t)`$ and $`\overline{\alpha }_\pm (t)`$, is independent of $`t`$ by construction. Note that for this time-dependent analysis, the decays of the CP even and CP odd eigenstates are not equivalent. The procedure may be repeated to give an independent measurement of $`\gamma `$ for each bin in $`t`$. Writing the amplitude triangles in the form (31), a cosmetic change which makes obvious the generalization of $`\alpha `$ and $`\overline{\alpha }`$ to $`\alpha _\pm (t)`$ and $`\overline{\alpha }_\pm (t)`$, requires a $`t`$-dependent choice of the CP transformation phase $`\xi `$. This is legitimate, since in Eq. (31) one combines amplitudes only at a fixed value of $`t`$. The necessity of determining the decay time $`t`$ means that such a measurement of $`\gamma `$ in the $`B_d`$ system would have to be performed at an asymmetric $`B`$ Factory operating at the $`\mathrm{{\rm Y}}(4S)`$. Although in principle the analysis could be performed by the BaBar or BELLE Collaborations, there are several difficulties. First, an accurate independent determination of $`\mathrm{sin}2\beta `$ must be available. Second, it is necessary to collect sufficient statistics to construct the amplitude triangles for individual bins in $`t`$. Third, in the case of $`B_dD^\pm \pi ^{}`$, $`D^{}\pi ^{}`$ or $`D^{()}\rho ^{}`$ the amplitude triangles are “squashed”, with one side shorter than the other two by a factor of order $`\lambda ^2`$. (They would not be squashed, however, if the analysis were performed instead for a mode such as $`B_dD^{()}K_S`$.) In summary, we have presented a new approach to extracting the CKM angle $`\gamma `$, employing an analysis which depends on tagging an initial $`B_s`$ or $`B_d`$ as a CP eigenstate. This theoretically clean method is free from dependence on unknown strong phases. In the $`B_s`$ case, the analysis is unique to an experiment performed at a very high luminosity $`e^+e^{}`$ collider operating at the $`\mathrm{{\rm Y}}(5S)`$ resonance. While there are no definite plans to upgrade any of the existing symmetric or asymmetric $`B`$ Factories to operate in this mode, we hope that the proposal outlined here will help rekindle interest in this possibility. It is a pleasure to thank Yossi Nir and Helen Quinn for helpful correspondence. Support was provided by the National Science Foundation under Grant PHY–9404057, by the Department of Energy under Outstanding Junior Investigator Award DE–FG02–94ER40869, and by the Research Corporation under Cottrell Scholar Award CS0362.
warning/0003/math0003212.html	ar5iv	text	# Motivic zeta functions of infinite dimensional Lie algebras ## 1. Introduction ### 1.1. In the present paper we associate motivic zeta functions to certain classes of infinite dimensional Lie algebra over a field $`k`$ of characteristic zero. Included in these classes are the important cases of loop algebras, affine Kac-Moody algebras, the Virasoro algebra and Lie algebras of Cartan type. These zeta functions take their values in the Grothendieck ring of algebraic varieties over $`k`$ and are built by encoding in some manner the $`k`$-subalgebras of a given codimension. This construction is done by an adaptation of the idea of the motivic Igusa zeta function recently introduced in . Let us recall the definition of the Grothendieck ring $``$ of algebraic varieties over $`k`$. This is the ring generated by symbols $`[S],`$ for each $`S`$ an algebraic variety over $`k`$, that is, a reduced separated scheme of finite type over $`k`$, with the relations 1. $`[S]=[S^{}]`$ if $`S`$ is isomorphic to $`S^{}`$; 2. $`[S]=[S\backslash S^{}]+[S^{}]`$ if $`S^{}`$ is closed in $`S`$; 3. $`[S\times S^{}]=[S][S^{}].`$ The idea leading to this construction comes from the analogy with the zeta functions capturing the subalgebra lattice of $`p`$-adic Lie algebras. More precisely, in the zeta function of a finite dimensional Lie algebra $`L_p`$ over $`_p`$ was defined as $`\zeta _{L_p}(s)`$ $`:=`$ $`{\displaystyle \underset{HL_p}{}}\|L_p:H\|^s`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}a_{p^n}(L_p)p^{ns},`$ where the first sum is taken over finite index $`_p`$-subalgebras, that is $`_p`$-submodules which define subalgebras in $`L_p`$ and $`a_{p^n}(L_p)`$ is the number of these of index $`p^n.`$ This is a well-defined series since there are only finitely many of each for a given finite index. Note that since $``$ is dense in $`_p,`$ finite index $`_p`$-subalgebras are the same as finite index $``$-subalgebras. The zeta function captures in some analytic manner the lattice of subalgebras. It was proved in the same paper that, when $`L_p`$ is additively isomorphic to $`_p^d`$, the zeta function is a rational function in $`p^s.`$ The proof depends on Denef’s results on $`p`$-adic definable integrals. Explicit formulae in terms of an associated resolution of singularities are given in for $`\zeta _{L_p}(s)`$ in the case that $`L_p=L_p`$ for some Lie algebra $`L`$ defined over $``$. In general an infinite dimensional Lie algebra over a field $`k`$ of characteristic zero contains an infinite number of subalgebras of each given finite codimension. So it doesn’t make sense to count them in a conventional way and try to encode this counting function in a Dirichlet or Poincaré series with integer coefficients. Instead, in the cases we consider in the present paper, the set of subalgebras of a fixed finite codimension form a constructible subset of some Grassmannian, so we may asssociate to it an element of the Grothendieck ring $``$, and we can construct the zeta function as a Poincaré series with coefficients in $``$. ### 1.2. An example : counting $`k[[t]]`$-submodules of $`k[[t]]^2`$ To illustrate the thinking behind this, it is instructive to see how we would count $`k[[t]]`$-submodules of finite codimension in the free $`k[[t]]`$-module $`L=k[[t]]^2.`$ Let $`e_1,e_2`$ be a basis for $`L`$, then each subalgebra $`H`$ of codimension $`1`$ can be represented by coordinates in a $`k[[t]]`$-basis with respect to this standard basis. We can record these coordinates in a matrix whose rows are the coordinates for the basis of the subalgebra $`H`$. By normalizing the matrix into upper triangular form we can choose a unique representative $$\left(\begin{array}{cc}m_{11}\hfill & m_{12}\hfill \\ 0\hfill & m_{22}\hfill \end{array}\right)$$ for each $`H.`$ We have two cases: 1. $`m_{22}t\left(k[[t]]\right)^{},`$ and $`m_{11}\left(k[[t]]\right)^{}.`$ To get a unique representative we can assume that $`m_{11}=1`$ and $`m_{22}=t.`$ By subtracting $`k[[t]]`$-combinations of the second row from the first we can choose $`m_{12}k.`$ So the subalgebras of this sort give a space of subalgebras which looks like $`k.`$ 2. If $`m_{22}\left(k[[t]]\right)^{},`$ and $`m_{11}t\left(k[[t]]\right)^{}`$ then we get a unique representative of the form $$\left(\begin{array}{cc}t\hfill & 0\hfill \\ 0\hfill & 1\hfill \end{array}\right).$$ This gives just one subalgebra of codimension $`1.`$ In the case of $`(_p^2,+)`$ this analysis showed us that there were $`p+1`$ subalgebras of index $`p.`$ Here we would like to say that there are $`𝔸_k^1+1`$ subalgebras, with $`𝔸_k^1`$ the affine line over $`k`$, since, for any field $`K`$ containing $`k`$, the $`K[[t]]`$-submodules of $`K[[t]]^2`$ are parametrized by $`K+1=𝔸_k^1(K)+1`$. For that to make sense one is naturally led to work in the Grothendieck ring $``$. So if we set $`𝐋:=[𝔸_k^1]`$ in $``$, we can interpret our analysis of the subalgebras in $`(\left(k[[t]]\right)^2,+)`$ as saying there are $`𝐋+1`$ subalgebras of codimension $`1`$. On the other hand if we consider counting all vector subspaces of codimension 1, i.e. ignoring the action of $`t,`$ then there are really too many subspaces to be able to assign an element of the Grothendieck ring $``$. In fact, this is similar to the situation when $`k=𝔽_p`$. There are $`1+p+\mathrm{}+p^n`$ $`𝔽_p[[t]]`$-submodules of index $`p^{n+1}`$ in $`\left(𝔽_p[[t]]\right)^2`$ but an infinite number of additive $`𝔽_p`$-subspaces of index $`p^{n+1}.`$ ### 1.3. Stability and the motivic zeta function The key to being able to assign an element of the Grothendieck ring to subalgebras of a fixed codimension is a notion we call stability. The context in which we shall work is that of $``$-filtered or $``$-filtered Lie algebras over a field $`k`$ of characteristic zero. These are infinite dimensional Lie algebras $`L`$ with a filtration $`L=L_i`$ indexed by $``$ or $``$ consisting of subalgebras $`L_i`$ of finite codimension in $`L_0.`$ By a class $`𝒳`$ of subalgebras of $`L`$, we shall mean the data, for every field $`K`$ which is a finite extension of $`k`$, consisting of a family $`𝒳(K)`$ of subalgebras of $`LK`$. We will sometimes, with a slight abuse, identify $`𝒳`$ with the union of the various sets $`𝒳(K)`$. Let $`A_n(𝒳)(K)`$ denote the set of subalgebras in $`𝒳(K)`$ of codimension $`n`$ in $`LK`$. We say that a class of subalgebras $`𝒳`$ of $`L`$ is stable if, for every $`n`$, there exists some $`f(n)`$ such that, for every $`K`$, if $`H`$ belongs to $`A_n(𝒳)(K)`$, then $`H`$ contains $`L_{f(n)}K`$. In this case it is possible to show for certain classes of subalgebras $`𝒳`$ that $`A_n(𝒳)`$ is a constructible subset of the Grassmannian $`\mathrm{Gr}(L/L_{f(n)}).`$ We can then define the concept of a motivic zeta function: $$P_{L,𝒳}(T):=\underset{n=0}{\overset{\mathrm{}}{}}[A_n(𝒳)]T^n$$ where the coefficients $`[A_n(𝒳)]`$ belong to $`.`$ This zeta function will be our interpretation of how to add the series $$\underset{H𝒳}{}T^{\mathrm{codim}H},$$ which is then an element of $`[[T]].`$ We also introduce a variant of this function, which has not been considered previously in the $`p`$-adic case, which counts subalgebras commensurable with the subalgebra $`L_0`$ of infinite codimension in a $``$-filtered Lie algebra $`L`$. A subalgebra $`H`$ is commensurable with $`L_0`$ if $`HL_0`$ has finite codimension in $`H`$ and $`L_0`$. We then define the codimension of $`H`$ in $`L_0`$ to be the sum of the codimension of $`HL_0`$ in $`H`$ and $`L_0`$. Our approach to proving the stability of subalgebras in a filtered Lie algebra is a new concept we call well-covered. This means essentially that elements in the $`n`$-th term of the grading can be realised in many ways as commutators of elements further up the grading. We prove that affine Kac-Moody Lie algebras (which are built from loop algebras $`Lk[[t]]`$, with $`L`$ a finite dimensional Lie algebra over $`k`$), the Virasoro algebra and Lie algebras of Cartan type are well-covered and hence have an associated motivic zeta function. On the other hand we show that a finitely generated free Lie algebra is not well-covered since an Engel element like $`(a,b,b,\mathrm{},b)`$ of length $`n`$ cannot be realised as a commutator in more than one way. Hence we can show that such Lie algebras don’t have an associated motivic zeta function. ### 1.4. Rationality of the motivic zeta function Once the motivic zeta function is defined, it is natural to consider the question of its rationality. By using results from and , we are able to prove the rationality of motivic zeta functions counting $`k[[t]]`$-subalgebras in $`k[[t]]`$-Lie algebras of the form $`L_kk[[t]]`$ where $`L_k`$ is a finite dimensional Lie algebra over $`k`$, a field of characteristic zero. Denote by $`_{\mathrm{loc}}`$ the ring $`[𝐋^1]`$ obtained by localization and define by $`[T]_{\mathrm{loc}}`$ the subring of $`_{\mathrm{loc}}[[T]]`$ generated by $`_{\mathrm{loc}}[T]`$ and the series $`(1𝐋^aT^b)^1`$ with $`a`$ and $`b.`$ ###### 1.5 Theorem. Let $`k`$ be a field of characteristic zero. 1. If $`L`$ is a finite dimensional free $`k[[t]]`$-Lie algebra of the form $`L=L_k_kk[[t]]`$, with $`L_k`$ a Lie algebra over $`k`$, and $`𝒳(K)`$ is the set of all $`K[[t]]`$-subalgebras of $`LK`$, then $`P_{L,𝒳}(T)`$ is well-defined and is rational, belonging to $`[T]_{\mathrm{loc}}`$. 2. If $`L`$ is a finite dimensional $`k((t))`$-Lie algebra and $`L_0`$ is a choice of some $`k[[t]]`$-Lie subalgebra of the form $`L_0=L_k_kk[[t]]`$, with $`L_k`$ a Lie algebra over $`k`$, and $`𝒳(K)`$ is the set of all $`K[[t]]`$-subalgebras of $`LK`$ commensurable with $`L_0K`$, then $`P_{L_0,𝒳}(T)`$ is well-defined and is rational, belonging to $`[T]_{\mathrm{loc}}`$ . The key to the proof of this Theorem is to express the motivic zeta function as a motivic integral as developed by Denef and the second author and . In fact we are naturally led to introduce the concept of a motivic measure on the infinite Grassmannian. ### 1.6. Implications for $`p`$-adic zeta functions If $`L`$ is a Lie algebra over $`,`$ then the motivic zeta function of $`L[[t]]`$ has implications for the zeta functions of $`L_p`$ as we range over primes $`p.`$ In particular we can show, by “taking the trace of Frobenius” of the motivic zeta function that the explicit formulae calculated in and alluded to in 1.1 for $`\zeta _{L_p}(s)`$ are canonical and hence independent of a choice of resolution and simplicial decomposition of a cone which are involved in the calculation. We also define the concept of the topological zeta function of $`L`$ which is the analogue of the topological Igusa zeta function defined in . This can be thought of as the limit as $`p1`$ of the explicit formulae expressing $`\zeta _{L_p}(s).`$ To show that this is well-defined we show that this is the same as taking the Euler characteristic of the associated motivic zeta function, similarly as in . Acknowledgements. We should like to thank the London Mathematical Society and the Institut Mathématique de Jussieu for grants which facilitated our collaboration. The first author would also like to thank the Royal Society for funding in the form of a University Research Fellowship. Notations and terminology. We shall use the notation $`=\{0,1,2,\mathrm{}\}.`$ For $`k`$ a field, we denote by $`𝔸_k^n`$ the affine space $`\mathrm{Spec}k[x_1,\mathrm{},x_n]`$. By an algebraic variety over $`k`$, we shall always mean a scheme of finite type over $`k`$, which is separated and reduced. Except explicitely stated otherwise, all Lie algebras we consider will be assumed to be Lie algebras over $`k`$. We shall consider a measure $`\stackrel{~}{\mu }`$ taking values in $`_{\mathrm{loc}}`$ and a measure $`\mu `$ taking values in the completion $`\widehat{}`$ of $`_{\mathrm{loc}}`$. ## 2. Definition of the motivic zeta function We fix a field $`k`$. In this section we explore some infinite dimensional Lie algebras for which one can define a motivic zeta function. ###### 2.1 Definition. (1) An $``$-filtered Lie algebra $`L`$ is a Lie algebra equipped with a filtration by subspaces $`L_j`$, $`j`$, of finite codimension with $`L_{j+1}L_j`$ satisfying $`(L_j,L_k)L_{j+k}`$ and $`_iL_i=0.`$ We shall assume for convenience that $`L_0=L.`$ (2) Similarly a $``$-filtered Lie algebra comes with a filtration by subspaces $`L_j`$, $`j`$, with $`L_{j+1}L_j`$ satisfying $`(L_j,L_k)L_{j+k}`$ where we assume that $`L_j`$ is commensurable with $`L_0`$ and $`L=_jL_j`$ and $`_iL_i=0.`$ The filtration defines a (formal) topology on $`L`$ defined by the neighbourhood base of zero consisting of subspaces $`L_j.`$ A subspace $`H`$ is commensurable with $`L_0`$ if the intersection $`HL_0`$ has finite codimension in both $`H`$ and $`L_0`$. We define the codimension of $`H`$ in $`L_0`$ in this case as the sum of the codimension of $`HL_0`$ in $`H`$ and $`L_0.`$ Note that in an $``$-filtered Lie algebra the $`L_j`$ are ideals but this is not necessarily the case in a $``$-filtered Lie algebra since $`(L_j,L_k)L_{kj}.`$ ###### 2.2 Remark. We can carry out the analysis of this section just under the assumption that the Lie algebra $`L`$ comes only with a formal topology of subspaces $`L_j`$ of finite codimension without insisting that $`(L_j,L_k)L_{j+k}`$. However essentially all the cases where we can get the theory up and running apply to the examples where the chain of subspaces is actually a filtration. However in any particular case it may not be necessary to assume anything this strong. It may be possible to show that subalgebras are a constructible subset of the finite Grassmannian $`\mathrm{Gr}(L/L_j).`$ ### 2.3. Examples of filtered Lie algebras (1) A good example of an $``$-filtration on $`L`$ is the lower central series. In this case $`L_0=L_1`$ and $`L_i=\gamma _i(L)=(\gamma _{i1}(L),L).`$ We are insisting that the filtration be of subalgebras of finite codimension. If the first layer $`L_1/L_2`$ is finite dimensional then this implies the same for the other layers $`L_i/L_{i+1}`$. If $`(L,L)`$ has infinite codimension then the theory won’t get up and running, since there will be too many subalgebras of codimension 1. Note that, by definition of an $``$-filtration, $`L_i`$ contains $`\gamma _i(L_1).`$ (2) Another example arises if $`L`$ has some $``$-grading, so that $`L=_{j0}L(j)`$ and suppose that $`L(j)`$ is finite dimensional then put $`L_k=_{jk}L_j.`$ (3) If $`L`$ has the structure of a $`k[[t]]`$-module then put $`L_k=t^kL.`$ (4) If $`L`$ comes with a $``$-grading $`L=_jL(j)`$ then we can define $`L_k=_{jk}L(j)`$ for $`k`$ and consider subalgebras commensurable with $`L_0`$. (5) Finally if $`L`$ is a $`k((t))`$-Lie algebra with a choice of $`k[[t]]`$-subalgebra $`L_0`$ we take $`L_k=t^kL_0`$ for $`k`$. ### 2.4. Constructible sets of subalgebras in Grassmannians It is the subalgebras of finite codimension that are closed with respect to the formal topology that we shall seek to count. Such subalgebras contain then some term $`L_j`$ of the neighbourhood base and can therefore be viewed as a point of the finite Grassmannian $`\mathrm{Gr}(L/L_j).`$ This follows because if $`H`$ is closed, $`H=_i(H+L_i)`$ and since $`H`$ has finite codimension, the chain $`H+L_i`$ must stabilize at some point. In the case of a closed commensurable subalgebra $`H`$ in a $``$-filtered Lie algebra, there exists some $`j`$ with $`L_jHL_j.`$ Hence these will be points of the finite Grassmannian $`\mathrm{Gr}(L_j/L_j).`$ For convenience we set $`X_j=\mathrm{Gr}(L/L_j)`$ or $`\mathrm{Gr}(L_j/L_j)`$ according to which situation we find ourselves. ###### 2.5 Remark. In section 4 we shall introduce the concept of infinite Grassmannians which can be thought of as the direct limit of these finite Grassmannians. These are considered for example in Section 7.2 (i) of . The set of subalgebras $`𝒳`$ can then be thought of as subsets of these infinite Grassmannians. We shall see later that in some cases it is possible to put a motivic measure directly on an infinite Grassmannian and express the motivic zeta functions of this section as integrals over $`𝒳`$ with respect to this measure. ###### 2.6 Remark. Note that not all subalgebras of finite codimension will be closed with respect to a chosen formal topology. For example in $`_ike_i`$ with respect to the formal topology $`L_j=_{ij}ke_i,`$ the subspace $$e_0e_1,e_1e_2,\mathrm{},e_ke_{k+1},\mathrm{}$$ has codimension 1 but contains no term of the filtration $`L_j.`$ ###### 2.7 Remark. There is a one-to-one correspondence preserving the codimension between closed subalgebras of finite codimension of $`L`$ and closed subalgebras of finite codimension of the completion of $`L`$ with respect to the formal topology. ### 2.8. Let $`Y`$ be an algebraic variety of finite type over $`k`$. By the set underlying $`Y`$ we shall always mean the set of closed points of $`Y`$. These points correspond to rational points of $`Y`$ over fields $`K`$ which are finite extensions of $`k`$. The boolean algebra of constructible subsets of $`Y`$ is the smallest family of subsets of $`Y`$ containing all Zariski closed subsets of $`Y`$ and stable by taking finite unions and complements. It follows from the definitions that to any constructible subset $`A`$ of $`Y`$ one can associate a canonical element $`[A]`$ of the Grothendieck ring $``$. We shall now view the finite Grassmannian $`X_j=\mathrm{Gr}(L/L_j)`$ or $`\mathrm{Gr}(L_j/L_j)`$ as an algebraic variety of finite type over $`k`$. Hence points of $`X_j`$ will be linear spaces defined over some field $`K`$ which is a finite extension of $`k`$. By a class $`𝒳`$ of subalgebras of $`L`$, we shall mean the data, for every field $`K`$ which is a finite extension of $`k`$, consisting of a family $`𝒳(K)`$ of subalgebras of $`LK`$. For every $`l,n`$, and every field $`K`$ which is a finite extension of $`k`$, we define the subsets$`:`$ $`A_n(𝒳)(K)`$ $`:=`$ $`\{H𝒳(K):\mathrm{codim}_{LK}H=n\}`$ $`A_{l,n}(𝒳)(K)`$ $`:=`$ $`\{HA_n(𝒳)(K):HL_lK\}.`$ Hence $`A_{l,n}(𝒳)(K)`$ can be thought of then as a subset of $`X_l(K)`$, and the union of these subsets when varying $`K`$ is a subset $`A_{l,n}(𝒳)`$ of the finite Grassmannian $`X_l`$. ###### 2.9 Definition. We call $`𝒳`$ a constructible class of subalgebras if $`A_{l,n}(𝒳)`$ is a constructible subset of $`X_l`$, for every $`l,n`$. ### 2.10. Examples of constructible classes of subalgebras For $`L`$ an $``$-filtered Lie algebra, we define $`𝒳^{}`$ as the class of closed subalgebras, i.e. $`𝒳^{}(K)`$ is the set of all closed subalgebras of $`LK`$, and $`𝒳^{}`$ as the class of closed subalgebras which are ideals, i.e. $`𝒳^{}(K)`$ is the set of all closed subalgebras of $`LK`$ which are ideals of $`LK`$. If $`L`$ is a $``$-filtered Lie algebra, the corresponding classes of commensurable subalgebras are defined as follows: for $`\{,\}`$, $`𝒳_0^{}(K)`$ is the set of subalgebras $`H`$ in $`𝒳^{}(K)`$ which are commensurable with $`L_0K`$. ###### 2.11 Lemma. Let $`L`$ be an $``$-filtered Lie algebra then for $`\{,\},`$ $`𝒳^{}`$ is a constructible class of subalgebras. ###### Proof. Choose a basis $`e_1,\mathrm{},e_r`$ for $`L/L_l.`$ Let $`\beta :𝔸_k^{2r}𝔸_k^r`$ denote the bilinear form defining the Lie bracket in $`L/L_l`$ with respect to this basis. Let $`\mathrm{Tr}_{r,k}`$ denote the $`k`$-variety of upper triangular matrices of order $`r`$. Hence the $`K`$-points of $`\mathrm{Tr}_{r,k}`$ are just the upper triangular matrices of order $`r`$ with coefficients in $`K`$. We have a morphism $$h:\mathrm{Tr}_{r,k}\mathrm{Gr}(L/L_l)$$ which takes the matrix $`\left(m_{ij}\right)`$ to the subspace spanned by $$\{m_{11}e_1+\mathrm{}+m_{1r}e_r,\mathrm{},m_{rr}e_r\}.$$ Let $`_{l,n}^{}`$ denote the inverse image of $`A_{l,n}(𝒳^{})`$ under this map. Since there is a finite partition of $`\mathrm{Gr}(L/L_l)`$ into locally closed subvarieties over which $`h`$ induces morphisms which are locally trivial fibrations for the Zariski topology, it suffices to show that $`_{l,n}^{}`$ is a constructible subset of $`\mathrm{Tr}_{r,k}`$. 1. a matrix $`\left(m_{ij}\right)`$ in $`\mathrm{Tr}_{r,k}(K)`$ defines a subspace of codimension $`n`$ in $`LK`$ if and only if the diagonal $`(m_{11},\mathrm{},m_{rr})`$ contains exactly $`n`$ zero entries. 2. a matrix $`\left(m_{ij}\right)`$ in $`\mathrm{Tr}_{r,k}(K)`$ defines a subalgebra of $`LK`$ if and only if for each $`1i<jr`$ there exist $`Y_{ij}^1,\mathrm{},Y_{ij}^rK`$ such that $$\beta (𝐦_i,𝐦_j)=\underset{k=1}{\overset{r}{}}Y_{ij}^k𝐦_k$$ where $`𝐦_i`$ denotes the $`i`$th row of $`\left(m_{ij}\right);`$ 3. a matrix $`\left(m_{ij}\right)`$ in $`\mathrm{Tr}_{r,k}(K)`$ defines an ideal of $`LK`$ if and only if for each $`1i,jr`$ there exist $`Y_{ij}^1,\mathrm{},Y_{ij}^rK`$ such that $$\beta (𝐞_i,𝐦_j)=\underset{k=1}{\overset{r}{}}Y_{ij}^k𝐦_k$$ where $`𝐞_i`$ denotes the $`r`$-tuple with 1 in the $`i`$-th entry and zeros elsewhere. The two conditions on a matrix $`\left(m_{ij}\right)`$ in (1) and (2) define $`_{l,n}^{}`$ as a constructible subset of $`\mathrm{Tr}_{r,k}`$. ∎ If $`L`$ is a finite dimensional $`k[[t]]`$-Lie algebra then we can also define the classes of subalgebras $`𝒳_t^{}`$ with $`𝒳_t^{}(K)`$ the set of closed subalgebras of $`LK`$ which are $`K[[t]]`$-submodules of $`LK`$ and $`𝒳_t^{}`$ defined by $`𝒳_t^{}(K)=𝒳_t^{}(K)𝒳^{}(K)`$. ###### 2.12 Lemma. Let $`L`$ be a finite dimensional $`k[[t]]`$-Lie algebra whose filtration consists of the ideals $`L_j=t^jL`$. Then for $`\{,\},`$ $`𝒳_t^{}`$ are constructible classes of subalgebras. ###### Proof. We just have to add the constructible condition: 1. a matrix $`\left(m_{ij}\right)`$ in $`\mathrm{Tr}_{r,k}(K)`$ defines a $`K[[t]]`$-submodule of $`LK`$ if and only if for each $`1ir`$ there exist $`Y_i^1,\mathrm{},Y_i^rK`$ such that $$t𝐦_i=\underset{k=1}{\overset{r}{}}Y_i^k𝐦_k.$$ We shall see later a better reason for these subalgebras being constructible when we realise $`𝒳_t^{}`$ as a semi-algebraic subset in a suitable arc space. ###### 2.13 Lemma. Suppose that $`L`$ is a $``$-filtered Lie algebra. Then $`𝒳_0^{}`$ is a constructible set of subalgebras. ###### Proof. We want to put a condition on a subspace of $`L_j/L_j`$ that it is a subalgebra of codimension $`n`$ in $`L_0.`$ The trouble is now that $`L_j`$ is not necessarily an ideal in $`L_j.`$ However we use the fact that $`(L_j,L_{2j})L_j.`$ Let $`e_0,e_1,\mathrm{},e_{s_1}`$ be a basis for $`L_{2j}/L_0`$ with $`e_0,e_1,\mathrm{},e_s`$ a basis for $`L_j/L_0`$, and $`e_1,\mathrm{},e_{r_1}`$ be a basis for $`L_0/L_{2j}`$ with $`e_1,\mathrm{},e_r`$ a basis for $`L_0/L_j.`$ Let $`\beta :𝔸_k^{s+r_1+1}\times 𝔸_k^{s+r_1+1}𝔸_k^{s_1+r+1}`$ be the bilinear map defining the Lie bracket from $`L_j/L_{2j}L_{2j}/L_j`$ with respect to these choices of basis. There is still a morphism $$\mathrm{Tr}_{s+r+1,k}\mathrm{Gr}(L_j/L_j)$$ which takes a matrix $`\left(m_{ij}\right)`$, $`sijr`$, to the subspace spanned by $$\{m_{s,s}e_s+\mathrm{}+m_{s,r}e_r,\mathrm{},m_{rr}e_r\}.$$ (1) a matrix $`\left(m_{ij}\right)`$ in $`\mathrm{Tr}_{s+r+1,k}(K)`$ defines a subspace of codimension $`n`$ in $`L_0K`$ if and only if the diagonal $`(m_{11},\mathrm{},m_{rr})`$ contains exactly $`n_1`$ zero entries and the diagonal $`(m_{s,s},\mathrm{},m_{00})`$ contains $`n_2`$ non zero entries and $`n_1+n_2=n`$. (2) the subspace $`H=m_{s,s}e_s+\mathrm{}+m_{s,r}e_r,\mathrm{},m_{rr}e_r+L_jK`$ defines a subalgebra of $`LK`$ if and only if, for $`si<jr`$ and for all $`\lambda =(\lambda _{r+1},\mathrm{},\lambda _{r_1}),\mu =(\mu _{r+1},\mathrm{},\mu _{r_1})K^{r_1r}`$ there exist $`Y_{ij}^s,\mathrm{},Y_{ij}^rK`$ such that $$\beta (𝐦_i,\lambda ,𝐦_j,\mu )=(0,\mathrm{},0,\underset{k=s}{\overset{r}{}}Y_{ij}^k𝐦_k)$$ where there are $`s_1s`$ zeros. The point is that we are going to be guaranteed $`(H,L_{2j})H`$ so we just need to check that $`H`$ is a subalgebra modulo $`L_{2j}`$ which is finite dimensional and therefore leads to a constructible subset. ∎ Proving that ideals in a $``$-filtered Lie algebra define a constructible set is slightly more problematic since we need to check that the action of all the $`L_j`$ stabilise the candidate ideal $`H.`$ However, if $`L`$ is finitely generated we need only check these finite number of generators: ###### 2.14 Lemma. Suppose that $`L`$ is a finitely generated $``$-filtered Lie algebra. Then $`𝒳_0^{}`$ is a constructible set of subalgebras. ###### Proof. Let $`f_1,\mathrm{},f_d`$ be a finite set of generators. Then, for any fixed $`j`$, there exists some $`N(j)`$ and $`M(j)`$ such that $`(f_i,L_{N(j)})L_j`$ and $`(f_i,L_j)L_{M(j)}.`$ So we can carry out the same analysis essentially as the previous lemma. Let $`e_0,e_1,\mathrm{},e_{s_1}`$ be a basis for $`L_{M(j)}/L_0`$ with $`e_0,e_1,\mathrm{},e_s`$ a basis for $`L_j/L_0`$, and $`e_1,\mathrm{},e_{r_1}`$ be a basis for $`L_0/L_{N(j)}`$ with $`e_1,\mathrm{},e_r`$ a basis for $`L_0/L_j.`$ Let $`\phi _i:𝔸_k^{s+r_1+1}𝔸_k^{s_1+r+1}`$ define the action of $`f_i`$ taking $`L_j/L_{N(j)}L_{M(j)}/L_j`$ with respect to these choices of basis. Then (2) the subspace $`H=m_{s,s}e_s+\mathrm{}+m_{s,r}e_r,\mathrm{},m_{rr}e_r+L_jK`$ defines an ideal of $`LK`$ if and only if it is a subalgebra of $`LK`$ and for $`sjr`$ and $`i=1,\mathrm{},d`$ for all $`\lambda =(\lambda _{r+1},\mathrm{},\lambda _{r_1})K^{r_1r}`$ there exist $`Y_{ij}^s,\mathrm{},Y_{ij}^rK`$ such that $$\phi _i(𝐦_i,\lambda )=(0,\mathrm{},0,\underset{k=s}{\overset{r}{}}Y_{ij}^k𝐦_k)$$ where there are $`s_1s`$ zeros. ∎ Let $`L`$ be a finite dimensional $`k((t))`$-Lie algebra and $`L_0`$ be some choice of a $`k[[t]]`$-submodule of the same dimension. Set $`L_j=t^jL_0`$ for $`j.`$ Note that in this case there are no commensurable ideals unless $`L`$ is abelian. The reason for this is as follows. Let $`e_1,\mathrm{},e_d`$ be a basis for the $`k[[t]]`$-submodule $`L_0.`$ Then commensurable ideals must be contained in the centre of $`L`$ since if $`HL_j`$ is an ideal and $`(e_i,h)=xL_k\backslash L_{k+1}`$ then $`(t^{kj1}e_i,h)=t^{kj1}xH.`$ But the centre $`Z(L)`$ is a $`k((t))`$-submodule. So if there exists a commensurable ideal $`H`$ then there exists some $`L_jHZ(L).`$ Hence $`Z(L)=L.`$ So the case of commensurable ideals in finite dimensional $`k((t))`$-Lie algebras is not an interesting one. In the same manner as above we can prove: ###### 2.15 Lemma. Suppose that $`L`$ is a finite dimensional $`k((t))`$-Lie algebra. Then $`𝒳_{t,0}^{}`$, defined by $$𝒳_{t,0}^{}(K)=\{H𝒳_0^{}(K):H\text{ is a }K[[t]]\text{-submodule of }LK\},$$ is a constructible class of subalgebras. ###### Proof. Again we just have to add the constructible condition: (3) a matrix $`\left(m_{ij}\right)`$ in $`\mathrm{Tr}_{s+r+1,k}(K)`$ defines a $`K[[t]]`$-submodule if and only if, for each $`sir`$, there exist $`Y_i^s,\mathrm{},Y_i^rK`$ such that $$t𝐦_i=\underset{k=s}{\overset{r}{}}Y_i^k𝐦_k.$$ ### 2.16. Stable classes of subalgebras and motivic zeta functions If $`𝒳`$ is a constructible class of subalgebras then $`A_{l,n}(𝒳)`$ defines an element $`\left[A_{l,n}(𝒳)\right]`$ of the Grothendieck ring $``$. The only way that this can have a limit as $`l\mathrm{}`$ is that the series stabilize at some point. ###### 2.17 Definition. We shall call a constructible class of subalgebras $`𝒳`$ a stable class of subalgebras if for every $`n`$ there exists some $`l`$ such that $`A_{l,n}=A_{k,n}`$ for all $`kl`$. When this is the case we set $`\left[A_n\right]=\left[A_{l,n}\right]`$ and we define the motivic zeta function of $`L`$ and $`𝒳`$ as $$P_{L,𝒳}(T)=\underset{n=0}{\overset{\mathrm{}}{}}\left[A_n\right]T^n.$$ Notice we can now say why it isn’t sensible to count $`k`$-subalgebras of finite codimension in $`k[[t]]^2`$ with trivial Lie structure since they are not a stable system of subalgebras. As we said above, there are too many such subalgebras. In fact $`𝒳^{}`$ (respectively $`𝒳^{}`$) will not be stable for all $``$-filtered or $``$-filtered Lie algebras $`L`$ for which $`\overline{(H,H)}`$ (respectively $`\overline{(L_0,H)})`$ has infinite codimension for any closed subalgebra $`H`$ (respectively ideal $`H)`$ of finite codimension. ### 2.18. Stability of ideals If we only want to count ideals then we can prove stability for all $``$-filtered Lie algebras, for which the closure of terms in the lower central series of $`L_1`$ have finite codimension: ###### 2.19 Theorem. Suppose $`L=_iL_i`$ is an $``$-filtered Lie algebra and $`\overline{\gamma _i(L_1)}`$ has finite codimension for all $`i`$ (where $`\overline{\gamma _i(L)}`$ is the closure of the $`i`$th term of the lower central series with respect to the formal topology defined by $`L_i)`$. Then $`𝒳^{}(L)`$ is a stable class of subalgebras. ###### Proof. Suppose $`H`$ is a closed ideal of codimension $`n`$ in $`L.`$ Let us show that each ideal $`H_{i1}L_i`$ of codimension $`n`$ in $`_{i1}L_i`$ contains some fixed term of the filtration. We can assume $`L_0=L_1`$. Let us consider $`H_i=\left(H\gamma _i(L)\right)\gamma _{i+1}(L)`$ where $`\gamma _i(L)`$ is the lower central series of $`L`$ (which might be distinct from the closure of the lower central series with respect to the filtration $`L_i)`$. The ideal $`H`$ is closed and hence contains some term of the filtration $`L_m\gamma _m(L).`$ There must be some $`1in+1`$ such that $`H_i=`$ $`\gamma _i(L)`$ else the ideal would have codimension greater than $`n.`$ Hence $`H+\gamma _{i+1}(L)\gamma _i(L).`$ So $`H\gamma _i(L)\gamma _i(L)mod\gamma _{i+1}(L).`$ Now $`H\gamma _{i+1}(L)(L,H\gamma _i(L))`$ since $`H`$ is an ideal. Hence $`H\gamma _{i+1}(L)\gamma _{i+1}(L)mod\gamma _{i+2}(L)`$ and therefore $`H+\gamma _{i+2}(L)H+\gamma _{i+1}(L)\gamma _i(L).`$ Continuing this analysis shows that $`H\gamma _{i+e}(L)\gamma _{i+e}(L)mod\gamma _{i+e+1}(L)`$ for all $`e`$ hence $`H+\gamma _{i+e+1}(L)H+\gamma _{i+e}(L)\gamma _i(L).`$ But $`H+\gamma _{i+e+1}(L)=H`$ for some $`e`$ since $`H`$ is a closed subalgebra of finite codimension. This implies that $`H\gamma _{n+1}(L)`$. Since $`H`$ is closed, $`H\overline{\gamma _{n+1}(L)}L_{f(n)}`$. Similarly, if $`K`$ is a finite extension of $`k`$, every closed ideal of codimension $`n`$ in $`LK`$ contains $`L_{f(n)}K`$, whence $`𝒳^{}(L)`$ is a stable class of subalgebras. ∎ ###### 2.20 Corollary. Suppose $`L=_iL_i`$ is an $``$-filtered Lie algebra and $`L_0=L_1.`$ Then $`𝒳^{}(L)`$ is a stable class of subalgebras if and only if $`\overline{\gamma _i(L)}`$ has finite codimension for all $`i.`$ Note that $`\overline{\gamma _i(L)}`$ has finite codimension for all $`i`$ if and only if $`\overline{(L_1,L_1)}`$ has finite codimension. ###### 2.21 Remark. There is a case missing in the above analysis: when $`L_0L_1,`$ $`\overline{(L_0,H)}`$ has finite codimension for all closed ideals $`H`$ of finite codimension but $`\overline{(L_1,L_1)}`$ has infinite codimension. ### 2.22. Stability of subalgebras We now describe several important categories of Lie algebras with stable classes of subalgebras. Note that for two classes of constructible subalgebras $`𝒳^{\mathrm{}}(L)𝒳^{\mathrm{}}(L)`$ once we have proved that $`𝒳^{\mathrm{}}(L)`$ is stable, then $`𝒳^{\mathrm{}}(L)`$ will also be stable. ###### 2.23 Theorem. (1) Let $`L`$ be a finitely dimensional Lie algebra over $`k[[t]]`$ and define the formal topology on $`L`$ via $`L_l=t^lL.`$ Then $`𝒳_t^{}`$ is a stable constructible class of subalgebras for $`\{,\}`$. (2) Suppose that $`L`$ is a finite dimensional $`k((t))`$-Lie algebra. Then $`𝒳_{t,0}^{}`$ for $`\{,\}`$ is a stable class of subalgebras. ###### Proof. This follows from the fact that a $`K[[t]]`$-submodule $`H`$ of codimension $`n`$ in $`LK`$ must contain $`L_nK`$. In (2) we apply this to $`H(L_0K)`$ which must have codimension bounded by $`n`$. Recall that the case $`𝒳_{t,0}^{}`$ either is empty or else $`L`$ is abelian and $`𝒳_{t,0}^{}=𝒳_{t,0}^{}.`$ ### 2.24. Loop algebras Examples of Lie algebras that are finite-dimensional as $`[[t]]`$-Lie algebras arise as the positive parts of loop algebras. These are defined as the complexified algebra of smooth maps from the circle $`S^1`$ (or more generally some compact complex manifold) into a finite dimensional $``$ -Lie algebra $`𝔤_{},`$ and denoted $`L𝔤_{}`$. We can decompose $`L𝔤_{}`$ into its Fourier components: $$L𝔤_{}=\underset{k}{}𝔤_{}t^k.$$ (One uses the notation $`L𝔤`$ to denote the set of loops, i.e. the set of smooth maps before we have complexified.) This is the decomposition into eigenspaces of the action of the circle which rotates the loops bodily. We define the positive part of the loop algebra to be $$L_0𝔤_{}=\underset{k0}{}𝔤_{}t^k.$$ This is then isomorphic to the $`[[t]]`$-Lie algebra $`𝔤_{}[[t]].`$ These examples are then graded Lie algebras, unlike a general $`[[t]]`$-Lie algebra arising from a bilinear form defined over $`[[t]]`$ on a finite dimensional $`[[t]]`$-module, which are not in general graded but just $``$-filtered. Note that $`L_0𝔤_{}((t))`$ is the $``$-subalgebra generated by the rational loops sometimes denoted $`L_{\mathrm{rat}}𝔤_{}`$ (see section 3.5 of ). Also it should be pointed out that $`L_{\mathrm{rat}}𝔤_{}`$ is $``$-filtered with $`L_{\mathrm{rat}}𝔤_{}=_it^iL_0𝔤_{}`$ whilst the Hilbert space of all loops $`L𝔤_{}`$ is something bigger containing $`L_{\mathrm{rat}}𝔤_{}`$ as a dense subspace. There is another important $``$-subalgebra $`L_{\mathrm{pol}}𝔤_{}=𝔤_{}[z,z^1]`$ generated by all the loops which are given by Laurent polynomials in $`z`$ and $`z^1.`$ This is the same as the Lie algebra of polynomial maps from $`^{}`$ into $`𝔤`$ \- $`^{}`$ is the complexification of the circle (see 2.2 of ). In the case that $`𝔤_{}`$ is a simple finite dimensional $``$-Lie algebra, the subalgebras of polynomial loops are what the Kac-Moody algebras are built from. Kac-Moody Lie algebras arise out of the existence of interesting central extensions that these loop algebras have. We shall see that affine Kac-Moody algebras are examples of infinite dimensional Lie algebras for whom the class of all subalgebras is a stable class. ### 2.25. Kac-Moody Lie algebras A Kac-Moody Lie algebra over a field $`k`$ of characteristic zero is defined in the first instance via a presentation associated to a generalized Cartan matrix (see ). The affine Kac-Moody algebras are those whose Cartan matrices have corank 1. They can then be realised as algebras whose derived group is the universal central extension of the Lie algebra $`L_{\mathrm{pol}}𝔤_k=𝔤_kk[t,t^1]`$ of polynomial loops from $`k^{}`$ into a simple finite dimensional Lie algebra $`𝔤_k`$. More generally if $`𝔤_k=𝔤_{1,k}\mathrm{}𝔤_{q,k}`$ is a decomposition of $`𝔤_k`$ into simple factors then the product of algebras $`k\stackrel{~}{}\stackrel{~}{L}𝔤_{j,k}`$ is an affine Kac-Moody Lie algebra where $`kd\stackrel{~}{}\stackrel{~}{L}𝔤_{j,k}`$ means the semi-direct product where the factor $`kd`$ is generated by the derivation $`d=t\frac{d}{dt}`$ of $`\stackrel{~}{L}𝔤_{j,k}`$ and $`\stackrel{~}{L}𝔤_{j,k}`$ is the universal central extension of the polynomial loop algebra $`L_{\mathrm{pol}}𝔤_{j,k}`$ which looks additively like $`kcL_{\mathrm{pol}}𝔤_{j,k}.`$ Not all are covered by this construction. The other affine Kac-Moody Lie algebras arise by taking the same construction associated to the twisted loop algebras. These are defined for a choice of outer automorphism $`\alpha `$ of $`𝔤_k`$ of finite order $`m`$ and for fields $`k`$ containing primitive $`m`$-th roots of unity. One replaces $`L_{\mathrm{pol}}𝔤_k`$ by its subalgebra $`L_{(\alpha )}𝔤_k`$ consisting of the loops $`fL_{\mathrm{pol}}𝔤_k`$ which are equivariant: $$f(\epsilon ^1t)=\alpha (f(t))$$ where $`\epsilon `$ is a primitive $`m`$-th root of unity. The automorphism gives rise to a decomposition of the finite dimensional Lie algebra: $$𝔤_k=\underset{j/m}{}𝔤_j,$$ where $`𝔤_j`$ is the eigenspace of $`\alpha `$ for the eigenvalue $`\epsilon ^j.`$ Conversely any $`/m`$-grading of $`𝔤_k`$ arises like this since the linear transformation of $`𝔤_k`$ given by multiplying the vectors of $`𝔤_j`$ by $`\epsilon ^j`$ is an automorphism $`\alpha `$ of $`𝔤_k`$ which has order $`m.`$ The twisted loop algebra then has the following description as a $``$-graded Lie algebra: $$L_{(\alpha )}𝔤_k=\underset{n}{}𝔤_{n\mathrm{mod}m}t^n.$$ (This is sometimes referred to as a loop algebra by algebraists because of the cyclic filtration and $`L_{\mathrm{pol}}𝔤_k`$ as the loop algebra with trivial $`C_m`$-grading. This seems to be a little confusing and perhaps arose out of a misunderstanding that the original name loop arose from the setting over $``$ of maps from the circle (hence the name loop) rather than a graded Lie algebra whose filtration somehow comes from ‘looping’ the finite graded Lie algebra $`𝔤.)`$ For more details see section 5.3 of . Note that in our context there is little difference between taking polynomial loops or the rational loop space because polynomial loops are dense in the rational loops and the lattices of closed subalgebras of finite codimension will be in one to one correspondence. ###### 2.26 Theorem. If $`L`$ is an affine Kac-Moody Lie algebra with no infinite dimensional abelian quotients and $`L_0`$ is its positive part, then for $`\{,\},`$ $`𝒳^{}(L_0)`$ and $`𝒳_0^{}(L)`$ are stable classes of subalgebras. ###### Proof. Let $`H`$ be a closed subalgebra of $`L_0`$ of codimension $`n.`$ Suppose $`𝔤_k=𝔤_{1,k}\mathrm{}𝔤_{q,k}`$ is the decomposition of the underlying simple Lie algebra. Then $$L_0=\underset{i=1}{\overset{q}{}}L_{(\alpha _i),0}𝔤_{i,k}kc_ikd_i$$ where $`c_i`$ is central, $`d_i`$ is a derivation and $`\alpha _i`$ is an automorphism of $`𝔤_k`$ which has order $`m`$ ($`k`$ is assumed to have primitive $`m`$-th roots of unity) and $$L_{(\alpha ),r}𝔤_k=\underset{nr}{}𝔤_{n\mathrm{mod}m}t^n.$$ We can restrict our attention to showing that a subalgebra $`H_i`$ of codimension $`n_i`$ in $`L_{(\alpha _i),0}𝔤_{i,k}`$ must contain some fixed term $`L_{(\alpha _i),f(n_i)}𝔤_{i,k}`$ of the filtration on $`L_{(\alpha _i),0}𝔤_{i,k}`$ depending only on $`n_i.`$ For convenience we drop the subscript $`i`$ and since we are only considering one twisted algebra at a time we put $`L_r=L_{(\alpha ),r}𝔤_k.`$ The assumption that $`L`$ has no infinite abelian sections means that $`𝔤_k`$ is not abelian and hence $`𝔤_k=(𝔤_k,𝔤_k)`$ is perfect. This implies in turn that $`(L_{im},L_{jm})=L_{(i+j)m}.`$ Let $$H_i=\left(HL_{im}\right)+L_{(i+1)m}/L_{(i+1)m}.$$ The codimension of $`H`$ is then equal to the sum of the codimensions of $`H_i`$ in $`L_{im}/L_{(i+1)m}.`$ Therefore $`H_i`$ $`L_{im}/L_{(i+1)m}`$ for at most $`n`$ values of $`i.`$ Therefore consider $`N2n,`$ there must be some value of $`i`$ with $`1iN/2`$ such that both $`H_i=L_{im}/L_{(i+1)m}`$ and $`H_{Ni}=L_{(Ni)m}/L_{(Ni+1)m}.`$ Hence $`H_N(H_i,H_{Ni})=L_{Nm}/L_{(N+1)m}`$ for all $`N.`$ Thus $`HL_{2nm}.`$ Note that an affine Kac-Moody Lie algebra is finitely generated hence $`𝒳_0^{}(L)`$ is constructible. The key to this proof is that in affine Kac-Moody Lie algebras, each element can be realised as a commutator in many different ways. In fact we can see that the argument applies to the following filtered Lie algebras: ###### 2.27 Definition. Call an $``$-filtered Lie algebra well-covered if it has the property that for each $`n`$ there exists some $`f(n)`$ such that $`L_{f(n)}=(L_i,L_{f(n)i})`$ for $`n+1`$ values of $`i<f(n)/2.`$ Call a $``$-filtered Lie algebra well-covered if its positive part is well-covered. The proof of the theorem above can then easily be adapted to prove: ###### 2.28 Theorem. (1) If the $``$-filtered Lie algebra $`L`$ is well-covered then $`𝒳^{}(L)`$ for $`\{,\}`$ are stable classes of subalgebras. (2) If the $``$-filtered Lie algebra $`L`$ is well-covered then $`𝒳_0^{}(L)`$ is a stable class of subalgebras. The affine Kac-Moody Lie algebras (or rather the underlying loop algebras) make up the bulk of the classification of the simple $``$-graded Lie algebras over an algebraically closed field of characteristic zero whose graded pieces have finite growth. According to the classification of such Lie algebras proved by Mathieu the remaining Lie algebras are those of Cartan type and the Virasoro algebra. ### 2.29. Virasoro algebra The Virasoro algebra is defined for any field $`k`$ of characteristic zero as the unique central extension of the Witt algebra. The Witt algebra $`𝔡`$ is defined by: $$𝔡:=\mathrm{Der}k[t,t^1]=\underset{j}{}kd_j$$ where $`d_j=t^{j+1}\frac{d}{dt}`$ with the following commutation relations: $$[d_i,d_j]=(ij)d_{i+j}.$$ The Lie algebra $`𝔡`$ has a unique (up to isomorphism) non-trivial central extension by a 1-dimensional centre, $`kc`$ say, called the Virasoro algebra $`\mathrm{Vir}`$, which is defined by the following commutation relations: $$[d_i,d_j]=(ij)d_{i+j}+(i^3i)\delta _{i,j}c\text{ where }i,j.$$ The Virasoro algebra, denoted $`\mathrm{Vir}`$, plays an important role in the representation theory of affine Kac-Moody Lie algebra and in quantum field theory. It is a $``$-graded subalgebra in $`\mathrm{Der}\stackrel{~}{L}𝔤_k,`$ the group of derivations of the algebra $`\stackrel{~}{L}𝔤_k,`$. When $`k=`$, it is the complexified Lie algebra of smooth vector fields on the circle. The positive part of this Lie algebra $`𝔡^+:=_{j1}kd_j`$ also arises as the graded Lie algebra of the so-called Nottingham group over $`k.`$ This is defined as the group whose underlying set is the set of formal power series $`tk[[t]]`$ and whose group operation is substitution of power series. It is a certain subgroup of the automorphism group of $`k[[t]],`$ called the wild automorphisms. ### 2.30. Cartan subalgebras The Cartan algebras are defined as follows. Let $`n1`$ and let $`𝐖_n`$ be the algebra of derivations of the polynomial ring $`R=k[X_1,\mathrm{},X_n].`$ So $`𝐖_n=RD_1+\mathrm{}+RD_n`$, where $`D_i=\frac{d}{dX_i}`$. Thus $`𝐖_n`$ acts on the Grassmann algebra of Kähler differential forms on $`R`$. We define three subalgebras of $`𝐖_n.`$ (1) The subalgebra $`𝐒_n`$, called the special algebra, consists of those derivations annihilating the differential form $`\nu =dX_1\mathrm{}dX_n.`$ We can describe this set of derivations as follows: $$𝐒_n=\{\underset{j=1}{\overset{n}{}}a_jD_j:\underset{j=1}{\overset{n}{}}D_j\left(a_j\right)=0\}.$$ It is additively spanned by $`D_{ij}(f)=D_j(f)D_iD_i(f)D_j`$, where $`1i,jn`$ and $`f`$ ranges over $`R.`$ (2) Suppose now that $`n=2m`$ is even then $`𝐇_{2m},`$ called the hamiltonian algebra, consists of derivations annihilating $`\omega =_{1im}dX_idX_{m+i}`$. Define the following notation: $`j^{}`$ $`=`$ $`j+m\text{ and }\sigma (j)=1\text{ if }1jm`$ $`j^{}`$ $`=`$ $`jm\text{ and }\sigma (j)=1\text{ if }m+1j2m.`$ Then we have the following description of the hamiltonian algebra: $$𝐇_{2m}=\{\underset{j=1}{\overset{2m}{}}a_jD_j:\sigma (j^{})D_i\left(a_j^{}\right)=\sigma (i^{})D_j\left(a_i^{}\right),1i,j2m\}.$$ Every element of $`𝐇_{2m}`$ can be represented as $`D_H(f)`$ for some $`fR`$ where $$D_H(f)=\underset{j=1}{\overset{2m}{}}\sigma (j)D_j(f)D_j^{}.$$ (3) Finally if $`n=2m+1`$ is odd then $`𝐊_{2m+1},`$ called the contact algebra, consists of derivations $``$ such that $`\alpha =f\alpha `$ for some polynomial $`f`$ where $$\alpha =dX_{2m+1}+\underset{1im}{}X_idX_{m+i}X_{m+i}dX_i.$$ For every derivation $`D=_{j=1}^na_j\frac{d}{dX_j}`$ define $$u(D)=\underset{j=1}{\overset{2m}{}}\sigma (j)X_j\frac{da_j^{}}{dX_n}\frac{da_n}{dX_n}.$$ Then $`D`$ is contained in $`𝐊_{2m+1}`$ if and only if for each $`1i2m`$ $$\sigma (i)a_i^{}\underset{j=1}{\overset{2m}{}}\sigma (j)X_j\frac{da_j^{}}{dX_i}\frac{da_n}{dX_i}=\sigma (i)X_i^{}u(D).$$ Every element of $`𝐊_{2m+1}`$ can be represented as $`D_K(f)=_{j=1}^{2m+1}f_jD_j`$ where $`f_j`$ $`=`$ $`X_jD_{2m+1}(f)+\sigma (j^{})D_j^{}(f)\text{ for }j2m`$ $`f_{2m+1}`$ $`=`$ $`2f{\displaystyle \underset{j=1}{\overset{2m}{}}}\sigma (j)X_jf_j^{}.`$ The Lie algebras $`𝐖_n,𝐒_n,𝐇_{2m}`$ and $`𝐊_{2m+1}`$ are called Lie algebras of Cartan type. The usual grading of the polynomial ring $`k[X_1,\mathrm{},X_n]`$ induces a grading of $`𝐖_n,𝐒_n`$ and $`𝐇_{2m}.`$ If $`n=2m+1`$ and $`m1`$ then there exists a unique grading of $`k[X_1,\mathrm{},X_{2m+1}]`$ such that $`X_1,\mathrm{},X_{2m}`$ are homogeneous of degree 1 and $`X_{2m+1}`$ is homogeneous of degree 2. The induced grading for $`𝐊_{2m+1}`$ is called the natural grading of $`𝐊_{2m+1}.`$ ###### 2.31 Theorem. The Virasoro algebra $`\mathrm{Vir}`$ and the Lie algebras of Cartan type $`𝐖_n,𝐒_n,𝐇_{2m}`$ and $`𝐊_{2m+1}`$ are well-covered. ###### Proof. This depends on the fact that any derivation $$X_1^{m_1}\mathrm{}X_i^{m_i+1}\mathrm{}X_n^{m_n}\frac{d}{dX_i}$$ with $`m_1+\mathrm{}+m_n=N`$ can be realised as a linear multiple of the Lie product $$(X_1^{r_1}\mathrm{}X_i^{r_i+1}\mathrm{}X_n^{r_n}\frac{d}{dX_i},X_1^{s_1}\mathrm{}X_i^{s_i+1}\mathrm{}X_n^{s_n}\frac{d}{dX_i})$$ where $`r_j+s_j=m_j`$ and $`r_is_i`$. For the Lie algebras $`\mathrm{Vir}`$ and $`𝐖_n`$ the result follows, since we can then realise the derivation in the grading of weight $`N`$ as a Lie product of elements of weight $`i`$ and $`Ni`$ for all values of $`i<N/2.`$ We shall have to work harder to get the corresponding result for the other Cartan algebras $`𝐒_n,𝐇_{2m}`$ or $`𝐊_{2m+1}`$, since we need to realise elements of weight $`N`$ as sums of commutators of elements of weight $`i`$ and $`Ni`$ which still stabilize the form defining the corresponding Cartan algebra. For $`𝐒_n`$ let us show why, for $`𝐗^𝐜=X_1^{c_1}\mathrm{}X_n^{c_n}`$ where $`𝐜=c_1+\mathrm{}+c_n=nN`$, that $`D_{ij}(𝐗^𝐜)`$ can be realised as a linear combination of commutators of elements of weight $`nNt`$ and $`t`$ for $`N`$ values of $`t.`$ We are going to use the following identity which can easily be checked. Suppose that $`b_j=0`$ then (2.31.1) $$[D_{ki}(𝐗^𝐚),D_{kj}(𝐗^𝐛)]=a_ib_kD_{kj}(𝐗^{𝐚+𝐛𝐞_i𝐞_k})+a_kb_kD_{ij}(𝐗^{𝐚+𝐛2𝐞_k})$$ Since $`𝐜=nN,`$ we know that there is some $`c_kN.`$ Case 1. Suppose firstly that $`ki,j.`$ Then choose $`𝐚`$ and $`𝐛`$ with $`𝐜=𝐚+𝐛2𝐞_k,`$ $`b_j=0`$, $`a_i=0`$ and $`a_k=c_k+1l`$ and $`b_k=l+1`$. For $`l=0,\mathrm{},c_k`$, we get $$[D_{ki}(𝐗^𝐚),D_{kj}(𝐗^𝐛)]=a_kb_kD_{ij}(𝐗^{𝐚+𝐛2𝐞_k})=\lambda D_{ij}(𝐗^𝐜)$$ with $`\lambda 0`$ and $`D_{ki}(𝐗^𝐚)`$ of weight $`a_1+\mathrm{}+a_{k1}+a_{k+1}+\mathrm{}+a_n+c_k+1l`$ and $`D_{kj}(𝐗^𝐛)`$ of weight $`b_1+\mathrm{}+b_{k1}+b_{k+1}+\mathrm{}+b_n+l+1.`$ Case 2. If $`k=i,`$ then choose some $`s\{1,\mathrm{},n\}\backslash \{i,j\}`$ and define $`𝐚`$ and $`𝐛`$ with $`𝐜=𝐚+𝐛2𝐞_s,`$ $`b_j=0`$, $`a_s=1`$ and $`a_i=c_i+1l`$ and $`b_i=l+1`$. Form $$[D_{si}(𝐗^𝐚),D_{sj}(𝐗^𝐛)]=a_ib_sD_{sj}(𝐗^{𝐚+𝐛𝐞_i𝐞_s})+a_sb_sD_{ij}(𝐗^{𝐚+𝐛2𝐞_s})$$ where note that $`a_s0b_s.`$ To deal with $`a_ib_sD_{sj}(𝐗^{𝐚+𝐛𝐞_i𝐞_s})`$ we can use the analysis of case 1 since the weight is concentrated still at $`c_i`$ but now we are considering a derivation $`D_{sj}`$ with $`is,j.`$ Hence choose $`𝐚^{}=𝐚+𝐞_i𝐞_s`$ then $`𝐚+𝐛𝐞_i𝐞_s=𝐚^{}+𝐛2𝐞_i`$, and $`a_s^{}=0`$, $`a_i^{}=c_i+2l`$ and $`b_i=l+1`$ where $`l=0,\mathrm{},c_i.`$ (Note that the weight at $`X_i`$ of $`D_{sj}(𝐗^{𝐚+𝐛𝐞_i𝐞_k})`$ is now $`c_i+1.)`$ Hence we have $$[D_{si}(𝐗^𝐚),D_{sj}(𝐗^𝐛)](a_ib_s)(a_i^{}b_i^{})^1[D_{is}(𝐗^𝐚^{}),D_{ij}(𝐗^𝐛)]=\lambda D_{ij}(𝐗^𝐜)$$ with $`\lambda 0,`$ $`D_{si}(𝐗^𝐚)`$ and $`D_{is}(𝐗^𝐚^{})`$ have weight $`a_1+\mathrm{}+a_{i1}+a_{i+1}+\mathrm{}+a_n+c_i+1l`$ and $`D_{sj}(𝐗^𝐛)`$ and $`D_{ij}(𝐗^𝐛)`$ have weight $`b_1+\mathrm{}+b_{i1}+b_{i+1}+\mathrm{}+b_n+l+1.`$ This proves then that $`𝐒_n`$ is well-covered. For the Hamiltonian algebra, we can use the following identity (see Chapter 4 Lemma 4.3 (1) of ) $`[D_H(𝐗^𝐚),D_H(X_i^{b_i})]`$ $`=`$ $`D_H\left(D_H(𝐗^𝐚)(X_i^{b_i})\right)`$ $`=`$ $`D_H\left({\displaystyle \underset{j=1}{\overset{2m}{}}}\sigma (j^{})D_j^{}(f)D_j(X_i^{b_i})\right)`$ $`=`$ $`D_H\left(\sigma (i^{})a_i^{}𝐗^{𝐚𝐞_i^{}}b_iX_i^{b_i1}\right).`$ Consider now realising $`D_H(𝐗^𝐜)`$ of weight $`Nn.`$ So there exists some $`i`$ such that $`c_iN.`$ We then put $`b_i=l+1`$ and $`𝐚=𝐜+𝐞_i^{}b_i+1`$ with $`l=0,\mathrm{},c_i`$ and hence can express $`D_H(𝐗^𝐜)`$ as a commutator of elements of the Hamiltonian algebra of weight $`Nnl`$ and $`l`$ for $`l=0,\mathrm{},c_i`$. Hence $`𝐇_{2m}`$ is well-covered. For the contact algebra $`𝐊_n=𝐊_{2m+1}`$, we shall use the following identity (see Chapter 4 Proposition 5.2 of ): $$[D_K(f),D_K(g)]=D_K\left(D_K(f)(g)2gD_{2m+1}(f)\right).$$ This algebra is slightly trickier than the previous algebras. We shall need to consider realising $`D_K(𝐗^𝐜)`$ of weight $`4Nn(n1)`$ as a linear combination of commutators of weight $`4Nn(n1)l`$ and $`l`$ for $`N`$ different values of $`l.`$ Recall that in this algebra $`X_n`$ has weight two whilst the other variables have weight one. There exists some $`i`$ with $`c_i2N(n1).`$ Case 1. Suppose that $`in.`$ We prove by induction on $`c_n`$ that we can express $`D_K(𝐗^𝐜)`$ as a linear combination of commutators of weight $`4Nn(n1)l`$ and $`l`$ for $`l=0,\mathrm{},N.`$ Put $`f=𝐗^𝐚`$ and $`g=X_i^{b_i}`$ where $`𝐚=𝐜+𝐞_i^{}l𝐞_i`$ and $`b_i=l+1.`$ Then (2.31.2) $$[D_K(𝐗^𝐚),D_K(X_i^{b_i})]=D_K\left(X_ia_n𝐗^{𝐚𝐞_n}b_iX_i^{b_i1}\right)+D_K\left(\sigma (i^{})a_i^{}𝐗^{𝐚𝐞_i^{}}b_iX_i^{b_i1}\right)$$ Suppose that $`c_n=a_n=0.`$ Then we are done. We then suppose by induction that we have proved the claim for $`c_n1.`$ Then the identity (2.31.2) shows how to express $`D_K(𝐗^𝐜)`$ as a linear combination of commutators as desired since the first expression on the right hand side of (2.31.2) can be dealt with by the inductive hypothesis. Case 2. Suppose now that $`i=n.`$ We show how to shift the weight from $`X_n`$ to the other variables such that eventually we are in a position to invoke case 1. We use the following identity $$[D_K(X_n^{b_n}),D_K(𝐗^𝐚)]=D_K\left(X_n^{b_n},𝐗^𝐚\right)$$ where $$\begin{array}{cc}\hfill X_n^{b_n},𝐗^𝐚& =D_K(X_n^{b_n})(𝐗^𝐚)2𝐗^𝐚D_n(X^{b_n})\hfill \\ & =\left(\left(\left(\underset{j=1}{\overset{2m}{}}a_j\right)2\right)b_n+2a_n\right)𝐗^𝐚X_n^{b_n1}\hfill \\ & \underset{j=1}{\overset{2m}{}}\sigma (j)b_na_n𝐗^{𝐚+(b_n2)𝐞_n+𝐞_j+𝐞_j^{}}.\hfill \end{array}$$ Note that we shall let $`b_n=l+1`$ for $`l=0,\mathrm{},N`$ and $`a_n=c_nl.`$ Since $`c_n2N(n1)`$ this will mean that $`\left(\left(\left(_{j=1}^{2m}a_j\right)2\right)b_n+2a_n\right)`$ is always non-zero. The choice of $`c_n2N(n1)`$ means that we can keep on applying the above identity to realise the expressions $`\sigma (j)b_na_n𝐗^{𝐚+(b_n2)𝐞_n+𝐞_j+𝐞_j^{}}`$ as commutators $`[D_K(X_n^{b_n^{}}),D_K(𝐗^𝐚^{})]`$ with an error term which has the values of $`c_i`$ increasing for $`in`$ and $`c_n`$ decreasing whilst still ensuring that $`c_n>N`$ so that we can choose $`b_n^{}=l+1`$ for $`l=0,\mathrm{},N.`$ Eventually the error terms will have some $`c_i`$ with $`c_i>N`$ and we can apply case 1 to finish the realisation. ∎ ###### 2.32 Corollary. If $`L`$ is a simple $``$-graded Lie algebra over an algebraically closed field of characteristic zero of finite growth then $`𝒳_0^{}(L)`$ and $`𝒳^{}(L_0)`$ for $`\{,\}`$ are stable classes of subalgebras. The proof above depends on the classification proved by Mathieu and a case by case analysis of each class of Lie algebra in the classification. It may be possible that there is a more direct argument based on the simplicity of the $``$-graded Lie algebra. Note that since $`L`$ is simple $`𝒳_0^{}(L)`$ is actually empty or consists of $`L`$ if $`L_0`$ has finite codimension in $`L`$. As we have pointed out, an $``$-filtered Lie algebra has many ideals so we cannot expect it to be simple. However in this context, the concept of being simple is replaced by that of being just infinite \- that is an infinite dimensional Lie algebra all of whose proper quotients are finite dimensional. It is conjectured by Shalev and Zelmanov (see section 6.5) that the just infinite $``$-filtered Lie algebras over an algebraically closed field of characteristic zero for which $`L_0=L_1`$ and $`dim(L_i/L_{i+1})`$ are uniformly bounded are in the following list: (1) $`L`$ is soluble; (2) The completion of $`L`$ is commensurable with the positive part of a loop algebra; (3) The completion of $`L`$ is commensurable to the completion of $`^+.`$ Note that only (2) and (3) will have stable subalgebras. ### 2.33. Non-stability of free graded Lie algebras Stability for infinite Lie algebras is the analogue of determining when a group has only a finite number of subgroups of each given index which can then be counted using a standard Dirichlet series. For groups the condition that the group be finitely generated (either as an abstract group or a topological group) was sufficient to ensure this condition. Note that being finitely generated will not suffice in the context of infinite dimensional Lie algebras as the following Theorem indicates: ###### 2.34 Theorem. Let $`L=_{i1}\gamma _i(L)/\gamma _{i+1}(L)`$ be the free two generated infinite dimensional graded Lie algebra where $`\gamma _i(L)`$ is the lower central series. Then for each $`i`$ there exist closed subalgebras of codimension 2 not containing $`\gamma _i(L).`$ ###### Proof. Let $`x`$ and $`y`$ be the free generators. If $`w`$ is any Lie word in $`x`$ and $`y`$ we define the length $`l(w)`$ to be the number of terms in $`w.`$ We take a Hall set $`H=\{w_i:i1\}`$ as a basis for $`L`$ (see II.2.10) which is defined as follows: (1) if $`w_iH`$ and $`w_jH`$ and $`l(w_i)<l(w_j)`$ then $`i<j;`$ (2) $`w_1=x,w_2=y`$ and $`w_3=(xy);`$ (3) an element $`w`$ of length $`3`$ belongs to $`H`$ if and only if it is of the form $`(a(bc))`$ with $`a,b,c`$ in $`H`$, $`(bc)H`$ and $`ba<bc`$ and $`b<c`$ where the ordering is defined by the ordering on the index set. For example the construction provided by Proposition 11 of II.2.10 starts with the following basis for each layer $`\gamma _i(L)/\gamma _{i+1}(L)`$: $$\begin{array}{cccc}\gamma _1(L)/\gamma _2(L)\hfill & w_1=x\hfill & w_2=y\hfill & \\ \gamma _2(L)/\gamma _3(L)\hfill & w_3=(xy)\hfill & & \\ \gamma _3(L)/\gamma _4(L)\hfill & w_4=(x(xy))\hfill & w_5=(y(xy))\hfill & \\ \gamma _4(L)/\gamma _5(L)\hfill & w_6=(x(x(xy)))\hfill & w_7=(y(x(xy)))\hfill & w_8=(y(y(xy)))\hfill \\ \gamma _5(L)/\gamma _6(L)\hfill & w_9=(x(x(x(xy))))\hfill & w_{10}=(y(x(x(xy))))\hfill & w_{11}=(y(y(x(xy))))\hfill \\ & w_{12}=(y(y(y(xy))))\hfill & w_{13}=((xy)(x(xy)))\hfill & w_{14}=((xy)(y(xy)))\hfill \end{array}$$ Let $`n_i1`$ be the dimension of $`\gamma _1(L)/\gamma _i(L).`$ Then we claim that the additive subspace of $`L`$ generated by the following basis elements is actually a subalgebra: $$H_i=\underset{l1,n_i}{}kw_l.$$ So this is a vector subspace of codimension 2 which skips the basis elements $`w_1=x`$ and $`w_{n_i}=(x(x\mathrm{}(xy)\mathrm{}))`$ where $`x`$ appears $`i1`$ times. The Hall basis has the property that each word $`w_l`$ has a unique decreasing factorization $`w_l=(w_jw_k)`$ with $`j<k<l`$ (see Corollary 4.7 of ). We prove by induction on the length of $`w_j`$ for $`1<j<n_i`$ and $`k1,n_i`$ that $`(w_jw_k)=_{l1,n_i}a_lw_l.`$ If $`w_j`$ has length $`1`$ then $`w_j=y`$ and by condition (3) $`(yw_k)`$ is in $`H.`$ The property of unique factorization implies it is not the element $`w_{n_i}=(xw_{n_{i1}}).`$ Suppose we have proved that $`(w_jw_k)=_{l1,n_i}a_lw_l`$ for all words $`w_j`$ of length less that $`m`$ and take $`w_j`$ of length $`m>1.`$ There is a unique decreasing factorization $`w_j=(w_{j_1}w_{j_2})`$ with $`l(w_{j_i})<l(w_j).`$ Using the Jacobi identity we can rewrite $`(w_jw_k)`$ $`=`$ $`((w_{j_1}w_{j_2})w_k)`$ $`=`$ $`(w_{j_1}(w_{j_2}w_k))(w_{j_2}(w_{j_1}w_k)).`$ Now we can use our induction hypothesis applied to $`w_{j_1}`$ and $`w_{j_2}.`$ Obviously the same argument applies to the $`d`$-generated free Lie algebra. So the existence of a motivic zeta function for an infinite dimensional Lie algebra does not appear to be so straightforward. The condition of being well-covered will suffice, but this is a special sort of property which depends on the internal workings of a Lie algebra. It would be nice to have a more transparent Lie theoretic condition which will ensure the stability of the class of subalgebras. We therefore raise the following: ###### 2.35 Problem. Determine a criterion for the class of all subalgebras of finite codimension to be stable in an $``$-filtered or $``$-filtered Lie algebra. ¿From now on we shall focus on the situation of a finite dimensional $`k[[t]]`$-Lie algebra and $`𝒳`$ a class of $`k[[t]]`$-subalgebras. ## 3. Motivic integration and rationality of Poincaré series In this section, we review material from and which will be used in the present paper. ### 3.1. Scheme of arcs For $`X`$ a variety over $`k`$, we will denote by $`(X)`$ the scheme of germs of arcs on $`X`$. It is a scheme over $`k`$ and for any field extension $`kK`$ there is a natural bijection $$(X)(K)\mathrm{Mor}_{k\mathrm{schemes}}(SpecK[[t]],X)$$ between the set of $`K`$-rational points of $`(X)`$ and the set of germs of arcs with coefficients in $`K`$ on $`X`$. We will call $`K`$-rational points of $`(X)`$, for $`K`$ a field extension of $`k`$, arcs on $`X`$, and $`\phi (0)`$ will be called the origin of the arc $`\phi `$. More precisely the scheme $`(X)`$ is defined as the projective limit $$(X):=\underset{}{\mathrm{lim}}_n(X)$$ in the category of $`k`$-schemes of the schemes $`_n(X)`$ representing the functor $$R\mathrm{Mor}_{k\mathrm{schemes}}(SpecR[t]/t^{n+1}R[t],X)$$ defined on the category of $`k`$-algebras. The existence of $`_n(X)`$ is well known (cf. ) and the projective limit exists since the transition morphisms are affine. We shall denote by $`\pi _n`$ the canonical morphism, corresponding to truncation of arcs, $$\pi _n:(X)_n(X).$$ The schemes $`(X)`$ and $`_n(X)`$ will always be considered with their reduced structure. ### 3.2. Semi-algebraic geometry ¿From now on we will denote by $`\overline{k}`$ a fixed algebraic closure of $`k`$, and by $`\overline{k}((t))`$ the fraction field of $`\overline{k}[[t]]`$, where $`t`$ is one variable. Let $`x_1,\mathrm{},x_m`$ be variables running over $`\overline{k}((t))`$ and let $`\mathrm{}_1,\mathrm{},\mathrm{}_r`$ be variables running over $``$. A semi-algebraic (resp. $`k[[t]]`$-semi-algebraic) condition $`\theta (x_1,\mathrm{},x_m;\mathrm{}_1,\mathrm{},\mathrm{}_r)`$ is a finite boolean combination of conditions of the form (1) $`\mathrm{ord}_tf_1(x_1,\mathrm{},x_m)\mathrm{ord}_tf_2(x_1,\mathrm{},x_m)+L(\mathrm{}_1,\mathrm{},\mathrm{}_r)`$ (2) $`\mathrm{ord}_tf_1(x_1,\mathrm{},x_m)L(\mathrm{}_1,\mathrm{},\mathrm{}_r)(modd)`$ (3) $`h(\overline{ac}(f_1(x_1,\mathrm{},x_m)),\mathrm{},\overline{ac}(f_m^{}(x_1,\mathrm{},x_m)))=0,`$ where $`f_i`$ are polynomials with coefficients in $`k`$ (resp. $`f_i`$ are polynomials with coefficients in $`k[[t]]`$), $`h`$ is a polynomial with coefficients in $`k`$, $`L`$ is a polynomial of degree $`1`$ over $``$, $`d`$, and $`\overline{ac}(x)`$ is the coefficient of lowest degree of $`x`$ in $`\overline{k}((t))`$ if $`x0`$, and is equal to 0 otherwise. Here we use the convention that $`\mathrm{}+\mathrm{}=\mathrm{}`$ and $`\mathrm{}\mathrm{}\mathrm{mod}d`$, for all $`\mathrm{}`$. In particular the condition $`f(x_1,\mathrm{},x_m)=0`$ is a semi-algebraic condition (resp. a $`k[[t]]`$-semi-algebraic condition), for $`f`$ a polynomial over $`k`$ (resp. over $`k[[t]]`$). A subset of $`\overline{k}((t))^m\times ^r`$ defined by a semi-algebraic (resp. $`k[[t]]`$-semi-algebraic) condition is called semi-algebraic (resp. $`k[[t]]`$-semi-algebraic). One defines similarly semi-algebraic and $`k[[t]]`$-semi-algebraic subsets of $`K((t))^m\times ^r`$ for $`K`$ an algebraically closed field containing $`\overline{k}`$. A function $`\alpha :\overline{k}((t))^m\times ^n`$ is called simple (resp. $`k[[t]]`$-simple) if its graph is semi-algebraic (resp. $`k[[t]]`$-semi-algebraic). Let $`X`$ be an algebraic variety over $`k`$. For $`x(X)`$, we denote by $`k_x`$ the residue field of $`x`$ on $`(X)`$, and by $`\stackrel{~}{x}`$ the corresponding rational point $`\stackrel{~}{x}(X)(k_x)=X(k_x[[t]])`$. When there is no danger of confusion we will often write $`x`$ instead of $`\stackrel{~}{x}`$. A semi-algebraic family of semi-algebraic subsets (resp. $`k[[t]]`$-semi-algebraic family of $`k[[t]]`$-semi-algebraic subsets) (for $`n=0`$ a semi-algebraic subset (resp. $`k[[t]]`$-semi-algebraic subset)) $`A_{\mathrm{}}`$, $`\mathrm{}^n`$, of $`(X)`$ is a family of subsets $`A_{\mathrm{}}`$ of $`(X)`$ such that there exists a covering of $`X`$ by affine Zariski open sets $`U`$ with $$A_{\mathrm{}}(U)=\{x(U):\theta (h_1(\stackrel{~}{x}),\mathrm{},h_m(\stackrel{~}{x});\mathrm{})\},$$ where $`h_1,\mathrm{},h_m`$ are regular functions on $`U`$ and $`\theta `$ is a semi-algebraic condition (resp. $`k[[t]]`$-semi-algebraic condition). Here the $`h_i`$’s and $`\theta `$ may depend on $`U`$ and $`h_i(\stackrel{~}{x})`$ belongs to $`k_x[[t]]`$. Let $`A`$ be a semi-algebraic subset (resp. $`k[[t]]`$-semi-algebraic subset) of $`(X)`$. A function $`\alpha :A\times ^n\{\mathrm{}\}`$ is called simple (resp. $`k[[t]]`$-simple) if the family of subsets $`\{x(X):\alpha (x,\mathrm{}_1,\mathrm{},\mathrm{}_n)=\mathrm{}_{n+1}\}`$, $`(\mathrm{}_1,\mathrm{},\mathrm{}_{n+1})^{n+1}`$, is a semi-algebraic family of semi-algebraic subsets (resp. a $`k[[t]]`$-semi-algebraic family of $`k[[t]]`$-semi-algebraic subsets) of $`(X)`$. An important fact is that if $`A`$ is a $`k[[t]]`$-semi-algebraic subset of $`(X)`$, then $`\pi _n(A)`$ is a constructible subset of $`_n(X)`$ (cf. Proposition 1.7 of ). ###### 3.3 Remark. Motivic integration is developed in for semi-algebraic subsets of $`(X)`$, when $`X`$ is an algebraic variety over $`k`$, and simple functions. In the paper , this is extended to what is called there $`t`$-semi-algebraic subsets of $`(X)`$ and $`t`$-simple functions, which are defined similarly as $`k[[t]]`$-semi-algebraic subsets of $`(X)`$ and $`k[[t]]`$-simple functions, but replacing $`k[[t]]`$ by $`k[t]`$. In fact, as mentioned in Remark 1.18 of , all results in §1 of (before 1.17) may be extended, with similar proofs, to cover the case of $`k[[t]]`$-semi-algebraic subsets and $`k[[t]]`$-simple functions. Hence when we shall quote a result from , we shall use its extension to $`k[[t]]`$-semi-algebraic subsets and $`k[[t]]`$-simple functions without further comment. ### 3.4. Motivic integration ¿From now on we assume $`X`$ is a smooth algebraic variety over $`k`$ of pure dimension $`d`$. Let $`A`$ be a $`k[[t]]`$-semi-algebraic subset $`(X)`$. We say $`A`$ is stable at level $`n`$ if $`A=\pi _n^1\pi _n(A)`$. Remark that if $`A`$ is stable at level $`n`$, then $`A`$ is stable at level $`m`$, for any $`mn`$. We say $`A`$ is stable if it is stable at some level. Denote by $`𝐁^t`$ the set of all $`k[[t]]`$-semi-algebraic subsets of $`(X)`$, and by $`𝐁_0^t`$ the set of all $`A`$ in $`𝐁^t`$ which are stable. Clearly there is a unique additive measure $$\stackrel{~}{\mu }:𝐁_0^t_{\mathrm{loc}}$$ satisfying $$\stackrel{~}{\mu }(A)=[\pi _n(A)]𝐋^{(n+1)d},$$ when $`A`$ is stable at level $`n`$. In fact, the condition of being $`k[[t]]`$-semi-algebraic is superfluous here. By the same formula one may define $`\stackrel{~}{\mu }(A)`$ for $`A`$ cylindrical at level $`n`$, i.e. subsets of $`(X)`$ of the form $`A=\pi _n^1(C)`$ with $`C`$ constructible. One says $`A`$ is cylindrical if it is cylindrical at some level. Let $`A`$ be in $`𝐁_0^t`$ and let $`\alpha :A`$ be a $`k[[t]]`$-simple function (or, more generally, assume $`A`$ and the fibers of $`\alpha `$ are cylindrical). By Lemma 2.4 of and Lemma A.3 of , $`\|\alpha \|`$ is bounded, and we can define (3.4.1) $$_A𝐋^\alpha 𝑑\stackrel{~}{\mu }:=𝐋^n\stackrel{~}{\mu }(\alpha ^1(n)),$$ the sum on the right hand side being finite. ### 3.5. Completion We now explain how one extends $`\stackrel{~}{\mu }`$ to non stable $`k[[t]]`$-semi-algebraic subsets by using a completion of $`_{\mathrm{loc}}`$. This is indeed similar to the use of real numbers for defining $`p`$-adic integrals. The material here will only be used in section 7. So we denote by $`\widehat{}`$ the completion of $`_{\mathrm{loc}}`$ with respect to the filtration $`F^m_{\mathrm{loc}}`$ where $`F^m_{\mathrm{loc}}`$ is the subgroup generated by $`\{\left[S\right]𝐋^i:idimSm\}`$. We will also denote by $`F^{}`$ the filtration induced on $`\widehat{}`$. In and the following is shown<sup>1</sup><sup>1</sup>1in fact loc. cit. also covers the case of singular varieties: There exists a unique map $`\mu :𝐁^t\widehat{}`$ satisfying the following three properties. 1. If $`A𝐁^t`$ is stable at level $`n`$, then $`\mu (A)=[\pi _n(A)]𝐋^{(n+1)d}`$. 2. If $`A𝐁^t`$ is contained in $`(S)`$ with $`S`$ a reduced closed subscheme of $`Xk[[t]]`$ with $`\mathrm{dim}_{k[[t]]}S<\mathrm{dim}X`$, then $`\mu (A)=0`$. 3. Let $`A_i`$ be in $`𝐁^t`$ for each $`i`$ in $``$. Assume that the $`A_i`$’s are mutually disjoint and that $`A:=_iA_i`$ is $`k[[t]]`$-semi-algebraic. Then $`_i\mu (A_i)`$ converges in $`\widehat{}`$ to $`\mu (A)`$. Moreover we have: 1. If $`A`$ and $`B`$ are in $`𝐁^t`$, $`AB`$ and if $`\mu (B)F^m\widehat{}`$, then $`\mu (A)F^m\widehat{}`$. This unique map $`\mu `$ is called the motivic volume on $`(X)`$ and is denoted by $`\mu _{(X)}`$ or $`\mu `$. For $`A`$ in $`𝐁^t`$ and $`\alpha :A\{\mathrm{}\}`$ a $`k[[t]]`$-simple function, one defines the motivic integral (3.5.1) $$_A𝐋^\alpha 𝑑\mu :=\underset{n}{}\mu (A\alpha ^1(n))𝐋^n$$ in $`\widehat{}`$, whenever the right hand side converges in $`\widehat{}`$, in which case we say that $`𝐋^\alpha `$ is integrable on $`A`$. If the function $`\alpha `$ is bounded from below, then $`𝐋^\alpha `$ is integrable on $`A`$, because of (4). ### 3.6. Rationality The following results are proved in section 5 of for semi-algebraic families of semi-algebraic subsets and simple functions. For simplicity results are stated here only for smooth varieties. ###### 3.7 Theorem. Let $`X`$ be a smooth algebraic variety over $`k`$ of pure dimension $`d`$. Let $`A_n`$, $`n^r`$, be a semi-algebraic family of semi-algebraic subsets of $`(X)`$ and let $`\alpha :(X)\times ^r`$ be a simple function. Assume that $`A_n`$ and the fibers of $`\alpha (\mathrm{\_},n):A_n`$ are stable, for every $`n^r`$. Then the power series $$\underset{n^r}{}T^n_{A_n}𝐋^{\alpha (\mathrm{\_},n)}𝑑\stackrel{~}{\mu }$$ in the variable $`T=(T_1,\mathrm{},T_r)`$ belongs to the subring of $`_{\mathrm{loc}}[[T]]`$ generated by $`_{\mathrm{loc}}[T]`$ and the series $`(1𝐋^aT^b)^1`$, with $`a`$ and $`b^r\{0\}`$. ###### 3.8 Theorem. Let $`X`$ be a smooth algebraic variety over $`k`$ of pure dimension $`d`$. Let $`A_n`$, $`n^r`$, be a semi-algebraic family of semi-algebraic subsets of $`(X)`$ and let $`\alpha :(X)\times ^r`$ be a simple function. Then the power series $$\underset{n^r}{}T^n_{A_n}𝐋^{\alpha (\mathrm{\_},n)}𝑑\mu $$ in the variable $`T=(T_1,\mathrm{},T_r)`$ belongs to the subring of $`\widehat{}[[T]]`$ generated by the image in $`\widehat{}[[T]]`$ of $`_{\mathrm{loc}}[T]`$, $`(𝐋^i1)^1`$ and $`(1𝐋^aT^b)^1`$, with $`i\{0\}`$, $`a`$, $`b^r\{0\}`$. ###### 3.9 Corollary. For any semi-algebraic subset $`A`$ of $`(X)`$, the measure $`\mu (A)`$ is in $`\overline{}_{\mathrm{loc}}[((𝐋^i1)^1)_{i1}]`$, where $`\overline{}_{\mathrm{loc}}`$ is the image of $`_{\mathrm{loc}}`$ in $`\widehat{}`$ ###### 3.10 Remark. By replacing $`T_i`$ by $`T_i𝐋^{m_i}`$ in Theorems 3.7 and 3.8, one sees that the condition “$`\alpha `$ takes values in $``$” may be replaced by the condition “$`\alpha `$ is bounded from below by a linear function of the $`^r`$-variable”. ###### 3.11 Remark. It seems quite likely that Theorems 3.7 and 3.8 remain true if semi-algebraic is replaced everywhere by $`k[[t]]`$-semi-algebraic and simple by $`k[[t]]`$-simple. However it does not seem that the proofs may be adapted directly to that more general situation, since there might be some “bad reduction at $`t=0`$”. ## 4. Motivic integration on the infinite Grassmannian All the material in 4.1 and 4.2 is contained in , or is directly adapted, replacing $`\mathrm{SL}`$ by $`\mathrm{GL}`$, from statements in . ### 4.1. We work over a field $`k`$ of characteristic $`0`$. We shall consider the infinite Grassmannian as a functor on the category of $`k`$-algebras. (A $`k`$-algebra will always be assumed to be associative, commutative and unitary.) A natural framework is that of $`k`$-spaces and $`k`$-groups in the sense of . By definition, a $`k`$-space (resp. a $`k`$-group) is a functor from the category of $`k`$-algebras to the category of sets (resp. of groups) which is a sheaf for the faithfully flat topology (see for a definition). The category of schemes over $`k`$ can be viewed as a full subcategory of the category of $`k`$-spaces. Schemes will always be assumed to be quasi-compact and quasi-separated. An important feature is that direct limits exist in the category of $`k`$-spaces, so we can say a $`k`$-space (resp. a $`k`$-group) is an ind-scheme (resp. an ind-group) if it is the direct limit of a directed system of schemes (resp. of group schemes). We fix a positive integer $`d`$. Let $`t`$ be an indeterminate. We consider the $`k`$-groups $`\mathrm{𝐆𝐋}_d(k[[t]])`$ and $`\mathrm{𝐆𝐋}_d\left(k((t))\right)`$ respectively defined by $`R\mathrm{GL}_d(R[[t]])`$ and $`R\mathrm{GL}_d\left(R((t))\right)`$. For $`n0`$, we denote by $`G_{(n)}(R)`$ the set of matrices $`A(t)`$ in $`\mathrm{GL}_d\left(R((t))\right)`$ such that both $`A(t)`$ and $`A(t)^1`$ have a pole of order $`n`$. This defines a subfunctor $`G_{(n)}`$ of $`\mathrm{𝐆𝐋}_d(k((t)))`$. One can show (cf. ), that the $`k`$-group $`\mathrm{𝐆𝐋}_d(k[[t]])`$ is an affine group scheme and that the $`k`$-group $`\mathrm{𝐆𝐋}_d(k((t)))`$ is an ind-group, being the direct limit of the sequence of schemes $`(G_{(n)})`$, $`n0`$. ### 4.2. The infinite Grassmannian $`𝒢r`$ may be defined as the quotient $$\mathrm{𝐆𝐋}_d(k((t)))/\mathrm{𝐆𝐋}_d(k[[t]])$$ in the category of $`k`$-spaces. It is the sheaf, for the faithfully flat topology, associated to the presheaf $`R\mathrm{GL}_d(R((t)))/\mathrm{GL}_d(R[[t]])`$. Now we have to explain why $`𝒢r`$ is indeed the infinite Grassmannian. For any $`k`$-algebra $`R`$, consider the set $`W(R)`$ of $`R[[t]]`$-submodules $`L`$ of $`R((t))^d`$ such that, for some $`n0`$, $`t^nR[[t]]^dLt^nR[[t]]^d`$ and $`L/t^nR[[t]]^d`$ is a projective $`R`$-module. By Proposition 2.3 of , the $`k`$-space $`𝒢r`$ is isomorphic to the functor $`RW(R)`$. Under this isomorphism, the group scheme action of $`\mathrm{𝐆𝐋}_d(k((t)))`$ on $`𝒢r`$ corresponds to the natural action of $`\mathrm{GL}_d(R((t)))`$ on $`R[[t]]`$-submodules of $`R((t))^d`$. Denote by $`𝒢r_{(n)}`$ the image of $`G_{(n)}`$ in $`𝒢r`$. Under the preceding isomorphism $`𝒢r_{(n)}`$ may be identified with the Grassmannian $`\mathrm{Gr}_t(t^nk[[t]]^d/t^nk[[t]]^d)`$, whose $`K`$-rational points, for $`K`$ a field containing $`k`$, are the $`K`$-linear subspaces of the finite dimensional $`K`$-vector space $`t^nK[[t]]^d/t^nK[[t]]^d`$ which are stable by multiplication by $`t`$; in particular $`𝒢r_{(n)}`$ is a projective variety. Furthermore, $`𝒢r`$ as a $`k`$-space is naturally isomorphic to the direct limit of the system of projective varieties $`𝒢r_{(n)}`$, hence $`𝒢r`$ is an ind-scheme. The canonical morphism of $`k`$-spaces $`\mathrm{\Theta }:\mathrm{𝐆𝐋}_d(k((t)))𝒢r`$ is a locally trivial fibration for the Zariski topology, i.e. $`𝒢r`$ is covered by open subsets over which $`\mathrm{\Theta }`$ is a product. ### 4.3. Denote by $`\mathrm{M}_d`$ the affine $`k`$-space of $`d`$ by $`d`$ matrices. There is canonical immersion of $`k`$-schemes $`\iota :\mathrm{𝐆𝐋}_d(k[[t]])(\mathrm{M}_d),`$ which identifies $`\mathrm{𝐆𝐋}_d(k[[t]])`$ with the open subscheme defined by $`\mathrm{det}0`$. In particular we can consider $`\mathrm{𝐆𝐋}_d(k[[t]])`$ as a semi-algebraic subset of $`(\mathrm{M}_d)`$. A subset of $`\mathrm{𝐆𝐋}_d(k[[t]])`$ will be called semi-algebraic (resp. $`k[[t]]`$-semi-algebraic) if it is semi-algebraic (resp. $`k[[t]]`$-semi-algebraic) as a subset of $`(\mathrm{M}_d)`$. The space $`\mathrm{𝐆𝐋}_d(k((t)))`$ being an ind-scheme, one can associate to it an underlying set, which by abuse we shall denote by the same letter. We shall say a subset $`A`$ of $`\mathrm{𝐆𝐋}_d(k((t)))`$ is bounded $`k[[t]]`$-semi-algebraic (resp. stable bounded $`k[[t]]`$-semi-algebraic, resp. cylindrical), if, for some integer $`n`$ in $``$, $`t^nA`$ is a $`k[[t]]`$-semi-algebraic (resp. stable $`k[[t]]`$-semi-algebraic, resp. cylindrical) subset of $`\mathrm{𝐆𝐋}_d(k[[t]])`$. We may extend the measures $`\mu `$ and $`\stackrel{~}{\mu }`$ to bounded $`k[[t]]`$-semi-algebraic, stable bounded $`k[[t]]`$-semi-algebraic and cylindrical subsets respectively, by defining $$\mu (t^nA)=𝐋^{d^2n}\mu (A)\text{and}\stackrel{~}{\mu }(t^nA)=𝐋^{d^2n}\stackrel{~}{\mu }(A),$$ for $`A`$ a $`k[[t]]`$-semi-algebraic (resp. stable $`k[[t]]`$-semi-algebraic or cylindrical) subset of $`\mathrm{𝐆𝐋}_d(k[[t]])`$, which is independent from the choice of the integer $`n`$. One should also remark that the measures $`\mu `$ and $`\stackrel{~}{\mu }`$ are invariant under $`\mathrm{𝐆𝐋}_d(k[[t]])`$-action. ### 4.4. For any integer $`m`$ in $``$, the functor $$R\{M\mathrm{GL}_d(R((t))):\mathrm{ord}_t\mathrm{det}M=m\}$$ defines a subspace of $`\mathrm{𝐆𝐋}_d(k((t)))`$, which we shall denote by $`\mathrm{𝐆𝐋}_d(k((t)))[m]`$. Clearly $`\mathrm{𝐆𝐋}_d(k((t)))[m]`$ is an ind-scheme. If $`H`$ is a linear subspace of $`t^nK[[t]]^d/t^nK[[t]]^d`$, for $`K`$ a field containing $`k`$, we set (4.4.1) $$\mathrm{index}(H):=dimt^nK[[t]]^d/Hnd.$$ We define $`𝒢r_{(n)}[m]`$ as the projective variety parametrizing subspaces in the Grassmannian $`\mathrm{Gr}_t(t^nk[[t]]^d/t^nk[[t]]^d)`$ which are of index $`m`$. For fixed $`m`$, the varieties $`𝒢r_{(n)}[m]`$ form an inductive system and we denote by $`𝒢r[m]`$ the corresponding ind-scheme. For every $`m`$, the fibration $`\mathrm{\Theta }:\mathrm{𝐆𝐋}_d(k((t)))𝒢r`$ restricts to a fibration $`\mathrm{\Theta }:\mathrm{𝐆𝐋}_d(k((t)))[m]𝒢r[m]`$. ### 4.5. Let $`A`$ be a subset of $`𝒢r`$. We say $`A`$ is gr-stable at level $`n`$ if $`A`$ is a constructible subset of $`𝒢r_{(n)}`$. Note that in this case $`A𝒢r[m]`$ is constructible in $`𝒢r_{(n)}[m]`$ for every $`m`$. Hence we can define the motivic measure $`\stackrel{~}{\mu }_{𝒢r}(A)`$ of $`A`$ as the element $$\stackrel{~}{\mu }_{𝒢r}(A):=\underset{m}{}\frac{[A𝒢r[m]]}{𝐋^{md}}$$ in $`_{\mathrm{loc}}`$, which makes sense, the sum on the right hand side being finite. Clearly if $`A`$ is gr-stable at level $`n`$ then $`\mathrm{\Theta }^1(t^nA)`$ is cylindrical at level $`2n`$. The relation with the previously defined measure $`\stackrel{~}{\mu }`$ is given by the following Proposition. ###### 4.6 Proposition. Let $`A`$ be a gr-stable subset of $`𝒢r`$. Then $$\stackrel{~}{\mu }_{𝒢r}(A)=(1𝐋^1)^1\mathrm{}(1𝐋^d)^1\stackrel{~}{\mu }(\mathrm{\Theta }^1(A)).$$ ###### Proof. We may assume $`A`$ is contained in $`𝒢r_{(n)}[m]`$. By Lemma 4.7 applied to $`r=m+nd`$ and $`p=2n`$, $$[\pi _{2n}(\mathrm{\Theta }^1(t^nA))]=[t^nA](𝐋^d1)\mathrm{}(𝐋^d𝐋^{d1})𝐋^{2nd^2d(m+nd)},$$ since $`[\mathrm{GL}_{d,k}]=(𝐋^d1)\mathrm{}(𝐋^d𝐋^{d1})`$. We deduce that $$\stackrel{~}{\mu }(\mathrm{\Theta }^1(t^nA))=[t^nA](1𝐋^1)\mathrm{}(1𝐋^d)𝐋^{nd^2}𝐋^{md}.$$ The result follows since $$\stackrel{~}{\mu }(\mathrm{\Theta }^1(t^nA))=\stackrel{~}{\mu }(t^n\mathrm{\Theta }^1(A))=𝐋^{nd^2}\stackrel{~}{\mu }(\mathrm{\Theta }^1(A))$$ and $`[t^nA]=[A]`$.∎ ###### 4.7 Lemma. For any integer $`p0`$, the morphism $$\pi _p(\mathrm{\Theta }^1(\mathrm{Gr}_t(k[[t]]^d/t^pk[[t]]^d)[r]))\mathrm{Gr}_t(k[[t]]^d/t^pk[[t]]^d)[r]$$ is a locally trivial fibration for the Zariski topology with fiber $`\mathrm{GL}_{d,k}\times 𝔸_k^{pd^2dr}`$. Here $`\mathrm{GL}_{d,k}`$ denotes the algebraic variety of invertible $`d`$ by $`d`$ matrices over $`k`$. ###### Proof. By taking a cover of the Grassmannian $`\mathrm{Gr}_t(k[[t]]^d/t^pk[[t]]^d)[r]`$ by open Schubert cells corresponding to different choices of bases of the lattice $`k[[t]]^d`$, one deduces the result from the following elementary Lemma 4.8.∎ ###### 4.8 Lemma. Let $`pm_d\mathrm{}m_10`$ be integers. Set $`r=_{1id}m_i`$. Let $`U`$ be the subscheme of $`(\mathrm{M}_d)`$ consisting of triangular matrices $`(a_{ij})`$ with $`a_{ii}=t^{m_i}`$, $`a_{ij}=0`$ for $`j<i`$, $`m_i\mathrm{ord}_ta_{ij}`$ and $`\mathrm{deg}_ta_{ij}<m_j`$ for $`j>i`$. Consider the morphism $`\phi :U\times \mathrm{𝐆𝐋}_d(k[[t]])U\times _p(\mathrm{M}_d)`$ sending $`(A,M)`$ to $`(A,\pi _p(AM))`$. Then the image $`W`$ of $`\phi `$ is an algebraic variety over $`k`$ which is isomorphic to $`U\times \mathrm{GL}_{d,k}\times 𝔸_k^{pd^2dr}`$ by an isomorphism compatible with the projection $`WU`$. Now we can express our zeta functions as integrals, defined as in 3.4.1, with respect to the measure $`\stackrel{~}{\mu }_{𝒢r}`$ on the Grassmannian: ###### 4.9 Theorem. Let $`𝒳`$ be a class of constructible $`k[[t]]`$-subalgebras. Then $$P_{L,𝒳}(𝐋^s)=_{𝒳\mathrm{Gr}_0^d(L)}𝐋^{\mathrm{index}(H)d\left(\mathrm{codim}H\right)s}𝑑\stackrel{~}{\mu }_{𝒢r}(H).$$ ###### Proof. Let $`𝒳_n`$ be the subalgebras of codimension $`n`$. Then $`𝒳_n\mathrm{Gr}^d(t^nL/t^nL)`$ and $`\left[𝒳_n\right]`$ $`=`$ $`{\displaystyle \underset{m}{}}\left[𝒳_n𝒢r[m]\right]`$ $`=`$ $`{\displaystyle \underset{m}{}}{\displaystyle \frac{\left[𝒳_n𝒢r[m]\right]}{𝐋^{md}}}𝐋^{md}`$ $`=`$ $`{\displaystyle \underset{m}{}}{\displaystyle _{𝒳_n𝒢r[m]}}𝐋^{\mathrm{index}(H)d}𝑑\stackrel{~}{\mu }_{𝒢r}(H)`$ $`=`$ $`{\displaystyle _{𝒳_n}}𝐋^{\mathrm{index}(H)d}𝑑\stackrel{~}{\mu }_{𝒢r}(H).`$ ## 5. Rationality of motivic zeta functions of infinite dimensional Lie algebras In this section we shall prove the following rationality result: ###### 5.1 Theorem. Let $`k`$ be a field of characteristic zero. 1. Let $`L`$ be a finite dimensional free $`k[[t]]`$-Lie algebra of the form $`L=L_k_kk[[t]]`$, with $`L_k`$ a Lie algebra over $`k`$. Let $`𝒳_t^{}`$ (resp. $`𝒳_t^{}`$) be the class of subalgebras such that $`𝒳_t^{}(K)`$ (resp. $`𝒳_t^{}(K)`$) is the set of all $`K[[t]]`$-subalgebras (resp. $`K[[t]]`$-ideals) of $`LK`$, for every field $`K`$ which is a finite extension of $`k`$. Then $`P_{L,𝒳_t^{}}(T)`$ for $`\{,\}`$ is rational, belonging to $`[T]_{\mathrm{loc}}`$. 2. Let $`L`$ be a finite dimensional $`k((t))`$-Lie algebra and $`L_0`$ be a choice of some $`k[[t]]`$-Lie subalgebra of the form $`L_0=L_k_kk[[t]]`$, with $`L_k`$ a Lie algebra over $`k`$. Let $`𝒳_{t,0}^{}`$ be the class of subalgebras such that $`𝒳_{t,0}^{}(K)`$ is the set of all $`K[[t]]`$-subalgebras $`LK`$ commensurable with $`L_0K`$, for every field $`K`$ which is a finite extension of $`k`$. Then $`P_{L,𝒳_{0,t}^{}}(T)`$ is rational, belonging to $`[T]_{\mathrm{loc}}`$. Let $`L`$ be a finite dimensional free $`k[[t]]`$-Lie algebra of dimension $`d`$. Choosing a basis allows us to identify $`L`$ additively with $`k[[t]]^d`$ and we may view the $`𝒳_t^{}`$ as a subset of the Grassmannian $`𝒢r`$. ###### 5.2 Theorem. With the preceding notations, $`\mathrm{\Theta }^1(𝒳_t^{})`$ is a $`k[[t]]`$-semi-algebraic subset of $`\mathrm{𝐆𝐋}_d(k[[t]])`$, for $`\{,\}`$, and (5.2.1) $`P_{L,𝒳_t^{}}(𝐋^s)`$ $`=`$ $`{\displaystyle _{𝒳_t^{}}}𝐋^{(ds)\mathrm{codim}H}𝑑\stackrel{~}{\mu }_{𝒢r}(H)`$ (5.2.2) $`=`$ $`(1𝐋^1)^1\mathrm{}(1𝐋^d)^1{\displaystyle _{\mathrm{\Theta }^1(𝒳_t^{})}}𝐋^{\left(ds\right)\mathrm{ord}_tM}𝑑\stackrel{~}{\mu }.`$ Furthermore, if $`L=L_k_kk[[t]]`$, with $`L_k`$ a Lie algebra over $`k`$, then $`\mathrm{\Theta }^1(𝒳_t^{})`$ is a semi-algebraic subset of $`\mathrm{𝐆𝐋}_d(k[[t]])`$. ###### Proof. The first line (5.2.1) follows from Theorem 4.9 and the remark that $`\mathrm{codim}H=\mathrm{index}(H)`$ if $`H`$ is a subalgebra of $`L.`$ Let $`A^{}=\mathrm{\Theta }^1(𝒳_t^{}).`$ Let $`M=(m_{ij})`$ be in $`\mathrm{𝐆𝐋}_d(k[[t]])`$ and write $`𝐦_i=m_{i1}e_1+\mathrm{}+m_{id}e_d,`$ for $`i=1,\mathrm{},d`$. We have $`MA^{}`$ if and only if, for every $`1i,jd`$, there exist $`Y_{ij}^1,\mathrm{},Y_{ij}^dK[[t]]`$ such that $`\beta (𝐦_i,𝐦_j)=_{k=1}^dY_{ij}^k𝐦_k`$, for some finite field extension $`K`$ of $`k`$. Here $`\beta `$ is the bilinear mapping $`L\times LL`$ corresponding to the product in $`L`$. Let $`C_j`$ denote the matrix whose rows are $`𝐜_i=\beta (e_i,e_j)`$. We then have $$\beta (𝐦_i,𝐦_j)=𝐦_i\left(\underset{l=j}{\overset{d}{}}m_{jl}C_l\right).$$ Then $`MA^{}`$ if and only if, for every $`1i,jd`$, one can solve the matrix equation (5.2.3) $$𝐦_i\left(\underset{l=j}{\overset{d}{}}m_{jl}C_l\right)=(Y_{ij}^1,\mathrm{},Y_{ij}^d)M$$ for $`Y_{ij}^1,\mathrm{},Y_{ij}^dK[[t]]`$, with $`K`$ a finite field extension of $`k`$. Let $`M^{}`$ denote the adjoint matrix, we can then rewrite (5.2.3) as $$𝐦_i\left(\underset{l=j}{\overset{d}{}}m_{jl}C_l\right)M^{}=(det(M)Y_{ij}^1,\mathrm{},det(M)Y_{ij}^d).$$ Let $`g_{ijk}(m_{rs})`$ denote the $`k`$-th entry of the $`d`$-tuple $`𝐦_i\left(_{l=j}^dm_{jl}C_l\right)M^{}`$. Then the set $`A^{}`$ has the following description: (5.2.4) $$A^{}=\{(m_{ij})(M_d):\mathrm{ord}_t(det(M))\mathrm{ord}_t(g_{ijk}(m_{rs}))\text{ for }i,j,k\{1,\mathrm{},d\}\}.$$ The set $`A^{}`$ is therefore $`k[[t]]`$-semi-algebraic. Let $`\mathrm{\Phi }^{}(M)`$ denote the conjunction of conditions $`\mathrm{ord}_t(det(M))\mathrm{ord}_t(g_{ijk}(m_{rs})).`$ Since $`\mathrm{ord}_t:(M_d)\left\{\mathrm{}\right\}`$ is a simple function this implies that $`A_n`$ is a semi-algebraic set and in particular is constructible. As we promised in section 2 this provides an alternative way to show the constructibility of this class of subalgebras. We prove that $`A^{}`$ is $`k[[t]]`$-semi-algebraic. We have $`MA^{}`$ if and only if, for every $`1i,jd`$, there exist $`Y_{ij}^1,\mathrm{},Y_{ij}^dK[[t]]`$, with $`K`$ some finite field extension of $`k`$, such that $`\beta (𝐦_i,𝐞_j)=_{k=1}^dY_{ij}^k𝐦_k.`$ Let $`g_{ijk}^{}(m_{rs})`$ denote the $`k`$-th entry of the $`d`$-tuple $`𝐦_iC_jM^{}`$. Then an argument similar to the above implies that the set $`A^{}`$ has the following description: $$A^{}=\{(m_{ij})(M_d):\mathrm{ord}_t(det(M))\mathrm{ord}_t(g_{ijk}^{}(m_{rs}))\text{ for }i,j,k\{1,\mathrm{},d\}\}.$$ Hence $`A^{}`$ is $`k[[t]]`$-semi-algebraic. Indeed it is defined by $`\mathrm{\Phi }^{}`$, the definable condition which is the conjunction of conditions $`\mathrm{ord}_t(det(M))\mathrm{ord}_t(g_{ijk}^{}(m_{rs}))`$. Equation (5.2.2) follows from Proposition 4.6 and the last statement is clear, since when $`L`$ is of the form $`L_k_kk[[t]]`$, the above definable conditions are all semi-algebraic. ∎ ###### 5.3 Remark. The following provides an easier integral expression in general to calculate. It is similar to the integral expression that was used in . Let $`Y`$ denote the variety $`\mathrm{Tr}_{d,k}`$ of $`d\times d`$ upper triangular matrices which is isomorphic to the affine space $`𝔸_k^{d(d+1)/2}`$. Let $`\stackrel{~}{\nu }`$ denote the measure on stable $`k[[t]]`$-semi-algebraic subsets of $`(Y)`$. For $`i=1,\mathrm{},d`$ define the simple function $`\alpha _i:(Y)`$ by $`\alpha _i(M)=\mathrm{ord}_t(m_{ii}).`$ Then $`(1𝐋^1)^1\mathrm{}(1𝐋^d)^1{\displaystyle _A^{}}𝐋^{\left(ds\right)\mathrm{ord}_t(M)}𝑑\stackrel{~}{\mu }`$ $`=`$ $`(1𝐋^1)^d{\displaystyle _{A^{}(Y)}}𝐋^{\mathrm{ord}_t(M)s}𝐋^{\alpha _1(M)+\mathrm{}+i\alpha _i(M)+\mathrm{}+d\alpha _d(M)}𝑑\stackrel{~}{\nu }.`$ ### 5.4. Commensurable subalgebras We turn now to the issue of commensurable subalgebras. The following discussion will also apply to counting commensurable subalgebras in a $`_p`$-Lie algebra, an issue which has previously not been discussed. We identify $`L`$ with $`k((t))^d`$ by a choice of a basis $`e_1,\mathrm{},e_d`$ for the $`k[[t]]`$-submodule $`L_0`$ which is then identified with $`k[[t]]^d`$. The mapping $`\mathrm{\Theta }:\mathrm{𝐆𝐋}_d(k((t)))𝒢r`$ sends a matrix $`M`$ to the lattice $`𝐦_1,\mathrm{},𝐦_d`$ spanned by $$𝐦_i=m_{ii}e_i+\mathrm{}+m_{id}e_d,$$ for $`i=1,\mathrm{},d`$. The index then of $`H=𝐦_1,\mathrm{},𝐦_d`$ (as defined in (4.4.1)) is just $`\mathrm{ord}_t(detM).`$ ###### 5.5 Remark. For subalgebras of $`L_0K`$ the index of $`H`$ is the same as the codimension. We might have considered encoding a zeta function of commensurable subalgebras according to index but note that there are an infinite number of commensurable submodules of index $`0`$ whilst there is only one of codimension $`0.`$ We now give a description of the codimension as a function of the entries of $`M.`$ Recall that if the entries of $`M`$ are in the field $`K`$ containing $`k`$, the codimension of $`H=𝐦_1,\mathrm{},𝐦_d`$ in $`LK`$ is defined as the sum of the codimension of $`H(L_0K)`$ in $`L_0K`$ and in $`H`$: ###### 5.6 Lemma. Let $`M`$ be a $`K`$-rational point of $`\mathrm{𝐆𝐋}_d(k((t)))`$ and define $$n_i=\mathrm{min}\{0,\mathrm{ord}_t(m_{ij}):j=1,\mathrm{},d\}.$$ Then the codimension of $`H=𝐦_1,\mathrm{},𝐦_d`$ in $`LK`$ is the sum of $`\mathrm{ord}_\mathrm{t}(det(m_{ij}t^{n_i}))`$ and $`_{i=1}^dn_i`$. ###### Proof. This follows from the fact that $`H(L_0K)=t^{n_1}𝐦_1,\mathrm{},t^{n_d}𝐦_d`$. Hence the codimension of $`H(L_0K)`$ is given as usual by $`\mathrm{ord}_t(det(m_{ij}t^{n_i})).`$ The codimension of $`H(L_0K)`$ in $`H`$ on the other hand is given by $`_{i=1}^dn_i`$.∎ We define then the function $$\mathrm{codim}(M):=\mathrm{ord}_t(det(m_{ij}t^{n_i}))+\underset{i=1}{\overset{d}{}}n_i=\mathrm{ord}_t(det(m_{ij}))+2\underset{i=1}{\overset{d}{}}n_i$$ which is a simple function. ### 5.7. We want to convert the Poincaré series into integrals over arc spaces. To do this we think of $`k((t))`$ as consisting of two pieces: $`k[[t]]`$ identified with $`k[[t]]`$ via $`vv`$ and $`k((t))\backslash k[[t]]`$ identified with $`tk[[t]]`$ via $`vv^1`$. Notice that the measure on $`k((t))\backslash k[[t]]`$ translates under this identification to the measure $`𝐋^{2\mathrm{o}\mathrm{r}\mathrm{d}_t(v)}\mu _{k[[t]]}`$ on $`tk[[t]].`$ This follows because the set $`t^nk[[t]]^{}`$ has measure $`𝐋^n`$ whilst $`t^nk[[t]]^{},`$ its image under $`vv^1,`$ has measure $`𝐋^n.`$ We can partition $`\mathrm{𝐆𝐋}_d(k((t)))`$ into $`2^{d^2}`$ subsets $`\mathrm{𝐆𝐋}_d(k((t)))_A`$ which we can index with an element of $`AM_d(\left\{\pm 1\right\})`$ so that $$\mathrm{𝐆𝐋}_d(k((t)))_A=\{M\mathrm{𝐆𝐋}_d(k((t))):\mathrm{ord}_t(m_{ij})0\text{ if and only if }A_{ij}=1\}.$$ We then identify each $`\mathrm{𝐆𝐋}_d(k((t)))_A`$ with a subset of $`(\mathrm{M}_d)`$ by the morphism $`\mathrm{\Xi }_A:M\left(m_{ij}^{A_{ij}}\right)`$. Note that $`\mathrm{\Xi }_A\left(\mathrm{\Xi }_A\left(M\right)\right)=M.`$ Now remark, that with the notations used in the proof of Proposition 5.2, a matrix $`M`$ is in $`\mathrm{\Theta }^1(𝒳_0^{})`$ if and only if $`\mathrm{\Phi }^{}(M)`$ is true where $`\mathrm{\Phi }^{}`$ was defined as the conjunction of statements $`\mathrm{ord}_t(det(M))\mathrm{ord}_t(g_{ijk}(m_{rs}))`$ appearing in (5.2.4). Hence the subset $`\mathrm{\Xi }_A\left(\mathrm{\Theta }^1(𝒳_{t,0}^{})\mathrm{𝐆𝐋}_d(k((t)))_A\right)`$ is a $`k[[t]]`$-semi-algebraic subset of $`(\mathrm{M}_d)`$ defined by the condition $`\mathrm{\Phi }^{}(\mathrm{\Xi }_AM)`$ is true. Furthermore, when $`L_0`$ is of the form $`L_k_kk[[t]]`$, the definable conditions considered are all semi-algebraic, hence $`\mathrm{\Xi }_A\left(\mathrm{\Theta }^1(𝒳_{t,0}^{})\mathrm{𝐆𝐋}_d(k((t)))_A\right)`$ is semi-algebraic in this case. Hence we have the following result: ###### 5.8 Theorem. With the preceding notations, $`\mathrm{\Xi }_A\left(\mathrm{\Theta }^1(𝒳_{t,0}^{})\mathrm{𝐆𝐋}_d(k((t)))_A\right)`$ is a $`k[[t]]`$-semi-algebraic subset of $`(\mathrm{M}_d)`$, and $`P_{L,𝒳_{t,0}^{}}(𝐋^s)`$ $`=`$ $`{\displaystyle _{\mathrm{\Theta }^1(𝒳_{t,0}^{})}}𝐋^{\mathrm{ord}_t(detM)d\left(\mathrm{codim}M\right)s}𝑑\stackrel{~}{\mu }`$ $`=`$ $`{\displaystyle \underset{AM_d(\left\{\pm 1\right\})}{}}{\displaystyle _{\mathrm{\Xi }_A\left(\mathrm{\Theta }^1(𝒳_{t,0}^{})\mathrm{𝐆𝐋}_d(k((t)))_A\right)}}𝐋^{\mathrm{ord}_t(det\mathrm{\Xi }_A\left(M\right))d\left(\mathrm{codim}\mathrm{\Xi }_A\left(M\right)\right)s}`$ $`\times 𝐋^{_{A_{ij}=1}2\mathrm{o}\mathrm{r}\mathrm{d}_tm_{ij}^{A_{ij}}}d\stackrel{~}{\mu }.`$ Furthermore, if $`L_0=L_k_kk[[t]]`$, with $`L_k`$ a Lie algebra over $`k`$, then $$\mathrm{\Xi }_A\left(\mathrm{\Theta }^1(𝒳_{t,0}^{})\mathrm{𝐆𝐋}_d(k((t)))_A\right)$$ is a semi-algebraic subset of $`(M_d)`$. ###### Proof. Everything has already been proven, except for (5.8), which follows from Proposition 4.6. ∎ ### 5.9. Proof of Theorem 5.1 Let us first prove (1). Since, by Theorem 5.2, we have (5.2.2) the result follows from Theorem 3.7, Remark 3.10 and the last statement in Theorem 5.2. Similarly to deduce (2) from Theorem 5.8, Theorem 3.7 and Remark 3.10, one needs to check that there exists an integer $`N`$ such that $$N\mathrm{codim}\mathrm{\Xi }_A\left(M\right)>\mathrm{ord}_t(det\mathrm{\Xi }_A\left(M\right))d+\underset{A_{ij}=1}{}2\mathrm{o}\mathrm{r}\mathrm{d}_tm_{ij}^{A_{ij}}.$$ This follows from the fact that $`\mathrm{codim}\mathrm{\Xi }_A\left(M\right)>\mathrm{ord}_t(det\mathrm{\Xi }_A\left(M\right))`$ and if $`A_{ij}=1`$ then $`\mathrm{ord}_tm_{ij}^1<n_i.`$ Hence if we take $`N=3d`$ we are done. ∎ ###### 5.10 Remark. If Theorem 3.7 remains true when semi-algebraic is replaced everywhere by $`k[[t]]`$-semi-algebraic and simple by $`k[[t]]`$-simple, as suggested in Remark 3.11, then Theorem 5.1 remains true with the same proof, without assuming that $`L`$ is obtained by extension of scalars from a Lie algebra $`L_k`$ defined over $`k`$. ### 5.11. Commensurable subalgebras in $`p`$-adic Lie algebras Note that the analysis above can also be applied to any finite dimensional $`_p`$-Lie algebra $``$ with a choice of $`_p`$-lattice $`L`$ where we define the commensurable zeta function as: $$\zeta _{,L}(s)=\underset{H𝒳(,L)}{}\|L+H:LH\|^s$$ where $`𝒳(,L)`$ denotes theset of $`_p`$-subalgebras commensurable with $`L`$. This is the analogue of what we did above where the index $`\|L+H:LH\|`$ corresponds to the codimension since $`\mathrm{codim}H=\mathrm{codim}(L_0+H:LH)`$ in the $`k[[t]]`$-setting. We shall consider later, see Corollary 7.8, the question of zeta functions of $`L_p,`$ where $`L`$ is $``$-Lie algebra, but here the analysis applies to a Lie algebra with no underlying $``$-structure. So it is worth recording the following result for $`\zeta _{,L}(s),`$ proved as we said by the same analysis as above replacing $`𝐋`$ by $`p`$, $`k[[t]]`$ by $`_p`$, $`k((t))`$ by $`_p`$, and Theorem 3.7 by Theorem 3.2 of . ###### 5.12 Theorem. Let $``$ be a finite dimensional $`_p`$-Lie algebra with a choice of $`_p`$-lattice $`L`$ such that $`L_p=.`$ Then $`\zeta _{,L}(s)`$ is a rational function in $`p^s.`$ ### 5.13. $`𝔽_p[[t]]`$-Lie algebras We mention that there is another open rationality question which fits into this context. This concerns Lie algebras in characteristic $`p.`$ For example the following is conjectured in section 5 of : ###### 5.14 Conjecture. Let $`L`$ be a finite dimensional Lie algebra over $`𝔽_p[[t]]`$ additively isomorphic to $`𝔽_p[[t]]^d.`$ Then $$\zeta _L(s)=\underset{HL}{}\|L:H\|^s$$ is a rational function in $`p^s`$ where the sum is taken over $`𝔽_p[[t]]`$-subalgebras $`H.`$ ### 5.15. The techniques involved in proving such a result in characteristic zero (quantifier elimination and Hironaka’s resolution of singularities) do not exist at present in characteristic $`p`$. However, in the rationality was proved for one class of $`𝔽_p[[t]]`$-algebras. Take $`L`$ to be an algebra additively isomorphic to $`^d`$ and set $`L_p=L𝔽_p[[t]],`$ then one can prove that for almost all primes $`\zeta _{L_p}(s)`$ is a rational function in $`p^s.`$ In fact one gets that $`\zeta _{L_p}(s)=\zeta _{L_p}(s)`$ for almost all primes. The proof actually follows the theme of this current paper, that there is an integral over the additive Haar measure on $`\left(𝔽_p[[t]]\right)^d`$ representing this zeta function $`𝔽_p[[t]]`$ which formally looks the same as the integral representing $`\zeta _{L_p}(s).`$ Macintyre then observed that the calculation of this latter integral will carry over formally to the former essentially for any primes for which one does not divide by in this calculation. See also Theorem 8.3.2 of for much more general results along these lines. ### 5.16. If we are going to consider the possible rationality of some of the other motivic zeta functions considered in this paper then it is necessary to introduce the concept of polynomial growth for $`(L,𝒳)`$. This is going to be a necessary requirement if we are to prove that the zeta function $`P_{L,𝒳}(T)`$ can be expressed as an element of $`[T]_{\mathrm{loc}}`$. ###### 5.17 Definition. We say that $`(L,𝒳)`$ has polynomial growth if there exists some $`d`$ such, that for all $`n`$, $`A_n(𝒳)`$ is a constructible set and $$dimA_n(𝒳)dn.$$ Recall a constructible set is of dimension $`m`$, if it may be partioned into finitely many locally closed pieces with dimension $`m`$. It is a necessary condition that $`(L,𝒳)`$ have polynomial growth for the zeta function to be rational, because the coefficients of a rational series satisfy a linear recurrences of finite order. For example, when $`L`$ is a finite dimensional $`k[[t]]`$-algebra additively isomorphic to $`k[[t]]^d`$ and $`𝒳`$ the class of commensurable $`k[[t]]`$-subalgebras, then $`A_n(𝒳)`$ is contained in $`𝒢r_{(n)}`$ whose dimension is bounded by $`2nd^2`$, the dimension of $`\mathrm{M}_d(t^nk[[t]]/t^nk[[t]])`$. ###### 5.18 Problem. Characterize the $``$-filtered or $``$-filtered infinite dimensional Lie algebras for which $`𝒳^{}`$ has polynomial growth. ### 5.19. Polynomial growth If we consider the algebra $`𝔡^+:=_{j1}kd_j`$, where $`(d_i,d_j)=\left(ij\right)d_{i+j}`$ and $`𝒳`$ is the class of all subalgebras, then this does not have polynomial growth. However, if we define $`𝒳^d`$ to be the class of $`d`$-generated subalgebras, we find that $`(𝔡^+,𝒳^d)`$ does have polynomial growth. ###### 5.20 Problem. Calculate $`P_{𝔡^+,𝒳^d}(T)`$ for $`d2.`$ Is it rational? More generally, $`(L,𝒳^d)`$ has polynomial growth if $`L`$ is a well-covered Lie algebra (where recall a subalgebra of codimension $`n`$ in $`LK`$ contains $`L_{f(n)}K`$ for some function $`f(n)`$) and $`dim(L/L_{f(n)})n^c`$ for some fixed $`c`$. This is true for example for all the simple graded Lie algebras of finite growth. However we are unable to prove that the corresponding zeta function encoding these $`d`$-generated subalgebras is rational. ## 6. $`k[[t]]`$-powered nilpotent groups ### 6.1. In subgroups of finite index in a torsion-free finitely generated nilpotent group $`G`$ were encoded in zeta functions. It was shown there how, for almost all primes $`p`$, there is a one-to-one correspondence between subgroups of finite index in the pro-$`p`$ completion $`\widehat{G}_p`$ of $`G`$ and subalgebras of finite index in $`L_p`$ for an associated Lie algebra $`L`$ over $`.`$ The pro-$`p`$ completion has the structure of a $`_p`$-powered nilpotent group. We can use the same ideas to show that the motivic zeta functions counting subalgebras in a finite dimensional $`k[[t]]`$-Lie algebra are also encoding the subgroup structure of a $`k[[t]]`$-powered nilpotent group $`G`$ where subgroups are taken to be $`k[[t]]`$-closed subgroups. An $`R`$-powered nilpotent group is a nilpotent group with an action of the ring $`R`$ satisfying various axioms which can be found in Chapter 6 of P. Hall’s notes on nilpotent groups . We shall be interested in a certain type of $`R`$-powered nilpotent group which arises when $`R`$ is a binomial ring in the following way. Take a torsion-free finitely generated nilpotent group $`G`$ with a choice of Malcev basis $`x_1,\mathrm{},x_d.`$ The group $`G`$ can then be identified with the set of elements $`𝐱(𝐚)=x_1^{a_1}\mathrm{}x_d^{a_d}`$ with $`𝐚=(a_1,\mathrm{},a_d)^d.`$ In P. Hall proved that there exist polynomials over $``$ in suitably many variables which define mappings $`\lambda `$ $`:`$ $`^d\times Z`$ $`\mu `$ $`:`$ $`^d\times ^d^d`$ such that for $`𝐚_1`$ and $`𝐚_2^d`$ and $`k`$ $`𝐱(𝐚)^k`$ $`=`$ $`𝐱(\lambda (𝐚,k))`$ $`𝐱(𝐚_1)𝐱(𝐚_2)`$ $`=`$ $`𝐱(\mu (𝐚_1,𝐚_2)).`$ Since these polynomials over $``$ take integer values, they can be represented as $``$-linear combinations of binomial polynomials. Suppose $`R`$ is a binomial ring, i.e. $`\left(\genfrac{}{}{0pt}{}{r}{n}\right)`$ is a well-defined element of $`R`$ for $`rR`$ and $`n`$. Then we can use the polynomials $`\lambda `$ and $`\mu `$ to define the structure of an $`R`$-powered nilpotent group on the set of elements $`𝐱(𝐚)=x_1^{a_1}\mathrm{}x_d^{a_d}`$ with $`𝐚=(a_1,\mathrm{},a_d)R^d.`$ Therefore to every finitely generated torsion-free nilpotent group $`G`$ we can define a $`k[[t]]`$-powered group which we shall denote $`G^{k[[t]]}.`$ There exists a $``$-rational representation of $`G^{k[[t]]}`$ embedding it as a subgroup of $`\mathrm{Tr}_n^1(k[[t]])`$, the upper triangular matrix group of dimension $`n`$ with $`1`$’s on the diagonal. The image of $`G^{k[[t]]}`$ under the map $`\mathrm{log}`$ defines a nilpotent $`k[[t]]`$-subalgebra $`=\mathrm{log}G^{k[[t]]}`$ of $`\mathrm{Tr}_n^0(k[[t]])`$, the upper triangular matrix group of dimension $`n`$ with $`0`$’s on the diagonal. (Note that $`\mathrm{log}`$ is a polynomial map since $`\left(U1\right)^n=0`$ if $`U\mathrm{Tr}_n^1(k[[t]]).)`$ The Lie algebra $``$ is $`L_{}k[[t]]`$ where $`L_{}`$ is the Lie algebra corresponding to $`G^{}`$ under the Malcev correspondence. In general the image of $`G`$ under $`\mathrm{log}`$ is not a Lie algebra. However there is some integer $`f`$ such that $`L=\mathrm{log}G^f`$ is a $``$-Lie algebra. This is the Lie algebra mentioned in the first paragraph used by Grunewald, Segal and Smith to set up a one-to-one correspondence between subgroups in $`G^_p=\widehat{G}_p`$ and $`L_p`$ for all primes not dividing $`f`$ (see Theorem 4.1 of ). A subgroup $`H`$ of $`G^{k[[t]]}`$ is $`k[[t]]`$-closed if $`h^rH`$ for all $`rk[[t]]`$ and $`hH.`$ We shall define the codimension of such a subgroup as the length $`n`$ of the maximal chain of $`k[[t]]`$-closed subgroups $`H_i`$ $`0in`$ with $`H_i`$ contained in $`H_{i1}`$ with infinite index and $`H_0=G^{k[[t]]}`$ and $`H_n=H.`$ The claim is then that $`\mathrm{log}`$ gives a way of getting from $`k[[t]]`$-closed subgroups to $`k[[t]]`$-subalgebras to prove the following: ###### 6.2 Theorem. There is a one-to-one correspondence between the set of $`k[[t]]`$-closed subgroups of $`G^{k[[t]]}`$ of finite codimension and the set of $`k[[t]]`$-subalgebras of finite codimension in $``$ which preserves the codimension. ###### Proof. We need to prove that the image of a $`k[[t]]`$-closed subgroup $`H`$ under the map $`\mathrm{log}`$ is a $`k[[t]]`$-subalgebra. The proof of Lemma 1 in Chapter 6 of can be applied in this setting to prove that $`m\mathrm{log}H\mathrm{log}H.`$ But since $`H`$ is $`k[[t]]`$-closed and $`k[[t]]`$ we have $`h^{1/m}H.`$ Hence $`\mathrm{log}H\mathrm{log}H,`$ i.e. $`\mathrm{log}H`$ is additively closed. Once we have this then Corollary 3 of Chapter 6 of implies that $`\mathrm{log}H`$ is also closed under the Lie bracket by applying the inverse of the Campbell-Hausdorff series. We have implicitly used the identity $`r\mathrm{log}h=\mathrm{log}h^r`$ above when $`r`$. This is true in fact for all values of $`rk[[t]].`$ This follows from the fact that there exist polynomials $`f_1(𝐚),\mathrm{},f_d(𝐚)`$ such that $`\mathrm{log}x_1^{a_1}\mathrm{}x_d^{a_d}=f_1(𝐚)\mathrm{log}x_1+\mathrm{}+f_d(𝐚)\mathrm{log}x_d.`$ Since $`r\mathrm{log}h=\mathrm{log}h^r`$ is true for all $`hG`$ and $`r`$ we get that the following is a polynomial identity for all $`𝐚k[[t]]^d`$ and $`r`$ $$f_i(\lambda (𝐚,r))=rf_i(𝐚).$$ Hence this must be a formal identity of polynomials which is then true for all values of $`rk[[t]].`$ Conversely, if $`M`$ is a $`k[[t]]`$-subalgebra of $``$ then $`\mathrm{exp}M`$ is well-defined on $`\mathrm{Tr}_n^0(k[[t]])`$ being a finite polynomial map. By using the Campbell-Hausdorff series one obtains directly that $`\mathrm{exp}M`$ is a subgroup and the identity $`r\mathrm{log}h=\mathrm{log}h^r`$ again implies that $`\mathrm{exp}M`$ is a $`k[[t]]`$-closed subgroup taking $`h=\mathrm{exp}m`$ and using the identity $`\mathrm{log}\mathrm{exp}(m)=m.`$ Since $`\mathrm{log}`$ and $`\mathrm{exp}`$ are bijective maps, we have a one-to-one correspondence as detailed in the statement of the Theorem. The correspondence preserves codimension since this is defined by lengths of chains of $`k[[t]]`$-closed subgroups or $`k[[t]]`$-subalgebras. ∎ We can use this to deduce the existence and rationality of a motivic zeta function associated to $`G^{k[[t]]},`$ a certain type of loop group. ###### 6.3 Corollary. Let $`G`$ be a finitely generated torsion-free nilpotent group and $`G^{k[[t]]}`$ the associated $`k[[t]]`$-powered nilpotent group. For any field $`K`$ which is a finite extension of $`k`$, let $`𝒳(K)`$ denote the set of $`K[[t]]`$-closed subgroups of $`G^{K[[t]]}`$ and let $`A_n(K)`$ denote the set of $`K[[t]]`$-closed subgroups of $`G^{K[[t]]}`$ of codimension $`n`$. There is a natural structure of constructible set on $`A_n=_KA_n(K)`$, hence $$P_{G^{k[[t]]},𝒳}(T):=\underset{n}{}[A_n]T^n$$ is a well-defined element of $`[[T]].`$ We call it the motivic zeta function of $`G^{k[[t]]}.`$ It is equal to the motivic zeta function $`P_{,𝒳_t^{}}(T)`$ of the associated $`k[[t]]`$-Lie algebra $`=\mathrm{log}G^{k[[t]]}=\mathrm{log}G^{}_{}k[[t]]`$. Hence $`P_{G^{k[[t]]},𝒳}(T)`$ is rational, belonging to $`[T]_{\mathrm{loc}}`$. We have a similar result for normal $`k[[t]]`$-closed subgroups in $`G^{k[[t]]}`$ since the one-to-one correspondence in Theorem 6.2 sends normal subgroups to ideals. ## 7. Motivic and $`p`$-adic cone integrals The motivic integrals we have defined to capture the subalgebras in a $`k[[t]]`$-algebra are examples of integrals that the first author and Grunewald studied in the $`p`$-adic setting in called cone integrals. We make the same definition in the motivic setting: ###### 7.1 Definition. We call an integral $$Z_{𝒟,\mathrm{geom}}(s):=_V𝐋^{\mathrm{ord}_t(f_0)s\mathrm{ord}_t(g_0)}𝑑\mu $$ a motivic cone integral if $`f_0(𝐱)`$ and $`g_0(𝐱)`$ are polynomials in $`k[[t]][x_1,\mathrm{},x_m]`$ and there exist polynomials $`f_i(𝐱),g_i(𝐱)`$, $`i=1,\mathrm{},l`$, in $`k[[t]][x_1,\mathrm{},x_m]`$ such that $$V=\{𝐱(𝔸_k^m):\mathrm{ord}_t(f_i(𝐱))\mathrm{ord}_t(g_i(𝐱))\text{ for }i=1,\mathrm{},l\}.$$ The set $`𝒟=\{f_0,g_0,f_1,g_1,\mathrm{},f_l,g_l\}`$ is called the set of cone integral data. ###### 7.2 Definition. Let $`L`$ be a $`k[[t]]`$-algebra additively isomorphic to $`\left(k[[t]]\right)^d`$ (or a $``$-Lie algebra additively isomorphic to $`^d).`$ Then we define $`Z_{L,\mathrm{geom}}(s)`$, the motivic cone integral associated to $`L`$, to be $`Z_{L,\mathrm{geom}}(s)`$ $`:=`$ $`{\displaystyle _A}𝐋^{s\mathrm{ord}_t(M)}𝑑\mu `$ $`=`$ $`Z_{𝒟,\mathrm{geom}}(s)`$ where $$A=\{(m_{ij})(M_d):\mathrm{ord}_t(det(M))\mathrm{ord}_t(g_{ijk}(m_{rs}))\text{ for }i,j,k\{1,\mathrm{},d\}\}$$ and $`𝒟`$ is the corresponding cone data. In section 5 we therefore established that $$P_{L,𝒳_t^{}}(𝐋^s)=(1𝐋^1)^1\mathrm{}(1𝐋^d)^1Z_{L,\mathrm{geom}}(sd).$$ The set $`V`$ in the definition of the motivic cone integral is $`k[[t]]`$-semi-algebraic and therefore measurable; the functions $`\mathrm{ord}_t(f_0)`$ and $`\mathrm{ord}_t(g_0)`$ are $`k[[t]]`$-simple functions. However, in this more general setting, we don’t know whether the fibres of $`\mathrm{ord}_t(f_0):V`$ and $`\mathrm{ord}_t(g_0):V`$ are weakly stable. This is why we need to use the completion $`\widehat{}`$ of $`_{\mathrm{loc}}`$ with respect to the filtration $`F^m_{\mathrm{loc}}`$, introduced in 3.5, and to replace the measure $`\stackrel{~}{\mu }`$ by the measure $`\mu `$. Hence, Theorem 3.8 gives us the following: ###### 7.3 Theorem. Assume the polynomials $`f_i(𝐱),g_i(𝐱)`$, $`i=0,\mathrm{},l,`$ have their coefficients in $`k`$. The motivic cone integral $`Z_{𝒟,\mathrm{geom}}(s)`$ belongs to the subring $`\overline{𝒩}`$ of the ring $`\widehat{}[[T]]`$ which is generated by the image in $`\widehat{}[[T]]`$ of $`_{\mathrm{loc}}[T]`$, $`(𝐋^i1)^1`$ and $`(1𝐋^aT^b)^1`$ with $`i`$, $`a`$, $`b\{0\}`$, where $`T=𝐋^s`$. In the algebra setting we had that the cone conditions were in fact stable and hence, by Theorem 3.7, the associated zeta function $`Z_{L,\mathrm{geom}}(s)`$ is an element of the subring of $`_{\mathrm{loc}}[[T]]`$, where $`T=𝐋^s.`$ However in the more general setting here, we are required to take infinite sums of elements of $`_{\mathrm{loc}}`$ which are then only defined in the associated complete ring $`\widehat{}.`$ At the moment it is not known if the image of $`_{\mathrm{loc}}`$ is the same as $`_{\mathrm{loc}}`$ since there might be some kernel. The infinite sums of elements of $`_{\mathrm{loc}}`$ and the terms like $`(𝐋^i1)^1`$ which don’t occur in Theorem 3.7 appear for example when we consider the constant term of these motivic zeta functions which take the following form: $$_{V\{𝐱:\mathrm{ord}_t(f_0(𝐱))=0\}}𝐋^{\mathrm{ord}_t(g_0)}𝑑\mu .$$ Even if $`g_0=0`$ we get terms coming from the cone conditions like $`(𝐋^i1)^1`$ as we shall see when we derive an explicit expression for motivic cone integrals. In the first author and Grunewald produced an explicit formula for $`p`$-adic cone integrals extending Denef’s explicit formula for the Igusa zeta function. We can do the same thing for the motivic cone integrals. As we shall explain, this will impact on how canonical the explicit formula for $`p`$-adic cone integrals is. It will also allow us to define a topological zeta function associated to $`p`$-adic cone integrals and ultimately to the zeta function of an algebra defined over $`.`$ ### 7.4. Assume the polynomials $`f_i(𝐱),g_i(𝐱)`$, $`i=0,\mathrm{},l,`$ have their coefficients in $`k`$. Let $`F=f_0g_0\mathrm{}f_lg_l`$ and denote by $`D`$ the divisor defined by $`F=0`$. Let $`(Y,h)`$ be a resolution of $`F`$. By this, we mean that $`Y`$ is a smooth and connected algebraic variety, $`h:YX=𝔸_k^m`$ is proper, that the restriction $`h:Yh^1(D)XD`$ is an isomorphism, and that $`(h^1(D))_{\mathrm{red}}`$ has only normal crossings as a subvariety of $`Y`$. Let $`E_i,`$ $`iT`$ be the irreducible (smooth) components of $`(h^1(D))_{\mathrm{red}}`$ where $`D`$ is the divisor defined by $`F=0`$ in $`X=𝔸^m.`$ For each $`iT,`$ denote by $`N_i(f_j)`$ and $`N_i(g_j)`$ the multiplicity of $`E_i`$ in the divisor of $`f_jh`$ and $`g_jh`$ on $`Y`$ and by $`\nu _i1`$ the multiplicity of $`E_i`$ in the divisor of $`h^{}dx`$, where $`dx`$ is a local non vanishing volume form. For $`iT`$ and $`IT,`$ we consider the schemes $`E_i^{}=E_i\backslash _{ji}E_j,E_I=_{iI}E_i`$ and $`E_I^{}=E_I\backslash _{jT\backslash I}E_j.`$ When $`I=\mathrm{},`$ we have $`E_{\mathrm{}}=Y.`$ Define a closed cone $`\overline{D_T}`$ as follows: $$\overline{D_T}:=\{(x_1,\mathrm{},x_t)_0^t:\underset{j=1}{\overset{t}{}}N_j(f_i)x_j\underset{j=1}{\overset{t}{}}N_j(g_i)x_j\text{ for }i=1,\mathrm{},l\},$$ where $`\mathrm{card}T=t`$ and $`_0=\{x:x0\}`$. Denote the lattice points in $`\overline{D_T}`$ by $`\overline{\mathrm{\Delta }_T},`$ i.e. $`\overline{\mathrm{\Delta }_T}=\overline{D_T}^t.`$ We can write this cone as a disjoint union of open simplicial pieces that we shall call $`R_k`$, $`k=0,1,\mathrm{},w.`$ We shall assume that $`R_0=(0,\mathrm{},0)`$ and that the next $`q`$ pieces are all the open one-dimensional edges in the cone $`\overline{D_T}`$: for $`k=1,\mathrm{},q`$, $$R_k=\{\alpha 𝐞_k=\alpha (q_{k1},\mathrm{},q_{kt}):\alpha >0\}.$$ Since these are all the edges, for any $`k\{0,\mathrm{},w\}`$ there exists some subset $`M_k\{1,\mathrm{},q\}`$ such that $$R_k=\{\underset{jM_k}{}\alpha _j𝐞_j:\alpha _j>0\text{ for all }jM_k\}.$$ Note that $`m_k:=\mathrm{card}M_kt.`$ We can choose the decomposition of the cone $`\overline{D_T}`$ so that for any $`k\{0,\mathrm{},w\}`$ $$R_k^t=\{\underset{jM_k}{}l_j𝐞_j:l_j_{>0}\text{ for all }jM_k\},$$ as follows from , p.123-124. Define for each $`k=1,\mathrm{},q`$ the following non-negative constants: (7.4.1) $`A_k`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{t}{}}}q_{kj}N_j(f_0)`$ $`B_k`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{t}{}}}q_{kj}\left(N_j(g_0)+\nu _j\right).`$ For each $`IT`$ define: $`D_I`$ $`=`$ $`\{(k_1,\mathrm{},k_t)\overline{D_T}:k_i>0\text{ if }iI\text{ and }k_i=0\text{ if }iT\backslash I\}`$ $`\mathrm{\Delta }_I`$ $`=`$ $`D_I^t.`$ So $`\overline{D_T}=_{IT}D_I`$, a disjoint union. For each $`IT`$ there is then a subset $`W_I\{0,\mathrm{},w\}`$ so that $$D_I=\underset{kW_I}{}R_k.$$ If $`kW_I`$, we set $`I_k=I`$. We can now state the explicit formula for motivic cone integrals: ###### 7.5 Theorem. Assume the polynomials $`f_i(𝐱),g_i(𝐱)`$, $`i=0,\mathrm{},l,`$ have their coefficients in $`k`$. The following equality holds in $`\overline{𝒩}`$: (7.5.1) $$Z_{𝒟,\mathrm{geom}}(s)=\underset{k=0}{\overset{w}{}}(𝐋1)^{\|I_k\|}𝐋^m[E_{I_k}^{}]\underset{jM_k}{}\frac{𝐋^{(A_js+B_j)}}{(1𝐋^{(A_js+B_j)})}.$$ ###### Proof. The proof works just the same as the one in for the explicit formula for $`p`$-adic cone integrals, using Lemma 3.4 of for performing the change of variable, cf. the proof of Proposition 2.2.2 of . ∎ ### 7.6. Our goal now is to explain the relationship between motivic cone integrals and $`p`$-adic cone integrals. We begin with some preliminaries. We assume $`k=`$ (in fact we could assume as well that $`k`$ is a number field or even only a field of finite type over $``$). For any variety $`X`$ over $``$, one can choose a model $`𝒳`$ of $`X`$ over $``$, and consider the number of points $`n_p(X)`$ of the reduction of $`𝒳`$ modulo $`p`$, for $`p`$ a prime number. Of course, for some prime numbers $`p`$, $`n_p(X)`$ may depend of the model $`𝒳`$, but, for a given $`X`$, the numbers $`n_p(X)`$ are well defined for almost all $`p`$. If we denote by $`𝒫`$ the set of all primes, the sequence $`n_p(X)`$ is well defined as an element of the ring $`^𝒫^{}`$, where, for any ring $`R`$, we set $`R^𝒫{}_{}{}^{}:=_{p𝒫}R/_{p𝒫}R`$. Moreover, counting points being additive for disjoint unions and multiplicative for products, the sequence $`n_p(X)`$ in $`^𝒫^{}`$ only depends on the class of $`X`$ in $``$ and may be extended to a ring morphism $`n:^𝒫^{}`$. Setting $`n_p(𝐋^1)=1/p`$, one may extend uniquely $`n`$ to a ring morphism $`n:_{\mathrm{loc}}^𝒫^{}`$. Furthermore, by Lemma 8.1.1 of , the morphism $`n`$ factors through the image $`\overline{}_{\mathrm{loc}}`$ of $`_{\mathrm{loc}}`$ in $`\widehat{}`$. Recall that we denoted by $`\overline{𝒩}`$ the subring of the ring $`\widehat{}[[T]]`$ which is generated by the image in $`\widehat{}[[T]]`$ of $`_{\mathrm{loc}}[T]`$, $`(𝐋^i1)^1`$ and $`(1𝐋^aT^b)^1`$ with $`i`$, $`a`$, $`b\{0\}`$. Hence, sending $`(𝐋^i1)^1`$ and $`(1𝐋^aT^b)^1`$ to $`(p^i1)^1`$ and $`(1p^aT^b)^1`$, respectively, one can uniquely extend $`n`$ to a ring morphism $`n:\overline{𝒩}[[T]]^𝒫^{}`$. ### 7.7. Assume $`k=`$ and that $`f_i`$ and $`g_i`$, for $`i=0,\mathrm{},l`$, are polynomials with coefficients in $``$. Then, for every prime $`p`$, we can consider the $`p`$-adic cone integral $$Z_{𝒟,p}(s)=_{V_p}p^{\mathrm{ord}_p(f_0)s\mathrm{ord}_p(g_0)}𝑑\mu _p,$$ with $$V_p=\{𝐱_p^m:\mathrm{ord}_p(f_i(𝐱))\mathrm{ord}_p(g_i(𝐱))\text{ for }i=1,\mathrm{},l\}.$$ Here $`\mathrm{ord}_p`$ stands for $`p`$-adic valuation and $`d\mu _p`$ denotes the additive Haar measure on $`_p^m`$, normalized so that $`_p^m`$ has volume 1. By , for almost all $`p`$, $$Z_{𝒟,p}(s)=\underset{k=0}{\overset{w}{}}(p1)^{\|I_k\|}p^m\|E_{I_k}^{}(𝔽_p)\|\underset{jM_k}{}\frac{p^{(A_js+B_j)}}{(1p^{(A_js+B_j)})}.$$ Hence Theorem 7.5 has the following corollary: ###### 7.8 Corollary. For almost all primes $`p`$, $$n_p(Z_{𝒟,\mathrm{geom}}(s))=Z_{𝒟,p}(s).$$ (In the left hand side $`T=𝐋^s`$, while in the right hand side $`T=p^s`$.) In particular, we have the following relationship between the motivic zeta function of a Lie algebra and the zeta functions counting finite index $`_p`$-subalgebras: ###### 7.9 Corollary. Let $`L`$ be a Lie-algebra of finite rank defined over $``$. For almost all $`p`$, $$n_p(P_{L[[t]],𝒳_t^{}}(𝐋^s))=\zeta _{L_p}(s).$$ ### 7.10. In fact, we should mention that a result such as Corollary 7.8 holds in much wider generality. Let $`k`$ be a finite extension of $`𝐐`$ with ring of integers $`𝒪`$ and $`R=𝒪[\frac{1}{N}]`$, with $`N`$ a non zero multiple of the discriminant. For $`x`$ a closed point of $`\mathrm{Spec}R`$, we denote by $`K_x`$ the completion of the localization of $`R`$ at $`x`$, by $`𝒪_x`$ its ring of integer, and by $`𝐅_x`$ the residue field at $`x`$, a finite field of cardinality $`q_x`$. Let $`f(x)`$ be a polynomial in $`k[x_1,\mathrm{},x_m]`$ (or more generally a definable function in the first order language of valued fields with variables and values taking place in the valued field and with coefficients in $`k`$) and let $`\phi `$ be a formula in the language of valued fields with coefficients in $`k`$, free variables $`x_1,\mathrm{},x_m`$ running over the valued field and no other free variables. Now set $`W_x:=\{y𝒪_x^m:\phi (y)\text{holds}\}`$ and define $$I_{\phi ,f,x}(s)=_{W_x}\|f\|_x^s\|dx\|_x,$$ for $`x`$ a closed point of $`\mathrm{Spec}𝒪`$. One should remark these integrals are of much more general nature than the previously considered cone integrals, since quantifiers are now allowed in $`\phi `$. In this more general setting we still want a geometrical object which specializes, for almost all $`x`$ to the local integrals $`I_{\phi ,f,x}(s)`$. In this aim is achieved, at the cost of replacing the “naive” ring $`_{\mathrm{loc}}`$ by a Grothendieck group of Chow motives<sup>2</sup><sup>2</sup>2For an indication of why Chow motives arise naturally in these questions, see 7.11. More precisely, it is shown in , that there exists a canonical rational function $`I_{\phi ,f,\mathrm{mot}}(T)`$ with coefficients in an appropriate Grothendieck ring of Chow motives, such that, for almost all closed points $`x`$ in $`\mathrm{Spec}𝒪`$, $`I_{\phi ,f,\mathrm{mot}}(T)`$ specializes - after taking the trace of the Frobenius on an étale realization and after setting $`T=q_x^s`$ \- to the $`p`$-adic integral $`I_{\phi ,f,x}(s)`$. ### 7.11. To what extent are the constructible sets $`E_{I_k}^{}`$ in Theorem 7.5 unique? The left hand side of (7.5.1) being canonical, it follows that the class of the right hand side is independent of the resolution. On the other hand it is unclear how one can deduce the motivic zeta function from the $`p`$-adic ones. The first problem is that distinct varieties over $``$ may have the same number of points in $`𝔽_p`$, which is in particular the case for isogenous elliptic curves over $``$. Much deeper is the fact, due to Faltings , that if, for almost all $`p`$, $`n_p(E)=n_p(E^{})`$, with $`E`$ and $`E^{}`$ elliptic curves over $``$, then $`E`$ and $`E^{}`$ are isogenous. Since isogenous elliptic curves define the same Chow motive, one sees Chow motives appearing in a natural way in the play. Let us recall (see for more details), that for $`k`$ a field, a Chow motive over $`k`$ is just a triple $`(S,p,n)`$ with $`S`$ proper and smooth over $`k`$, $`p`$ an idempotent correspondence with coefficients in $``$ on $`X`$, and $`n`$ in $``$. One endows Chow motives over $`k`$ with the structure of a pseudo-abelian category which we denote by $`\mathrm{Mot}_k`$ and we denote by $`K_0(\mathrm{Mot}_k)`$ its Grothendieck group which may also be defined as the abelian group associated to the monoid of isomorphism classes of objects in $`\mathrm{Mot}_k`$ with respect to the natural sum $``$. There is also a tensor product on $`\mathrm{Mot}_k`$ inducing a product on $`K_0(\mathrm{Mot}_k)`$ which provides $`K_0(\mathrm{Mot}_k)`$ with a natural ring structure. Assume now that $`k`$ is of characteristic zero. By a result of Gillet and Soulé and Guillén and Navarro Aznar there exists a unique morphism of rings $$\chi _c:K_0(\mathrm{Mot}_k)$$ such that $`\chi _c([S])=[h(S)]`$ for $`S`$ projective and smooth, where $`h(S)`$ denotes the Chow motive associated to $`S`$, i.e. $`h(S)=(S,\mathrm{id},0)`$. Let us still denote by $`𝐋`$ the image of $`𝐋`$ by $`\chi _c`$. Since $`𝐋=[(\mathrm{Spec}k,\mathrm{id},1)]`$, it is invertible in $`K_0(\mathrm{Mot}_k)`$, hence $`\chi _c`$ can be extended uniquely to a ring morphism $$\chi _c:_{\mathrm{loc}}K_0(\mathrm{Mot}_k).$$ ###### 7.12 Problem. Assume $`k=`$ (one can ask a similar question when $`k`$ is a number field, or more generaly of finite type over $``$). Is it true that for $`\alpha `$ in $`_{\mathrm{loc}}`$ if $`n(\alpha )=0`$ then $`\chi _c(\alpha )=0`$? A positive answer to this question would imply that the knowledge of the motivic integrals may be deduced from that of the corresponding $`p`$-adic integrals. ### 7.13. Since in the simplest examples, as the ones given in §9, the zeta functions $`\zeta _{L_p}(s)`$ are just rational functions of $`p`$ and $`p^s`$ and the corresponding motivic zeta functions rational functions of $`𝐋`$ and $`𝐋^s`$, we give an example showing that non trivial motives may indeed appear. Let $`L`$ be the class two nilpotent Lie algebra over $``$ of dimension 9 as a free $``$-module given by the following presentation: $$L=\begin{array}{c}x_1,x_2,x_3,x_4,x_5,x_6,y_1,y_2,y_3:(x_1,x_4)=y_3,(x_1,x_5)=y_1,(x_1,x_6)=y_2\\ (x_2,x_4)=y_2,(x_2,x_6)=y_1,(x_3,x_4)=y_1,(x_3,x_5)=y_3\end{array}$$ where all other commutators are defined to be 0. This Lie algebra was discovered by the first author to answer negatively the following question asked by Grunewald, Segal and Smith in : if $`L`$ is a Lie algebra over $``$, do there exist finitely many rational function $`\mathrm{\Phi }_1(X,Y),\mathrm{},\mathrm{\Phi }_r(X,Y)`$ such that for each prime $`p`$ there exists $`i\{1,\mathrm{},r\}`$ and $`\zeta _{L_p}^{}(s)=\mathrm{\Phi }_i(p,p^s)\mathrm{?}`$ Let $`E`$ be the elliptic curve $`Y^2=X^3X.`$ One can show, by a direct calculation (see ), that there exist two non zero rational functions $`P_1(X,Y)`$ and $`P_2(X,Y)(X,Y)`$ such that, for almost all primes $`p`$, $$\zeta _{L_p}^{}(s)=P_1(p,p^s)+\|E(𝔽_p)\|P_2(p,p^s).$$ By the result of Faltings already alluded to, it follows that the Chow motive of the curve $`E`$ is canonically attached to the Lie algebra $`L`$. It seems also quite natural to guess, but this has to be checked, that the calculation of the associated motivic cone integral gives $$P_1(𝐋,𝐋^s)+[E]P_2(𝐋,𝐋^s).$$ To see where the elliptic curve is hidden in this Lie algebra, consider $`det((x_i,x_{i+3}))=0`$ which gives an equation for the projective version of $`E.`$ ### 7.14. There is one case in which one can deduce an expression for the motivic zeta function from the corresponding expression for the $`p`$-adic zeta functions: ###### 7.15 Theorem. Let $`L`$ be a Lie algebra over $``$ additively isomorphic to $`^d.`$ Suppose that there exist two rational functions $`\mathrm{\Phi }(X,Y)`$ and $`\mathrm{\Psi }(X,Y)(X,Y)`$ such that, for almost all primes $`p,`$ $`\zeta _{L_p}(s)=\mathrm{\Phi }(p,p^s)`$ and $`P_{L[[t]],𝒳_t^{}}(𝐋^s)=\mathrm{\Psi }(𝐋,𝐋^s).`$ Then $`\mathrm{\Phi }(X,Y)=\mathrm{\Psi }(X,Y).`$ ###### Proof. This follows from Corollary 7.9 and the fact that if $`\mathrm{\Phi }(p,p^s)=\mathrm{\Psi }(p,p^s)`$ for almost all primes $`p`$ then $`\mathrm{\Phi }(X,Y)=\mathrm{\Psi }(X,Y).`$ ## 8. Topological zeta functions ### 8.1. In a topological zeta function was associated to each polynomial $`f[X]`$. The function was defined from the explicit formula in terms of a resolution of singularities for the associated local Igusa zeta functions of $`f`$. It can be interpreted as considering the explicit formula as $`p1.`$ The heart of the paper was the proof that this definition was independent of the resolution. The proof given there depended essentially on using the $`p`$-adic analysis of these local zeta functions and the associated explicit formulae. However the current motivic setting gives a proof of the independence of the resolution without going to the local zeta functions as we shall explain. In the setting of groups and Lie algebras it was unclear that such a topological zeta function could be associated to every nilpotent group or $``$-Lie algebra. This depended on producing an explicit formula. Since now provides such an explicit formula for zeta functions of nilpotent groups and Lie algebras and more generally for any cone integral, we can define the associated topological zeta function and use the setting of the current paper to prove that it is independent of the resolution and subdivision of the cone. We now proceed similarly as in section 2.3 of to define topological zeta functions. Let $`[𝐋^s]_{\mathrm{loc}}^{}`$ be the subring of $`[𝐋^s]_{\mathrm{loc}}`$ generated by the ring of polynomials $`_{\mathrm{loc}}[𝐋^s]`$ and by the quotients $`(𝐋1)(1𝐋^{Nsn})^1`$ for $`(N,n)^2\backslash (0,0)`$. Recall we denoted by $`\overline{}_{\mathrm{loc}}`$ the image of $`_{\mathrm{loc}}`$ in $`\widehat{}`$. In particular $`\overline{}_{\mathrm{loc}}`$ is a quotient without $`(𝐋1)`$-torsion of $`_{\mathrm{loc}}`$. By expanding $`𝐋^s`$ and $`(𝐋1)(1𝐋^{Nsn})^1`$ into series in $`𝐋1,`$ one gets a canonical morphism of algebras $$\phi :[𝐋^s]_{\mathrm{loc}}^{}\overline{}_{\mathrm{loc}}[s][(Ns+n)^1]_{(N,n)^2\backslash (0,0)}[[𝐋1]],$$ where $`[[𝐋1]]`$ denotes completion with respect to the ideal generated by $`𝐋1`$. Taking the quotient of $`\overline{}_{\mathrm{loc}}[s][(Ns+n)^1]_{(N,n)^2\backslash (0,0)}[[𝐋1]]`$ by the ideal generated by $`𝐋1,`$ one obtains the evaluation morphism: $$\begin{array}{c}\mathrm{ev}_{𝐋=1}:\overline{}_{\mathrm{loc}}[s][(Ns+n)^1]_{(N,n)^2\backslash (0,0)}[[𝐋1]]\hfill \\ \hfill \left(\overline{}_{\mathrm{loc}}/𝐋1\right)[s][(Ns+n)^1]_{(N,n)^2\backslash (0,0)}.\end{array}$$ We denote by $`\chi _{\mathrm{top}}(X)`$ the usual Euler characteristic of a smooth, projective $`k`$-scheme (say in étale $`\overline{_l}`$ -cohomology). By §6.1 of , this factorises through a morphism $`\chi _{\mathrm{top}}:\overline{}_{\mathrm{loc}}`$. Since $`\chi _{\mathrm{top}}(𝐋)=1`$, this induces then a morphism $$\chi _{\mathrm{top}}:\left(\overline{}_{\mathrm{loc}}/𝐋1\right)[s][(Ns+n)^1]_{(N,n)^2\backslash (0,0)}[s][(Ns+n)^1]_{(N,n)^2\backslash (0,0)}.$$ ###### 8.2 Definition. The topological cone zeta function $`Z_{𝒟,\mathrm{top}}(s)`$ associated to the cone data $`𝒟`$ is defined to be $$Z_{𝒟,\mathrm{top}}(s):=\left(\chi _{\mathrm{top}}\mathrm{ev}_{𝐋=1}\phi \right)\left(Z_{𝒟,\mathrm{geom}}(s)\right).$$ ###### 8.3 Definition. Let $`L`$ be a $`k[[t]]`$-algebra additively isomorphic to $`\left(k[[t]]\right)^d`$ or a $``$-Lie algebra additively isomorphic to $`^d.`$ Let $`L`$ be an algebra additively isomorphic to $`^d.`$ We define the topological zeta function $`Z_{L,\mathrm{top}}(s)`$ of $`L`$ to be $$Z_{L,\mathrm{top}}(s):=\left(\chi _{\mathrm{top}}\mathrm{ev}_{𝐋=1}\phi \right)\left(Z_{L,\mathrm{geom}}(s)\right).$$ These definitions make sense since $`Z_{𝒟,\mathrm{geom}}(s)`$ and $`Z_{L,\mathrm{geom}}(s)`$ belong to $`[𝐋^s]_{\mathrm{loc}}^{}`$ by Theorem 7.5. We can get an explicit formula for the topological zeta functions from that for the associated motivic zeta function. We use the notation defined prior to the statement of Theorem 7.5. We define the subset $`W_{\mathrm{top}}\{0,\mathrm{},w\}`$ such that $`kW_{\mathrm{top}}`$ if and only if $`\|I_k\|=\|M_k\|.`$ These are the indices of pieces of the cone whose dimension is equal to the dimension of the vector space spanned by all those basis elements with non-zero coefficients somewhere in the piece. ###### 8.4 Proposition. Let $`𝒟`$ be a set of cone integral data. Then with the notation above: $$Z_{𝒟,\mathrm{top}}(s)=\underset{kW_{\mathrm{top}}}{}\chi _{\mathrm{top}}(E_{I_k}^{})\underset{jM_k}{}\frac{1}{(A_js+B_j)}.$$ Some pieces $`R_j`$ of the cone may be missing in this expression. Note that in the special case of the Igusa zeta function we do see all the pieces $`R_j`$ since the cone there is the full positive quadrant. Nevertheless, in the more general setting of a cone integral, the constants $`A_j`$ and $`B_j`$ ($`j\{1,\mathrm{},q\}`$) will appear in this formula somewhere whenever the one-dimensional edge $`R_j`$ sits on an open simplicial piece of the cone of dimension corresponding to the number of non-zero entries of $`R_j`$. ## 9. Examples ### 9.1. Free abelian Let $`L=^d.`$ Then $`P_{Lk[[t]],𝒳_t^{}}(s)`$ $`=`$ $`(1𝐋^s)^1\mathrm{}(1𝐋^{s+d1})^1`$ $`Z_{L,\mathrm{geom}}(s)`$ $`=`$ $`{\displaystyle \frac{(1𝐋^1)}{(1𝐋^{sd})}}\mathrm{}{\displaystyle \frac{(1𝐋^d)}{(1𝐋^{s1})}}`$ $`Z_{L,\mathrm{top}}(s)`$ $`=`$ $`{\displaystyle \frac{1}{(s+d)\mathrm{}(s+1)}}.`$ The calculation of the integral in the $`p`$-adic setting translates immediately to a calculation of the motivic integral, proving the above identities. ### 9.2. Heisenberg Let $$L=\left(\begin{array}{ccc}0\hfill & \hfill & \hfill \\ & 0\hfill & \hfill \\ & & 0\hfill \end{array}\right).$$ Then $`P_{Lk[[t]],𝒳_t^{}}(s)`$ $`=`$ $`(1𝐋^s)^1(1𝐋^{1s})^1(1𝐋^{22s})^1(1𝐋^{32s})^1(1𝐋^{33s})`$ $`Z_{L,\mathrm{geom}}(s)`$ $`=`$ $`{\displaystyle \frac{(1𝐋^1)(1𝐋^2)(1𝐋^3)(1𝐋^{3s6})}{(1𝐋^{s3})(1𝐋^{s2})(1𝐋^{2s4})(1𝐋^{2s3})}}`$ $`Z_{L,\mathrm{top}}(s)`$ $`=`$ $`{\displaystyle \frac{3}{2\left(s+3\right)\left(s+2\right)\left(2s+3\right)}}.`$ Using Proposition 5.3, $$Z_{L,\mathrm{geom}}(s)=\frac{(1𝐋^1)(1𝐋^2)(1𝐋^3)}{(1𝐋^1)^3}_A𝐋^{\mathrm{ord}_t(x)(s+2)}𝐋^{\mathrm{ord}_t(y)(s+1)}𝐋^{\mathrm{ord}_t(z)s}𝑑\mu $$ where $$A=\{(x,y,z)k[[t]]^3:\mathrm{ord}_t(xy)\mathrm{ord}_t(z)\}.$$ The polynomial associated to the cone integral data already defines a non-singular space consisting only of normal crossings. So no resolution of singularities is required. The associated cone $`\overline{D_T}`$ is the span of the four one-dimensional pieces where we record the associated $`A_j`$ and $`B_j`$: $`R_1`$ $`=`$ $`(1,0,1),A_1=2,B_1=4`$ $`R_2`$ $`=`$ $`(1,0,0),A_2=1,B_2=3`$ $`R_3`$ $`=`$ $`(1,0,1),A_3=2,B_3=3`$ $`R_4`$ $`=`$ $`(0,1,0),A_4=1,B_4=2.`$ The following is a table of the pieces $`R_k`$ $`k\{0,\mathrm{},w\}`$ corresponding to a suitable decomposition of this cone where $`R_{i_1}\mathrm{}R_{i_r}`$ denotes the open piece spanned by $`R_{i_1},\mathrm{},R_{i_r}`$: $$\begin{array}{cccc}R_k\hfill & \|I_k\|\hfill & \|M_k\|\hfill & E_{I_k}^{}\hfill \\ 0\hfill & 0\hfill & 0\hfill & \left(𝐋1\right)^3\hfill \\ R_1\hfill & 2\hfill & 1\hfill & \left(𝐋1\right)\hfill \\ R_2\hfill & 1\hfill & 1\hfill & \left(𝐋1\right)^2\hfill \\ R_3\hfill & 2\hfill & 1\hfill & \left(𝐋1\right)\hfill \\ R_4\hfill & 1\hfill & 1\hfill & \left(𝐋1\right)^2\hfill \\ R_1R_2\hfill & 2\hfill & 2\hfill & \left(𝐋1\right)\hfill \\ R_1R_3\hfill & 3\hfill & 2\hfill & 1\hfill \\ R_2R_3\hfill & 3\hfill & 2\hfill & 1\hfill \\ R_2R_4\hfill & 2\hfill & 2\hfill & \left(𝐋1\right)\hfill \\ R_3R_4\hfill & 2\hfill & 2\hfill & \left(𝐋1\right)\hfill \\ R_1R_2R_3\hfill & 3\hfill & 3\hfill & 1\hfill \\ R_2R_3R_4\hfill & 3\hfill & 3\hfill & 1\hfill \end{array}$$ Using our explicit formula it is possible then to derive the expression for $`Z_{L,\mathrm{geom}}(s)`$ given above. To derive the expression for $`Z_{L,\mathrm{top}}(s)`$ note that $`\chi _{\mathrm{top}}(\left(𝐋1\right)^i)=0`$ for $`i>0.`$ Hence those pieces $`R_k`$ for which $`\|I_k\|=\|M_k\|`$ and $`\chi _{\mathrm{top}}(E_{I_k}^{})0`$ are $`R_1R_2R_3`$ and $`R_2R_3R_4.`$ This gives the following expression for $`Z_{L,\mathrm{top}}(s):`$ $`Z_{L,\mathrm{top}}(s)`$ $`=`$ $`{\displaystyle \frac{1}{(2s+4)(s+3)(2s+3)}}+{\displaystyle \frac{1}{(s+3)(2s+3)(s+2)}}`$ $`=`$ $`{\displaystyle \frac{3}{2\left(s+3\right)\left(2s+3\right)\left(s+2\right)}}.`$ Of course one could directly deduce this by applying the function $`\chi _{\mathrm{top}}\mathrm{ev}_{𝐋=1}\phi `$ to the final expression for $`Z_{L,\mathrm{geom}}(s)`$. However it is instructive to see the explicit formula for $`Z_{L,\mathrm{top}}(s)`$ at work. ### 9.3. $`𝔰l_2`$ Let $$L=\{\left(\begin{array}{cc}a\hfill & b\hfill \\ c\hfill & a\hfill \end{array}\right):(a,b,c)^3\}.$$ Then $`P_{Lk[[t]],𝒳_t^{}}(s)`$ $`=`$ $`(1𝐋^s)^1(1𝐋^{1s})^1(1𝐋^{22s})^1(1𝐋^{12s})^1(1𝐋^{13s})`$ $`Z_{L,\mathrm{geom}}(s)`$ $`=`$ $`{\displaystyle \frac{(1𝐋^1)(1𝐋^2)(1𝐋^3)(1𝐋^{3s8})}{(1𝐋^{s3})(1𝐋^{s2})(1𝐋^{2s4})(1𝐋^{2s5})}}`$ $`Z_{L,\mathrm{top}}(s)`$ $`=`$ $`{\displaystyle \frac{3s+8}{2\left(s+3\right)\left(s+2\right)^2\left(2s+5\right)}}.`$ Using Proposition 5.3 we have the following description of $`Z_{L,\mathrm{geom}}(s)`$ $$Z_{L,\mathrm{geom}}(s)=\frac{(1𝐋^1)(1𝐋^2)(1𝐋^3)}{(1𝐋^1)^3}_W𝐋^{\mathrm{ord}_t(a)(s+2)}𝐋^{\mathrm{ord}_t(x)(s+1)}𝐋^{\mathrm{ord}_t(z)s}𝑑\mu $$ where $$W=\{\left(\begin{array}{ccc}a\hfill & b\hfill & c\hfill \\ 0\hfill & x\hfill & y\hfill \\ 0\hfill & 0\hfill & z\hfill \end{array}\right)\mathrm{Tr}_3(R):\begin{array}{c}v(x)v(4zb)\\ v(x)v(4by)\\ v(zx)v(ax^2+4cxy4by^2)\end{array}\}.$$ The polynomial associated to the cone integral data is singular in this case. There have been several calculations of the zeta function of $`𝔰l_2(_p)`$ for $`p>2`$ (see and ) which give the answer above with $`𝐋`$ replaced by $`p`$. However, it is not possible to deduce directly that the above expression is a correct formula for the corresponding motivic zeta function, cf. the discussion at the end of Section 7. However, recently, the first author and Ph.D. student Gareth Taylor , have used the approach of applying resolution of singularities to the non-singular polynomial associated to the cone integral data of the integral above to deduce the expression for the $`p`$-adic integrals. The advantage of this approach is that it formally translates into a calculation of the corresponding motivic zeta function. Alternatively, one can observe that the resolution of the polynomial gives rise to $`E_I`$ which are generated out of $`𝐋`$, i.e. we don’t get any exotic varieties, so that there exists a rational function $`\mathrm{\Psi }(X,Y)(X,Y)`$ such that the motivic zeta function is $`\mathrm{\Psi }(𝐋,𝐋^s)`$. Hence, by Corollary 7.9, one is able to deduce that the rational function describing the $`p`$-adic expressions is the same as $`\mathrm{\Psi }(X,Y).`$
warning/0003/astro-ph0003360.html	ar5iv	text	# Galaxy Morphologies in the Cluster CL1358+62 at z=0.331footnote 11footnote 1Based on observations with the NASA/ESA Hubble Space Telescope and the W. M. Keck Observatory ## 1 Introduction The WFPC2 camera on the Hubble Space Telescope (HST) has made it possible to determine the morphologies of galaxies at intermediate redshift and beyond. It has been known for some time that the photometric and spectral properties of galaxies in intermediate redshift clusters differ from galaxies in nearby clusters; the population of blue, star-forming galaxies and post-starburst galaxies is larger at intermediate redshift, see e.g., Butcher & Oemler (1984), Dressler (1987), Gunn & Dressler (1988), Couch & Sharples (1987), and Dressler & Gunn (1999). Given the correlation between the spectral and morphological properties of galaxies (Morgan & Mayall (1957), Morgan & Osterbrock (1969), Kennicutt (1992)), we might be able to detect a corresponding evolution of galaxy morphology in intermediate redshift clusters. However, the evolution of morphology is likely to be more subtle than spectral evolution, since galaxies of the same morphological type can have significantly different star formation rates (e.g. Jansen et al. (1999)). In the most ambitious study of this sort with WFPC2 to date, the MORPHS group have classified over 1200 galaxies in 10 clusters at 0.37$`<`$z$`<`$0.56 (Smail et al. (1997)). They find that S0 galaxies are less common than in low redshift clusters and that the ratio of S0’s to E’s within a radius of $``$600 kpc (for H<sub>0</sub>=50, q<sub>0</sub>=0.5) decreases with redshift, falling from 2 in low redshift clusters to less than 0.5 at $`z=0.5`$. Andreon, Davoust & Heim (1997) (see also Andreon (1998)) have independently classified galaxies in a WFPC2 image of Cl0939+4713, one of the less concentrated clusters in the MORPHS sample. They find a ratio of S0’s to E’s of $``$2, quite comparable to a low-$`z`$ reference sample in the Coma Cluster. However, they classify 40-50% of the galaxies in Cl0939+4713 as spirals (S), in contrast with 20-30% S in a comparable region of the Coma Cluster. The sample of galaxies in Cl0939+4713 is relatively small ($``$70), and redshifts are available for less than one third of these. Couch et al. (1998) present a study of the morphological types in three clusters at $`z`$=0.31, also using WFPC2 images. There is substantial overlap between the authors of this paper and the MORPHS group, and the two groups have attempted to adopt a consistent morphological system. At $`z`$=0.31, Couch et al. find an excess of S’s, with an abundance at small radii ($``$400 kpc for $`H_0`$=50, $`q_0`$=0.5) approximately twice that in low-$`z`$ reference clusters. However, averaged over the three clusters, within 400 kpc, they find a ratio of S0’s to E’s at most slightly depressed relative to regions of comparable galaxy surface density in low-$`z`$ clusters. <sup>3</sup><sup>3</sup>3It is important to account for the morphology-density relation when we consider the morphological content of clusters. At high galaxy densities, typically found in the cores of clusters, the low-$`z`$ reference population becomes increasingly dominated by E’s, Dressler et al. (1997). Note, however, that the average $`z`$ of the MORPHS clusters is larger (0.46). Lubin et al. (1998) report the morphological types in two more distant clusters at $`z`$0.9. One cluster, CL0023+04, appears to be composed of two low velocity dispersion groups, and contains predominantly S’s. The other, CL1604+43, with a velocity dispersion of $``$1200 km s<sup>-1</sup>, contains $``$76% early-types. In the latter case, the S0/E ratio is found to be 1.7$`\pm `$0.9. This result is sensitive to the assumed morphological composition of the foreground/background population, but is evidence that the S0/E ratio does not decline smoothly with $`z`$. Our approach is complementary to the Dressler et al. (1997), Andreon, Davoust & Heim (1997), and Couch et al. (1998) studies which predominantly describe the galaxy morphologies in cluster cores. We use mosaics of WFPC2 fields to study the galaxy population in a larger region (allowing larger galaxy samples per cluster), and we have acquired large numbers of spectra of cluster galaxies. The spectra remove ambiguity about cluster membership and allow us to directly connect the morphological and spectral properties of the galaxies. Our sample of clusters is x-ray selected, with x-ray luminosities exceeding 4$`\times `$10<sup>44</sup> erg s<sup>-1</sup> in the 0.2–4.5 keV band. In CL1358+62 at $`z`$=0.33, we have drawn a complete sample of 518 galaxies to a magnitude limit $`I`$=22 from a WFPC2 mosaic image of CL1358+62 covering 53 square arcminutes. Spectra for 276 of the 518 galaxies in the morphological sample were previously obtained at the Multiple Mirror and William Herschel Telescopes. The color-magnitude relation of the 194 spectroscopically confirmed cluster members in the HST mosaic (3 are fainter than $`I`$=22) has been previously described in van Dokkum et al. (1998). The spectroscopic properties of 232 cluster members (some outside the HST mosaic), as well as the cluster dynamics have been described in Fisher et al. (1998). Our objectives in this paper are fourfold. (1) We introduce the morphological classification techniques that we will apply to our entire sample of clusters. (2) We classify the galaxies in CL1358+62 at z=0.33, comparing our classifications with those of an experienced external researcher, Alan Dressler. Our deep sample with two independent classifications provides a useful assessment of the scatter between WFPC2 visual morphological classifications at intermediate redshifts. (3) We describe the robust, classifier-independent conclusions and explore the physical implications of the differences between the two sets of classifications. (4) We connect the spectral and morphological types of the cluster galaxies. The paper is organized in the following fashion. In $`\mathrm{\S }`$ 2, we describe the photometric catalog from which the morphological sample was drawn. The two set of morphological classifications are discussed in $`\mathrm{\S }`$ 3. The morphological composition of the cluster and evidence for morphological evolution are presented in $`\mathrm{\S }`$ 4. The connection between the morphological and spectral properties of the galaxies is made in $`\mathrm{\S }`$ 5. $`\mathrm{\S }`$ 6 contains a brief discussion and conclusions. ## 2 Photometric Catalog Our photometric catalog, from which we draw galaxies for morphological classification, is derived from the HST F814W mosaic image. The techniques used to construct the mosaic are described in van Dokkum et al. (1998). Our use of the mosaic image, instead of individual WFPC2 CCD frames, slightly compromises the accuracy of the photometry, but considerably simplifies the source detection problem by eliminating most of the boundaries. Because our goal is to select a sample of galaxies for morphological classification to a consistent magnitude limit, rather than to perform precision photometry, this is a beneficial tradeoff. We used the SExtractor package, described by Bertin & Arnouts (1996), to detect the galaxies and perform the photometry. We use a detection and analysis threshold of 24.5 mag arcsec<sup>-2</sup>, and a zeropoint of 30.546 (3600 s exposures, May 1997 WFPC2 SYNPHOT update) to convert from instrumental magnitudes to a Cousins $`I`$ magnitude<sup>4</sup><sup>4</sup>4 Throughout the paper we consistently refer to a Cousins $`I`$ magnitude, which is very close to the natural F814W system referred to a Vega zeropoint (Holtzmann et al. (1995)). The adopted zeropoint is the average of the value for the 3 WFPC2 CCDs, which vary by $``$$`\pm `$0.01 magnitudes. We use the SExtractor total magnitude estimator (see Bertin & Arnouts (1996)), which is insensitive to the analysis threshold. In order to allow convenient comparisons with the results of the MORPHS group (Dressler et al. (1997)), we calculate a conversion of observed $`I`$ magnitudes (for $`z<`$0.6) to a $`V`$ magnitude in the rest frame. Some of the MORPHS classifications used WFPC2 F702W filter data, so we also derive a consistent conversion from F702W magnitudes to a rest frame $`V`$. The intent in Dressler et al. (1997) was to work to a consistent limit of $`M_V`$=-20, but a transcription error from Table 9 of Holtzmann et al. (1995) led to the adoption of a deeper limit of $`M_V`$$``$-19. (Here and throughout we use $`M_V`$ to refer to a rest frame $`V`$ absolute magnitude.) In contrast with Dressler et al. (1997), we also apply an evolutionary correction to allow a closer comparison with low $`z`$ clusters. We adopt the Dressler et al. (1997) cosmological model (H<sub>0</sub>=50 km s<sup>-1</sup> Mpc<sup>-1</sup>, q<sub>0</sub>=0.5) We use four numbers in addition to the distance modulus to convert from the WFPC2 magnitudes to $`M_V`$: (1) the conversion from the Vega-referenced WFPC2 natural filter system to Cousins filter bands, (2) the $`z`$-dependent $`K`$ correction to transform the observed Cousins $`R`$ and $`I`$ magnitudes to the rest frame, (3) the estimated rest frame galaxy colors, and (4) a $`z`$-dependent evolutionary correction. Each of these numbers depends on the spectral energy distribution of the galaxies, so the accuracy of this procedure is limited. We use these numbers, however, only to choose sample limiting magnitudes appropriately scaled with $`z`$. We adopt the approximate expressions described below for these conversions. We first fit a linear relation to the synthetic transformation of F702W to $`R`$ as a function of $`VR`$ (valid for $`VR`$$`<`$1.5), using results from Fig. 10 of Holtzmann et al. (1995): $`RF702W=0.31(VR)`$ (1a) We take the $`VR`$ colors of galaxies as a function of $`z`$ from Frei & Gunn (1994), with a morphological mix of 70% E and 30% Sbc to convert (1a) to: $`RF702W=0.31(0.53730.9132z+10.17z^211.21z^3)`$ (1b) To a good approximation (0.1 mag), Fig. 9 of Holtzmann et al. (1995) shows: $`IF814W=0`$ (1c) We fit polynomials to the average $`K`$ corrections of Frei & Gunn (1994) and Poggianti (1997), using 70% E and 30% Sbc or Sc contributions: $`K_R=0.4293z+3.807z^22.903z^3`$ (2a) $`K_I=0.4910z+0.4836z^2`$ (2b) From Frei & Gunn (1994) we take: $`V_{rest}R_{rest}=0.53`$ (3a) $`V_{rest}I_{rest}=1.10`$ (3b) From the the early-type galaxy fundamental plane study of Kelson et al. (1997) we derive an evolutionary correction: $`EC_{V_{rest}}=0.77z`$ (4) The corrections are applied in the following fashion: $`I=M_V+DM+EC_{V_{rest}}(V_{rest}I_{rest})+K_I`$ Here, $`DM`$ is the distance modulus. Applying these expressions to CL1358+62, with a distance modulus at $`z=0.3283`$ of 41.62, we find that the MORPHS limit of $`M_V`$=-20 corresponds to an observed $`I`$=20.5. Our photometric catalog completeness extends below $`I=23`$, but we have limited our morphological sample to $`I=22`$ to increase the reliability of the morphological classifications. At $`I=22`$, we attain a signal-to-noise ratio (S/N) that is very similar to the S/N attained at the $`I=23`$ classification limit adopted by the MORPHS group for their deeper images (Smail et al. (1997)). The standard deviation of the sky noise in our (3600 s) CL1358+62 image has an equivalent surface brightness of 24.4 mag arcsec<sup>-2</sup>, as compared with 25 mag arcsec<sup>-2</sup> typical for the ($``$12600 s) MORPHS images. A crude scaling relation can be derived by assuming that all galaxies have the same surface brightness, and that our criterion is that classifications of equal reliability require the same number of pixels at the same S/N. For a shallower image, we can effectively bin the image of a brighter galaxy that is $`n`$ times larger $`n\times n`$ pixels to achieve a S/N that is $`n`$ times better. To recover a sky S/N deficit of 0.6 mag, we need to set a limit $``$1.2 mag brighter. ## 3 Morphological Classification There are 518 galaxies in the morphological catalog after removing stars and image artifacts. In the rest frame, the filter central wavelength corresponds to $``$6100 Å. We use the F814W images for morphological classification because they are deeper than the F606W images. Dressler et al. (1997) have used either F702W and F814W images for classification at comparable redshifts (between 0.37 and 0.41). Below, we describe two sets of morphological classifications. The first set was carried out by the authors of the paper using the techniques described in $`\mathrm{\S }`$3.1. We realized that our classifications would be of greater value if we could compare them with morphological classifications from an external expert. Alan Dressler (AD) kindly agreed to independently classify the entire sample. Dressler is a member of the MORPHS group, and most importantly, has classified a large reference sample of low-$`z`$ cluster galaxies. The strength of the AD classifications is that the same experienced classifier has classified the galaxies at low and intermediate redshifts. Consistency has obvious benefits when searching for evolution. However, it is important to keep in mind the difficulty of classifying these relatively faint and small galaxies from WFPC2 images. The brighter cluster ellipticals and S0’s in CL1358+62 ($`I19`$) have effective radii (half-light radii), r<sub>e</sub>, of 0.5<sup>′′</sup> to 0.6<sup>′′</sup>. The faintest galaxies in the sample ($`I`$22) have r<sub>e</sub>$``$0.3<sup>′′</sup>. The 50% encircled energy radius of WFPC2 stellar images measured from our frames is about 0.17<sup>′′</sup>, so the numbers of meaningful information elements are small: typically 20 to 40. For this reason, we cannot be sure that agreement on a common morphological system is the most significant issue. Even the most experienced classifier may not account for all the systematic differences between the low $`z`$ (photographic) and the intermediate $`z`$ WFPC2 images. By comparing the two sets of morphological classifications we will be able to discern which aspects of the classification are most robust. ### 3.1 DF/MF/PvD Classifications The 518 sample galaxies were independently classified along the revised Hubble sequence by DF, MF and PvD from 96$`\times `$96 pixel “postage-stamps” drawn from both the mosaic and original individual CCD images. We referred frequently to Sandage (1961) and Sandage & Tammann (1987), as well as to artificially redshifted digital images drawn from the Nearby Field Galaxy Survey (Jansen et al. (1999)). We assigned numerical types as follows: -5 (elliptical), -4 (elliptical or S0), -2 (S0), 0 (S0 or Sa), 1 (Sa), 3 (Sb), 5 (Sc), 7 (Sd), 9 (Sm), 10 (Im) and 99 (peculiar or merger). Intermediate types (2, 4, 6, and 8) were also assigned. For 80% of the galaxies, the three independent classifications span a range of three or fewer numerical types (we consider -4, -2 and 0 to be adjacent numerical types), and agree exactly for 34% of the sample. For 6% of the sample, the object is considered unclassifiable by one or more of the authors, or else the classifications disagree wildly. In these cases, we assign a numerical type of 999. The remaining 14% are classified type 1 or later by all the authors, but with a broader range of classifications. Table 1 lists the combination rules used to assign types where the agreement is not exact, including type 15 for indeterminate late types. ### 3.2 AD Classifications Following the completion of the DF/MF/PvD classifications, AD independently classified the CL1358+62 galaxies according to techniques described in Smail et al. (1997). AD does not assign galaxies to the DF/MF/PvD intermediate types -4 (E or S0) or 0 (S0 or Sa) in the same fashion, preferring to subdivide these into E/S0 (-4), S0/E (-3), S0/Sa (-1) and Sa/S0 (0). When discussing the cluster population in broad terms, AD’s types -5 and -4 are combined into E, types -3, -2 and -1 into S0, and etc. DF/MF/PvD split the contents of the -4 bin equally into E and S0, and the 0 bin equally into S0 and Sa. When converting AD’s descriptive types into numerical types, we have placed a few galaxies into the merger (99) bin where an individual catalog object corresponds to an interacting or merging pair of objects. Augustus Oemler, another member of the MORPHS group, independently classified 307 galaxies from the CL1358+62 sample ($`I`$=22 limit) to check if AD’s classifications adhere to the MORPHS system. Oemler’s classifications agree very well with AD’s overall, with $``$10% more S0’s by number, and a corresponding decrease in the numbers of E’s. The number of spirals is identical in both classifications. ### 3.3 Quantitative Comparison of the Classifications Like any measurement, classifications will suffer from random and systematic errors. We can estimate these errors by comparing multiple sets of classifications. Ideally, such a comparison should be based on classifications using independent imaging data. Here, both sets of classifiers (DF/MF/PvD and AD) worked from the same data set, so we may underestimate the errors. Nonetheless, this comparison is quite interesting. Figure 1 shows the difference between DF/MF/PvD and AD morphological types for individual galaxies as a function of magnitude. As expected, the differences increase at fainter magnitudes. To avoid creating artifical gaps in Figure 1, we condensed and shifted the numerical types of the early type galaxies for this presentation: -3 for E’s, -2 for E/S0’s and S0/E’s, -1 for S0’s, 0 for S0/Sa and Sa/S0. For types Sa and later, the “normal” types were used. Galaxies which were classified as merger, peculiar, or unclassified by one group were ignored. Figure 2 is a scatter diagram comparing the two sets of classifications. The scatter is dominated by random differences, but there is some evidence for systematic differences for the early type galaxies. We return to this point below. The scatter between the classifications has been measured by calculating the mean absolute deviation, and normalizing it to the mean absolute deviation of a gaussian with an rms of 1. This measurement of the scatter is much more robust than the RMS of the differences. The result is shown in Figure 3. The scatter is $``$1 for bright galaxies, but increases strongly faintwards of $`I=21`$. The appendix describes how classification errors might systematically bias our population estimates; unfortunately we cannot simply calculate correction factors. In what follows, we restrict most of our analysis to the subset of galaxies with $`I<21`$. ### 3.4 Classification of Ellipticals and S0s Figure 2 shows a systematic difference in the DF/MF/PvD and AD classifications of early type galaxies. This is not surprising as the division between E and SO is a difficult problem in visual classifications. We might hope that looking at two objectively determined structural parameters for the galaxies, ellipticity and the bulge/total light ratio, might be helpful in resolving this issue. For example, Smail et al. (1997) compare the ellipticity distributions of the MORPHS E and S0 intermediate-$`z`$ galaxies with those from an Andreon et al. (1996) study of Coma Cluster galaxies as a consistency check of the MORPHS classifications. We have therefore examined the structural properties of the CL1358+62 galaxies measured by the Medium Deep Survey (MDS) group<sup>5</sup><sup>5</sup>5The MDS catalog is based on observations with the NASA/ESA HST, obtained at the STSCI, operated by AURA., Ratnatunga, Griffiths & Ostrander (1999). The MDS group has published structural properties for 70% of the CL1358+62 galaxies in our morphological catalog. The missing 30% are located near frame edges or in one mosaic frame that is not yet included in the MDS database. The MDS group has chosen the best fitting of four structural models for the objects in their catalogs: star, disk, bulge, or disk + bulge, accounting for the point spread function of HST. If a galaxy model is chosen, the MDS group fit the ellipticity of the disk, bulge or disk and bulge separately as appropriate. We focus on two structural parameters derived from these fits: the ratio of bulge to total light, and the weighted ellipticity. The ratio of bulge to total light follows trivially from the best fit model, or from the fits to the disk + bulge model. We calculate the weighted ellipticity from the MDS disk and bulge ellipticities, weighting by the fractions of light in the disk and bulge. In Figures 4 and 5 we plot the distributions of the bulge to total light ratio for the E and S0 galaxies with $`I<21`$ in the DF/MF/PvD and AD samples, respectively. Excepting the differences in numbers, the distributions for both E’s and S0’s look remarkably similar for the two sets of classifications. In Figures 6 and 7 we plot the distribution of weighted ellipticities for the same E and S0 galaxies ($`I<21`$), again showing the DF/MF/PvD and AD samples independently. Caution must be exercised when comparing to these to other measures of ellipticity since the MDS numbers are corrected for the HST PSF. A larger fraction of the AD E’s have ellipticities exceeding 0.2, but the mean AD E ellipticity is 0.22, only slightly larger than the DF/MF/PvD mean of 0.17. The ellipticities of the two samples of S0’s are similar: AD finds a mean ellipticity of 0.46 and DF/MF/PvD find a mean ellipticity of 0.40. The sizable difference in the numbers of E and S0 galaxies found by the two sets of classifiers is not reflected in a large difference between the structural parameters for the two samples of E’s or S0’s. We conclude that although the total number of early-type galaxies is well established, the SO to E ratio in CL1358+62 is uncertain. ## 4 Morphological Composition of CL1358+62 and Evidence for Evolution We have learned that the differences between the two sets of classifications rise steeply below $`I=21`$, suggesting this as a practical limit for our morphological study. This is also conveniently close to the effective limit of our spectroscopic completeness. We have redshifts for 277<sup>6</sup><sup>6</sup>6In one case, a spectrum was obtained of a pair of galaxies separated by 0.7<sup>′′</sup>, one (#1295) at $`I=21.17`$ and the other (#1297) at $`I=20.92`$. We have assigned the measured velocity (cz=99657) to both galaxies. Galaxies #1481 and #1483 may have both fallen within the spectrograph slit. We assign the measured velocity (cz=99792) to #1483, which is 1.2 mag brighter. of the 518 galaxies in the morphological sample, and 191 of these are cluster members by the criteria given in Fisher et al. (1998): 0.31461 $`<`$ z $`<`$ 0.34201. The spectroscopic completeness is 89% for the galaxies brighter than $`I=20.5`$ (corresponding to $`M_V=20)`$, falling to 59% for $`20.5<I<21`$, 33% for $`21<I<21.5`$, and 9% for galaxies with $`21.5<I<22`$. Table 2 lists positions, $`I`$ magnitudes, both sets of morphological classifications, and radial velocities for galaxies brighter than $`I`$=21. ### 4.1 Morphological Composition We may study the morphological composition of the subsample of known members, which is nearly complete to the MORPHS’s depth of $`M_V=20`$, without concern about background galaxy subtraction. The subsample of known cluster members to $`I=20.5`$ (or $`M_V=20`$, see $`\mathrm{\S }`$ 2) contains 138 galaxies. DF/MF/PvD classify 27$`\pm `$4% of these as E, 44$`\pm `$6% as S0, 29$`\pm `$5% as S, with 1 unclassified galaxy. (The errors here and in the following discussions account only for the Poisson statistics of the number of galaxies per classification bin.) AD classifies 35$`\pm `$5% as E, 38$`\pm `$5% as S0, and 27$`\pm `$4% as S. In both cases, the total early type population is $``$72%. To $`M_V`$=-20, then, the only difference between the two groups of classifiers is the relative numbers of E’s and S0’s. We consider also the complete photometric sample to a depth of $`I=21`$ (or $``$$`M_V=19.5`$), where a small background correction is necessary. To $`I=21`$ our sample contains 298 galaxies, for which we have 236 redshifts (79.2% completeness). Of the 236 galaxies with redshifts, 65$`\pm `$8 are nonmembers, yielding a foreground/background count of 82$`\pm `$10 after correcting for the spectroscopic completeness. We determine the morphological composition of the foreground/background galaxies directly from our spectroscopic sample, which contains 86 foreground/background galaxies. DF/MF/PvD classify 3% E, 13% S0, 77% S, and 6% mergers. AD classifies 11% E, 9% S0, 71% S, and 9% mergers. We average these classifications, and adopt a background composition of 7% E, 11% S0, 74% S and 8% mergers. For comparison, Dressler et al. (1997) adopt a morphological composition of 10% E, 10% S0, and 80% S. After correcting for background, to $`I=21`$, DF/MF/PvD find a population of 25$`\pm `$4% E’s, 46$`\pm `$5% S0’s, 29$`\pm `$6% S’s, and 0$`\pm `$1% mergers. AD finds 36$`\pm `$4% E’s, 36$`\pm `$5% S0’s, 28$`\pm `$6% S’s, and 1$`\pm `$1% mergers. These results are indistinguishable from those for the brighter spectrosopic sample, with $``$72% early-types in both sets of classifications. ### 4.2 Morphological Evolution We search for morphological evolution in CL1358+62 by comparing its population with that of equivalent low-$`z`$ clusters. Judging which low-$`z`$ clusters are equivalent is somewhat uncertain, but we take as an approximation low-$`z`$ clusters with a similar number of galaxies within a fixed metric aperture, allowing us to correct for the effects of the morphology-density relation. We use the nearby cluster catalog of Dressler (1980), reanalyzed and summarized in Dressler et al. (1997) as a benchmark. We use the subset of high concentration clusters (10 of 55) in the low-$`z`$ sample for comparison, because CL 1358+62 was selected for its high x-ray luminosity and has a concentration index C$``$0.49 (Fabricant, McClintock & Bautz (1991)). The data for the 10 high-concentration clusters are plotted in Fig. 12 of Dressler et al. (1997). There are an average of $``$63 cluster galaxies within a radius of 1450 kpc in these clusters to $`M_V`$=-20.4. In CL1358+62, there are $``$114 cluster members within this radius to $`M_V`$=-20.4, so CL1358+62 is richer than the average cluster in the low-$`z`$ sample. However, the low-$`z`$ sample does contain a high-concentration cluster as rich as CL1358+62: the Coma Cluster. Furthermore, the density difference between CL1358+62 and the average low-$`z`$ reference cluster, 0.3 dex, is comparable to the bin size in the morphology-density relation plots in Dressler et al. (1997). The comparison between the morphological composition of the CL1358+62 sample to $`M_V`$=-20 and the low-$`z`$ reference sample is made in Table 3 and Figure 8. The DF/MF/PVD S0/E ratio, 1.6$`\pm `$0.3, differs at only 1.4$`\sigma `$ confidence from the low-$`z`$ reference sample ratio of 2.1$`\pm `$0.2. AD’s ratio, 1.1$`\pm `$0.2, differs from the the low-$`z`$ reference sample ratio at 3.5$`\sigma `$ confidence. Despite this difference in the S0/E ratio, we stress that both classifications for CL1358+62 yield fractions of early types (E+S0), $``$72%, and late types (S), $``$28%, that are identical within the errors to the low-$`z`$ reference sample. We can therefore draw a robust conclusion from the two sets of morphological classifications: CL1358+62 does not contain an elevated population of spiral galaxies compared with low $`z`$ reference clusters. This contrasts with the results of Andreon, Davoust & Heim (1997) and and Couch et al. (1998) for other intermediate $`z`$ clusters. Our work suggests that the early type/spiral classifications are likely to be secure, implying that the populations of intermediate $`z`$ clusters vary significantly, even after accounting for the effects of the low $`z`$ morphology-density relation. This conclusion about the spiral population in CL1358+62 limits the range of physical models for evolution in CL1358+62, even allowing for our uncertainty about the E/S0 classifications. As we mentioned earlier, the AD classifications have the strong advantage of the same classifier at low and intermediate $`z`$. However, given the different character of the low-$`z`$ photographic images and the intermediate-$`z`$ WFPC2 images, we must acknowledge the possibility of systematic, redshift-dependent classification uncertainties. Figure 1 provides some reason to be cautious about this issue. Because we do not understand in detail the reasons for the differences between the two sets of visual classifications, we cannot be positive that the two sets of classifications bound our uncertainties. However, the best we can do at present is to leave the issue of E/S0 classifications open, and to explore the consequences of both sets of classifications below. The DF/MF/PvD classifications would imply that cluster evolution from $`z=0.33`$ to the present does not affect the cluster morphological composition or its morphology-density relation within the observed 1.4 Mpc radius aperture. Figure 9 shows this relation for CL1358+68, binning the background subtracted data to I=22 ($`M_V`$=-18.5) radially about the dominant central galaxy. Since we are looking only for radial trends, using the deeper sample is appropriate here. The average galaxy density in each of four radial bins (0–1, 1–2, 2–3, and 3–5 arcmin) is the abscissa for this histogram. The galaxy densities in Figure 9 have been normalized to a $`M_V`$=-20 limit to allow comparison with the low-$`z`$ reference sample. We find that the morphology-density relation for CL1358+62 is indistinguishable within the errors from the low-$`z`$ relation (to $``$$`M_V`$=-20) shown in Figure 3 of Dressler et al. (1997). The AD classifications suggest an evolutionary mechanism that decreases the fraction of E’s, increases the fraction of S0’s, while leaving the fraction of spirals unchanged. Assuming that there is no plausible mechanism for directly converting E’s into S0’s, we can exclude models that transform the observed $`z`$=0.33 cluster spirals into S0’s without accretion of additional galaxies. In order to transform the $`z`$=0.33 AD morphological mix into the low-$`z`$ reference sample population while accreting the smallest number of galaxies, the cluster population would increase by 50%. Approximately 70% of the accreted galaxies would become S0’s by the present day, and 30% spirals. ## 5 The Morphological-Spectral Connection While a great deal has been learned about the galaxy population in intermediate redshift clusters from relatively small samples of spectra, we must remember that even present-day clusters of galaxies are a heterogeneous group, differing widely in their degree of virialization. Because cluster relaxation may drive galaxy evolution, the range in cluster galaxy populations may be large at any redshift. For this reason, it is desirable to connect the morphological and spectral properties of a large sample of galaxies in each of a number of clusters directly. We have spectral classifications from Fisher et al. (1998) sorting each of the galaxies with spectra into one of the four categories: (1) absorption lines only, (2) emission lines present, (3) emission lines plus strong Balmer absorption lines, and (4) k+a (also called E+A). Category (4) contains galaxies with the normal absorption lines of E/S0 galaxies plus strong Balmer absorption lines. Emission line galaxies have \[OII\] 3727 Å emission with equivalent width (EW) $`>`$5 Å. If H<sub>δ</sub> absorption of $`>`$4 Å EW is detected for an emission line galaxy, the galaxy is classified as emission plus Balmer lines. Galaxies with \[(H<sub>δ</sub> \+ H<sub>γ</sub> \+ H<sub>β</sub>)/3\] EW greater than 4 Å, but \[OII\] emission with $`<`$5 Å EW, are classified as k+a. We refer the reader to Fisher et al. (1998) for a more complete summary of the spectral properties of these galaxies and a comparison with spectra of low-$`z`$ galaxies. Table 4 and Figure 10 summarize the comparison between the spectral and morphological properties for the 191 cluster members in common with Fisher et al. (1998), using the DF/MF/PvD classifications. In rough terms, the bulk of the CL1358+62 galaxies follow the morphological-spectral correlation expected for bright field galaxies at low redshift: the preponderance of the E and S0 galaxies have pure absorption line spectra, while the fraction of galaxies with emission lines rises for the late type spirals. Dressler et al. (1999) found a similar behavior for galaxies in the MORPHS sample; see also Poggianti (1999). We do know, however, that the percentage of galaxies with emission lines and strong Balmer absorption lines, $``$19%, is higher than the $``$6% content of these galaxies in comparable low $`z`$ clusters (Dressler (1987), see also Fisher et al. (1998)). If we use the AD classifications for the CL1358+62 galaxies, these conclusions do not change significantly. The most interesting difference between the two sets of morphological classifications is that DF/MF/PvD classify all the E+A galaxies as types S0 to Sb, while AD classifies these galaxies as having a wider range of morphologies from E to Sbc. ## 6 Discussion and Conclusions For CL1358+62, we have acquired a unique data set including a large mosaic of HST fields and extensive spectroscopy that allows us to unambiguously determine cluster membership for galaxies with $`M_V<20`$. We have directly compared the morphological classifications of two sets of classifiers for the galaxies in CL1358+62. The two sets of classifiers agree that (to a limit of $`M_V`$=-20) the fraction of early type galaxies (and therefore spirals) in this cluster at $`z`$=0.33 is indistinguishable from the fraction in comparable low-$`z`$ clusters. In contrast, previous workers, Andreon, Davoust & Heim (1997) and Couch et al. (1998), who also studied WFPC2 images of clusters at z$``$0.3, found an elevated population of spirals compared with low-$`z`$ reference samples. Because our work confirms the reliability of early-type/spiral classifications from intermediate $`z`$ WFPC2 observations, we conclude that this is evidence for a dispersion in the evolution of intermediate-$`z`$ clusters. The two groups of classifiers differ on ratio of E to S0 galaxies in CL1358+62. DF/MF/PvD find a population of S0 galaxies (S0/E=1.6$`\pm `$0.3) that is within 1.4$`\sigma `$ of the low-$`z`$ reference sample, while AD finds a significantly smaller ratio (1.1$`\pm `$0.2). This systematic difference is most likely related to the fact that the transition between S0’s and intermediate luminosity E’s is rather gradual. Many of the intermediate luminosity E’s are thought to have disks, e.g. Scorza et al. (1998), Rix & White (1990), and Jorgensen & Franx (1994). It may only be possible to resolve this issue by direct model fitting to images at low and intermediate $`z`$. Even though we conclude that we have not reliably determined the ratio of S0’s to E’s among the early types, our work significantly restricts possible evolutionary models. If we accept the MF/DF/PvD classifications, evolution must preserve the fraction of E’s, S0’s and S’s as well as the morphology-density relation. If the AD classifications are correct, evolution must decrease the fraction of E’s and increase the fraction of S0’s while maintaining the fraction of S’s. A possible mechanism for driving the evolution of the morphological mix in this latter fashion is accretion of additional galaxies from the spiral-rich infall region that become predominantly S0’s. The cluster population within an $``$1 Mpc radius must increase by a minimum of 50% from $`z`$=0.5 to the present day in order convert an intermediate $`z`$ population rich in E’s to a low-$`z`$ population rich in S0’s. It will be interesting to see whether such accretion can be produced in simulations of cluster formation. In most cluster formation scenarios, massive clusters form by the merging of pre-existing massive clusters with (presumably) similar populations of early type galaxies. It may therefore be difficult to double the ratio of S0 to E galaxies. We compare our morphological classifications to our previous spectral classifications and conclude that the morphologies of the spectrally “active” galaxies are as might be expected from the low-$`z`$ field population: the galaxies with emission lines are predominently spirals and the k+a (or E+A) post-starburst galaxies are typically early type disk galaxies (S0–Sb). We wish to acknowledge the generous contributions of Alan Dressler to this paper, including his independent classifications and insightful comments. We thank Gus Oemler for checking the classifications and for thoughtful comments on the manuscript, which led to several improvements. We thank Margaret Geller for a critical reading of an earlier version of the manuscript and helpful comments. Our referee, Ian Smail, and our editor, Greg Bothun, both made insightful comments that helped us clarify the paper. APPENDIX SYSTEMATIC EFFECTS OF MEASUREMENT ERRORS The effects of classification errors on the distribution of types can be important. Quantifying these effects is difficult because morphological classification is a subjective procedure, but we can gain some insight by considering simple models for the errors. We begin by assuming that the numerical type is based on a one dimensional measurement with a simple, constant error. This simple model would imply that any peaks in the distribution of types would be softened. If we adopt nominal intrinsic fractions of S:S0:E of 0.20:0.53:0.27, as is approximately correct for the inner 600 kpc of low redshift clusters (Dressler et al. (1997)), the errors will automatically decrease the fraction of S0’s, and enhance the fraction of E’s and S’s. We can calculate an upper limit to the loss of S0’s for our nominal 20:53:27 (S:S0:E) population with this model for the errors. We assume that the types are scattered with a mean absolute deviation (MAD) of 1, and all intermediate types are divided equally between adjacent types. The outcome of this experiment is a distribution of 0.27:0.38:0.34 (S:S0:E). In this case, all galaxies scattered beyond the normal type boundaries were assigned to the boundary type (i.e. E). A MAD of 1 may be an overestimate as the difference in the types assigned by the two groups of classifiers is of this order. The intrinsic errors are $`\sqrt{2}`$ smaller if the errors are independent. For these smaller errors the resulting distribution would be 0.24:0.43:0.33 (S:S0:E). In either case, the results serve to illustrate that the systematic effects can be very significant. Let us consider a second, more physical model for the visual classification errors. Here, we assume that morphological type is based on two independent variables with continuous distributions: bulge-to-total light fraction ($`f_b`$), and asymmetric features due to spiral arms ($`A`$). This is very similar to the quantitative classification devised by Abraham et al. (1996). Galaxies with low $`A`$ will be classified as early type ($`t<0`$) and then divided into E’s or S0’s based on whether $`f_b`$ is above or below a critical value. It has been argued that most $`L`$ ellipticals have faint disks, e.g. Rix & White (1990) and Jorgensen & Franx (1994). Similarly, the spiral classification will be based on a combination of $`f_b`$ and $`A`$. If the intrinsic distribution of $`f_b`$ is flat, then errors in $`f_b`$ will not change the ratio of S0’s to E’s. However, if the intrinsic distribution of $`f_b`$ is peaked, the errors will have a systematic effect. The sign depends on the details of the intrinsic $`f_b`$ distributions and the size of the errors. Contributing to this uncertainty, the errors can also be asymmetric, if for example, faint extended disks are missed in noisy data. The situation is similar for errors in $`A`$, where E’s and Sa’s now share a boundary. Extensive simulations are required to estimate the systematic effects of limited S/N on classification errors, taking into account the the point spread function of WFPC. We will undertake this effort in a future paper.
warning/0003/gr-qc0003003.html	ar5iv	text	# Close limit of grazing black hole collisions: non-spinning holes ## I Introduction The phrase “collision of black holes” has an aura of a mysterious and exotic happening that is not far from the reality of such an event. A black hole is not an ordinary object defined by the amount and properties of the material of which it is made. Rather it is a region from which no signal can escape. The surface, the black hole horizon, bounding this region is defined by the formal “no escape” property. Unlike the surface of an ordinary object the horizon has no local properties that would be sensed by an observer with the bad fortune to fall inward through it. A collision of two holes is the process in which two no-escape regions merge to become a single, larger, region of no escape. In the last few years such mergers have become the focus of much research attention, for two not entirely independent reasons. The first reason is the development of numerical relativity . General relativity, Einstein’s theory of gravity, sets the dynamics of space-time via a set of nonlinear partial differential equations of such complexity that analytic solutions have been limited to two classes: solutions of high symmetry, or solutions based on approximation techniques, such as linearized weak field theory. The study of Einstein’s equations on computers has been viewed as the key to finding more general asymmetric strong field solutions and it was natural for this key to be be applied to black hole collisions. Black holes are incontrovertibly strong field regions, but single isolated black holes are stationary solutions of Einstein’s theory, and the simplifying symmetry of time independence allows for closed form well-understood solutions . Collisions of black holes, on the other hand, are necessarily nonstationary as well as being crucially strong-field events. It is known that the collision will result in a single final black hole and in the generation of gravitational waves carrying off some of the mass energy originally associated with the holes. But this is all that is known with certainty. The nature of the merging of the horizons, in the general collision, is not even qualitatively understood. A reasonably complete understanding awaits progress in numerical relativity, and the wait has been longer than anticipated. The solution of general black hole collisions on computers has proved to be remarkably difficult. There is, however, a class of cases in which reliable answers are available. If the collision is a ’head-on’ collision along a straight line, then there is rotational symmetry about the line of the collision. Though the collision is still highly dynamic and nonlinear, the simplifications afforded by this symmetry reduce the computational demands sufficiently that the collision could successfully be simulated even in the mid 1970’s, and run with good reliability in the mid 1990’s. The simplification of head-on collisions, however, masks some of the physics of the most interesting types of collisions, the fully three dimensional collisions at the end point of the inspiral of a mutually orbiting pair of black holes. The second development that directed attention to black hole collisions is the advent of sensitive gravitational wave detectors. In the next few years, several interferometric gravitational wave observatories (the LIGO project in the US, the VIRGO and GEO projects in Europe and the TAMA project in Japan ) may be capable of detecting gravitational waves. Whether near term searches are successful will depend more than anything else on the strength of astrophysical sources. Attributes of a good generator of gravitational waves include strong gravitational fields and high velocities, so black hole processes are a natural source to consider. It is astrophysically plausible that black holes form binary associations with other objects, including other black holes . Due to the loss of energy by the emission of gravitational radiation, the separation and period of the binary orbits would decrease. If the binary consists of two black holes, the inspiral would end with a rapid strong field merger that has the potential to be a powerful source of detectable gravitational waves . The whole process of inspiral generates gravitational radiation, but in the early large-separation stages the radiation is relatively weak and is reasonably well described by Newtonian gravity theory and Post-Newtonian extensions of it . It is only the final strong field merger that could in principle produce a powerful burst of gravitational waves, but at this point only one parameter of the burst is reliably known. The characteristic frequency of the waves is inversely proportional to the mass of the final black hole formed, and works out to be on the order of $`10^3`$ Hz for a 10$`M_{}`$ hole, a typical expected mass of a “stellar” sized hole. For supermassive holes of mass $`10^6M_{}`$ typical of galactic nuclei, the waves would be less than 1 Hz. The maximum sensitivity of the next generation of gravitational wave detectors occurs at frequency around 100 Hz and the detectors will be ideally suited to waves from a black hole with mass of several hundred $`M_{}`$. Some recent observations offer indirect evidence that black holes in this range may exist. If they do not, then the detection of the collision of black holes may require the deployment of space-based detectors sensitive to the low frequency waves produced by supermassive holes. The ratio of the masses in a binary determines both how difficult it is to analyze, and how exciting it is as a potential source. If the mass of a black hole $`M_1`$ is much larger than the mass of its binary companion $`M_2`$, then the smaller mass object can be treated as a perturbation to the well understood spacetime of the larger mass black hole. The equations that describe perturbations are linear, and hence relatively easily dealt with in general. In the specific case of perturbations to black hole spacetimes, the techniques of calculation were worked out in the 1970s and resulted in the Regge-Wheeler and Zerilli equations for perturbations of Schwarzschild (nonrotating) black holes, and in the Teukolsky equation for perturbations of Kerr (rotating) black holes. The relatively easily analyzed “particle limit” case $`M_2M_1`$ may be of interest in connection, say, with neutron stars merging with supermassive black holes, but this process cannot give the hoped for high power. It is easy to show the gravitational wave power generated scales in the masses as $`(M_2/M_1)^2`$. High power requires roughly equal masses, and this means the simplifications of the particle limit do not apply to the most interesting sources. If not directly applicable to equal mass inspiral, the clarity of the particle limit can, at least, help us to formulate questions about the nature of the endpoint of inspiral, like the existence of a last stable circular orbit. As a particle orbits a black hole it reaches a radius at which it can no longer stably orbit with slowly decreasing radius and it begins a rapid inward plunge. For the inspiral of two roughly equal mass holes it can be imagined that the binary gradually spirals inward or that it reaches a point at which a discontinuous plunge begins. If the late orbits are being degraded rapidly enough by the emission of gravitational radiation, there might not even be any meaning to late stage “stability.” This uncertainty about even the qualitative nature of the late stage of the inspiral is related to an important, but totally unresolved, question: How does the inspiraling binary shed enough angular momentum to form a black hole? In a relativist’s units in which $`c=G=1`$, a black hole must have a total angular momentum $`J`$ that is limited by the maximum angular momentum $`J=M^2`$ that a rotating (Kerr) black hole can possess. Until the binary pair is close, its angular momentum will be above this limit, but technical considerations limit the rate at which angular momentum can be shed in gravitational waves at very late stages. If both black holes of the pair are rapidly rotating with angular momentum in the same direction the shedding appears to present a barrier to the formation of the final single black hole. It is possible that even the qualitative details of the late stage inspiral depend on the angular momentum of the inspiraling binary. The set of possibilities is considerable and the answers are important both to our understanding of nonlinear gravitational interactions and to an understanding of gravitational wave sources. Real answers will require advances in numerical relativity that will be several years in coming, but interest in the questions justifies approximation methods that can help, even slightly, to close some of the wide open questions. We take such an approach here. We offer an estimate of the gravitational radiation generated during the late stage of inspiral of two black holes. Our method involves a number of assumptions and limitations that constrain its applicability and reliability, but for all its shortcomings it is one step towards a complete understanding. The approximation method we use, the “close limit,” takes advantage of the property of a black hole horizon. Late in the merger of the binary the single horizon of the final black hole engulfs the entire binary. All the complex structure of the binary will be inside that final horizon, and cannot influence spacetime outside the horizon. It is only what is outside the horizon that can generate gravitational waves that can be detected by distant observers. Since the ultimate fate of the merger is a stationary black hole, it follows that sufficiently late in the merger what is outside the hole will be a perturbation of the final stationary hole. Thus the gravitational waves generated during the latest stage of inspiral can be computed using the techniques of perturbations of black hole spacetimes, with the Zerilli, Regge-Wheeler, and Teukolsky equations. To understand how the close limit method is to be used, it is necessary to consider the general problem addressed by numerical relativity. Einstein’s field equations are divided into “initial value” equations and equations of time evolution . The initial value equations determine the nature of spacetime at a chosen initial moment. The solutions of these initial value equations are the initial values for the remaining differential equations of Einstein’s theory, the equations that determine the spacetime (including its gravitational wave content) to the future of the initial time. The two tasks of numerical relativity are first to find an initial value solution representing a moment in the life of the colliding holes, and second to find the future spacetime for those initial values. The more computationally difficult task is that of evolving to the future and the codes that accomplish this task tend to be unstable for long time evolutions. For long evolutions to be avoided, the initial value solutions must be chosen to be a moment late in the life of the inspiral. If that moment is late enough, the close limit method can be brought to bear and evolution can be carried out with the stable linearized equations of perturbation theory. But choosing too late a starting moment for evolution creates a new difficulty. The connection of an initial value solution to a “sensible” physical configuration for the binary is reasonably secure only if the binary pair is well separated. At close separations, the gravitational field of each of the binary holes strongly affects the other hole, and the individual mass, individual angular momentum, and physical separation of the holes lose clear meaning. The problem then requires navigating between the Scylla of numerical instabilities for evolution, and the Charybdis of uncertain initial conditions. By using a very late initial moment and linearized evolution, the close limit method completely avoids the former hazard. There are reasons beyond speculation to believe that close limit evolutions give useful answers. Numerical relativity results are available for axisymmetric head-on collisions. These represent evolution of a number of initial value solutions, in particular the closed form solution due to Misner, containing a single parameter representing the initial separation of equal mass holes in units of the mass of the spacetime (in $`c=G=1`$ units). This separation index defines a parameterized family of initial value solutions. Choices of this parameter can be made corresponding to large or small initial separation. When numerical relativity and close limit results are compared it is seen that agreement is excellent for small initial separations, and is surprisingly good even when the initial configuration is not close enough for a horizon to engulf the entire binary. Arguments can be made also, that the gravitational waves calculated in the late stage of inspiral are not highly sensitive to details of initial data. Particularly interesting in this regard is work by Abrahams and Cook. In the past several years the close limit method has been extensively studied for head-on collisions of boosted and spinning holes and compared with the results of numerical relativity. Most notably, second order perturbation theory has been developed for the close limit method. In this process of comparison much has been learned about the strengths and limitations of the close limit method, with the goal of applying the method to problems that cannot yet be handled with numerical relativity. The present work represents the first example of this. We report here the results of the application of the close limit method for the three dimensional problem of the late stage inspiral of two black holes. We will use the close limit method for the initial data families constructed by the Bowen and York method and the associated “punctures” families . It is known that these families possess an artificial radiation content when one considers black holes that are close, but such content is also known to be moderate . An important advantage of these methods is that they are typically the starting point for numerical relativity, and thus close limit evolution of these starting points can be compared with the numerical evolution of these same initial data when such evolutions become available. The most important disadvantage, for our purposes, is that the Bowen-York family does not include the Kerr solution, the solution for a rotating hole. This precludes finding a family of initial value solutions that goes, in the limit of small initial separation, to a Kerr black hole. With Bowen-York initial value solutions, then, we cannot consider a collision that will result in a rapidly rotating hole. Rather, we limit our attention to collisions involving a modest amount of total angular momentum and consider the angular momentum as well as the initial separation to be a perturbation of a nonrotating final hole. It it should also be mentioned that currently fashionable astrophysical scenarios suggest that the individual holes might not carry a significant amount of spin in realistic black hole collisions. The organization of this paper is as follows: in the next section we review the method for obtaining the initial data and describe the approximations involved. In the following two sections we discuss how to set up the perturbative formalism geared towards evolution. Since the collisions have net angular momentum we will evolve them both as a perturbation of a rotating and a non-rotating black hole. The comparison of both approaches is given in the subsequent section and we will see that insight is gained by treating the problem in two different ways. We end with a discussion of the results in terms of waveforms and radiated energies and we describe a puzzle in the calculation of the angular momentum radiated. For the reader who wishes to be spared all the details, we summarize our results in a brief punchline: the final ringdown of the inspiraling collision of two non-spinning black holes is unlikely to radiate more than $`1\%`$ of the mass of the system or more than $`0.1\%`$ of its angular momentum in gravitational waves. ## II Initial data To evolve a spacetime in general relativity, one needs to provide initial data, a 3-geometry $`g_{ab}`$ and an extrinsic curvature $`K_{ab}`$, that solve Einstein’s equations on some starting hypersurface (i.e., at some starting time). For two black holes, this is an easy task if the holes are far apart, since one can superpose the solutions for two individual holes ignoring their interactions. When the black holes are close on the initial hypersurface, the astrophysically correct initial data is the solution corresponding to what would have evolved during the binary inspiral, but such an evolution cannot be computed with present day techniques. One must therefore use a somewhat artificial initial data solution that is a best guess at a representation of close black holes. The need for such a guess is one of the sources of uncertainty in our result. ### A Summary of the Bowen–York construction: The initial value equations for general relativity are, $`^a(K_{ab}g_{ab}K)`$ $`=`$ $`0`$ (1) $`{}_{}{}^{3}RK_{ab}K^{ab}+K^2`$ $`=`$ $`0`$ (2) where $`g_{ab}`$ is the spatial metric, $`K_{ab}`$ is the extrinsic curvature and $`{}_{}{}^{3}R`$ is the scalar curvature of the three metric. If we propose a 3-metric that is conformally flat $`g_{ab}=\varphi ^4\widehat{g}_{ab}`$, with $`\widehat{g}_{ab}`$ a flat metric, and $`\varphi ^4`$ the conformal factor, and we use a decomposition of the extrinsic curvature $`K_{ab}=\varphi ^2\widehat{K}_{ab}`$, and assume maximal slicing $`K_a^a=0`$, the constraints become, $`\widehat{}^a\widehat{K}_{ab}`$ $`=`$ $`0`$ (3) $`\widehat{}^2\varphi `$ $`=`$ $`{\displaystyle \frac{1}{8}}\varphi ^7\widehat{K}_{ab}\widehat{K}^{ab},`$ (4) where $`\widehat{}`$ is a flat-space covariant derivative. To solve the momentum constraint, we start with a solution that represents a single hole with linear momentum $`P`$ , $$\widehat{K}_{ab}^{\mathrm{one}}=\frac{3}{2r^2}\left[2P_{(a}n_{b)}(\delta _{ab}n_an_b)P^cn_c\right].$$ (5) In this expression for the conformally related extrinsic curvature at some point $`x^a`$, the quantity $`n_b`$ is a unit vector, in the “base” flat space with metric $`\widehat{g}_{ab}`$, directed from a point representing the location of the hole to the point $`x^a`$. The symbol $`r`$ represents the distance, in the flat base space, from the point of the hole to $`x^a`$. It is straightforward to show that the solution of the Hamiltonian constraint corresponding to eq. (5) corresponds to a spacetime with ADM momentum $`P_a`$. The next step is to modify this to represent holes centered at $`x=\pm L/2`$ in the conformally flat metric. Since the momentum constraint is linear, we can simply add two expressions of the above form, $$\widehat{K}_{ab}^{\mathrm{two}}=\widehat{K}_{ab}^{\mathrm{one}}\left(xxL/2,P_y=P\right)+\widehat{K}_{ab}^{\mathrm{one}}\left(xx+L/2,P_y=P\right).$$ (6) We will choose in further expressions to use a polar coordinate system in the flat space determined by $`\widehat{g}_{ab}`$ centered in the mid-point separating the two holes and label the polar coordinates as $`(R,\theta ,\varphi )`$. So $`R`$ will be the distance in the flat space from the midpoint between the holes. To solve the Hamiltonian constraint 4, we introduce an approximation, (the slow approximation) which we will show is enough for our purposes. In fact, in this approximation the solution for the conformal factor turns out to be the familiar Misner solution if one chooses the topology of the slice to have a single asymptotically flat region, or the Brill–Lindquist solution if there are three asymptotically flat regions. ### B The slow approximation We assume that the black holes are initially close, and that the initial momentum $`P`$ is small. We denote by $`\stackrel{}{n}^+`$ and $`\stackrel{}{n}^{}`$ the normal vectors corresponding, respectively, to the one hole solutions at $`x=+L/2`$ and at $`x=L/2`$, and we define $`R`$ to be the distance to a field point, in the flat conformal space, from the point midway between the holes. For large $`R`$, the normal vectors $`\stackrel{}{n}^+`$ and $`\stackrel{}{n}^{}`$ almost cancel. More specifically $`\stackrel{}{n}^+=\stackrel{}{n}^{}+O(L/R)`$. A consequence of this is that the total initial $`\widehat{K}^{ab}`$ is first order in $`L/R`$, and its ($`R,\theta ,\phi `$ coordinate basis) components can be written as $$\widehat{K}_{ab}=\frac{3PL}{8R^3}\left[\begin{array}{ccc}8\mathrm{sin}^2\theta \mathrm{sin}2\phi & 0& 8R\mathrm{sin}^2\theta \\ 0& R^2(5+\mathrm{cos}2\theta )\mathrm{sin}2\phi & 2R^2\mathrm{sin}2\phi \mathrm{sin}2\theta \\ 8R\mathrm{sin}^2\theta & 2R^2\mathrm{sin}2\phi \mathrm{sin}2\theta & R^2\mathrm{sin}^2\theta \mathrm{sin}2\phi (1+3\mathrm{cos}2\theta )\end{array}\right].$$ (7) This solution for $`\widehat{K}_{ab}`$ is first order, both in $`P`$ and $`L`$. Thus the source term in the Hamiltonian constraint is quadratic in $`P`$. If we choose to find a solution to the conformal factor to first order in $`P`$ (which should give us a good approximation in the case of slowly moving holes), we can ignore this quadratic source term. So now, the Hamiltonian constraint looks like the one for zero momentum, which is simply the Laplace equation. A well known solution to this, is the Misner solution . This solution, is characterized by a parameter $`\mu _0`$ which describes the separation of the two throats. We can relate this parameter to the conformal distance $`L`$ in the following way , $$L/M=\frac{\mathrm{coth}\mu _0}{2\mathrm{\Sigma }_1}\mathrm{\Sigma }_1\underset{n=1}{}\frac{1}{\mathrm{sinh}n\mu _0}.$$ (8) To clarify: in the slow approximation we are considering, the data we use in our simulations consists of the extrinsic curvature proposed by Bowen and York and the conformal factor due to Misner. This might appear as odd, since the conformal factor of Misner is “symmetrized” through the throats and the extrinsic curvature due to Bowen and York is not. What we do is not inconsistent, it is just a different (and perhaps from a certain point of view less natural) choice of boundary conditions for the fields. In practice, in the close limit and to first order in perturbation theory, the conformal factor of Misner differs from that of Brill and Lindquist by a numerical factor that can be absorbed in the definition of the separation of the holes . Some readers may be disturbed by the slow approximation, since in the computation of certain quantities, for instance the ADM mass, the higher order terms in the expansion in terms of the momentum are crucial. We have already discussed this in detail in previous head-on simulations . The bottomline is that to get an accurate estimate of the ADM mass for high values of the momentum one indeed needs a full solution of the Hamiltonian constraint and not a “slow approximation” solution. For the values of the separations and the momenta we will consider in this paper ($`a<0.5`$) the ADM mass computed with the slow approximation and the one computed with the full solution differ by less than $`10\%`$ so we will ignore this difference. We must now map the coordinates of the initial value solution to the coordinates for the Schwarzschild/Kerr (in the vanishing spin limit) background. To do this, we interpret the $`R`$ as the isotropic radial coordinate of a Schwarzschild spacetime, and we relate it to the usual Schwarzschild radial coordinate $`r`$ by $`R=(\sqrt{r}+\sqrt{r2M})^2/4`$. From this we arrive at the following expression for the components of the extrinsic curvature, $$K_{ab}=\frac{3PL}{8r^3}\left[\begin{array}{ccc}\frac{8\mathrm{sin}^2\theta \mathrm{sin}2\phi }{12M/r}& 0& \frac{8r\mathrm{sin}^2\theta }{\sqrt{12M/r}}\\ 0& r^2(5+\mathrm{cos}2\theta )\mathrm{sin}2\phi & 2r^2\mathrm{sin}2\phi \mathrm{sin}2\theta \\ \frac{8r\mathrm{sin}^2\theta }{\sqrt{12M/r}}& 2r^2\mathrm{sin}2\phi \mathrm{sin}2\theta & r^2\mathrm{sin}^2\theta \mathrm{sin}2\phi (1+3\mathrm{cos}2\theta )\end{array}\right].$$ (9) Here we have used the fact that $$\varphi ^2\varphi _{\mathrm{Mis}}^2\varphi _{\mathrm{Schw}}^2=r/R=\frac{1}{\sqrt{12M/r}}\frac{dr}{dR}.$$ (10) ## III The close limit as a perturbation of a Schwarzschild hole In this paper we will evolve the initial data we just constructed using the perturbative evolution equations for linearized first order perturbations: the Zerilli equation in the case of a Schwarzschild background and the Teukolsky equation in the case of a Kerr background. We need to construct the initial data for these equations in terms of the metric and extrinsic curvature we discussed above. In this section we discuss the setup of initial data and evolution of the problem as a perturbation of a Schwarzschild black hole, using the Zerilli–Regge–Wheeler formalism. ### A Setting up the initial data for the Zerilli function Given the three metric and the extrinsic curvature, one can explicitly construct the zeroth and first order term of a power series expansion in a fiducial time variable $`t`$ of the space-time metric. From this expression one can read off the appropriate coefficients of the multipolar expansion of the metric in the Regge–Wheeler notation. The only nonvanishing perturbations at $`t=0`$ are, $`H_2[\mathrm{}=2,m=\pm 2]`$ $`=`$ $`K[\mathrm{}=2,m=\pm 2]=\sqrt{{\displaystyle \frac{6\pi }{5}}}{\displaystyle \frac{8ML^2}{\sqrt{r}(\sqrt{r}+\sqrt{r2M})^5}}`$ (11) $`H_2[\mathrm{}=2,m=0]`$ $`=`$ $`K[\mathrm{}=2,m=0]=2\sqrt{{\displaystyle \frac{\pi }{5}}}{\displaystyle \frac{8ML^2}{\sqrt{r}(\sqrt{r}+\sqrt{r2M})^5}}`$ (12) We compute the time derivative of these quantities, using the extrinsic curvature $`K_{ij}`$ obtained in the last section. The nonvanishing ones are, $`{\displaystyle \frac{H_2[\mathrm{}=2,m=2]}{t}}`$ $`=`$ $`i24\sqrt{{\displaystyle \frac{\pi }{30}}}PL{\displaystyle \frac{\sqrt{r2M}}{r^3\sqrt{r}}}`$ (13) $`{\displaystyle \frac{H_2[\mathrm{}=2,m=2]}{t}}`$ $`=`$ $`{\displaystyle \frac{H_2[\mathrm{}=2,m=2]}{t}}`$ (14) $`{\displaystyle \frac{K[\mathrm{}=2,m=2]}{t}}`$ $`=`$ $`i\sqrt{30\pi }PL{\displaystyle \frac{\sqrt{r2M}}{r^3\sqrt{r}}}`$ (15) $`{\displaystyle \frac{K[\mathrm{}=2,m=2]}{t}}`$ $`=`$ $`{\displaystyle \frac{K[\mathrm{}=2,m=2]}{t}}`$ (16) $`{\displaystyle \frac{G[\mathrm{}=2,m=2]}{t}}`$ $`=`$ $`i\sqrt{{\displaystyle \frac{6\pi }{5}}}PL{\displaystyle \frac{\sqrt{r2M}}{r^3\sqrt{r}}}`$ (17) $`{\displaystyle \frac{G[\mathrm{}=2,m=2]}{t}}`$ $`=`$ $`{\displaystyle \frac{G[\mathrm{}=2,m=2]}{t}}`$ (18) where $`i`$ is the imaginary unit. (Here we are using the standard conventions for the spherical harmonics. Notice that the $`m=2`$ and $`m=2`$ perturbations are individually complex, but when they are added to give the total perturbation the resulting function of $`t,r,\theta `$ and $`\phi `$ is real, as of course it must be. We also have an odd parity contribution, $$\frac{^{\text{odd}}h_1[\mathrm{}=1,m=0]}{t}=8\sqrt{3\pi }\frac{PL}{r^2}.$$ (19) This perturbation represents the difference between the Kerr solution that represents the rotating space-time and the Schwarzschild background used in the perturbative approach. To first order it decouples from all other perturbations, and in fact is unchanging in time, corresponding to the conservation of angular momentum to first order in the perturbations. The change over time in the quantity induced by second order perturbations, will be discussed below in connection with the radiation of angular momentum. The Zerilli function is defined by (see for instance ), $`\psi _{(\mathrm{},m)}`$ $`=`$ $`{\displaystyle \frac{2r(r2M)}{\mathrm{}(\mathrm{}+1)(\lambda r+3M)}}\left[H_{2}^{}{}_{(\mathrm{},m)}{}^{}r{\displaystyle \frac{K_{(\mathrm{},m)}}{r}}{\displaystyle \frac{r3M}{r2M}}K_{(\mathrm{},m)}\right]`$ (21) $`+{\displaystyle \frac{r^2}{\lambda r+3M}}\left[K_{(\mathrm{},m)}+(r2M)\left({\displaystyle \frac{G_{(\mathrm{},m)}}{r}}{\displaystyle \frac{2}{r^2}}h_{1}^{}{}_{(\mathrm{},m)}{}^{}\right)\right].`$ Therefore, for $`t=0`$ we have $$\psi _{(2,m)}(0,r)=\frac{r(r2M)}{3(2r+3M)}\left[H_{2}^{}{}_{(2,m)}{}^{}(0,r)r\frac{K_{(2,m)}(0,r)}{r}\right]+\frac{r}{3}K_{(2,2)}(0,r).$$ (22) and $`\dot{\psi }_{(2,m)}(0,r)`$ $`=`$ $`{\displaystyle \frac{r(r2M)}{3(2r+3M)}}\left[\dot{H_2}_{(2,m)}(0,r)r{\displaystyle \frac{\dot{K}_{(2,m)}(0,r)}{r}}+3r{\displaystyle \frac{\dot{G}_{(2,m)}}{r}}\right]`$ (24) $`+{\displaystyle \frac{r}{3}}\dot{K}_{(2,2)}(0,r).`$ After some simplifications we have the initial data for the Zerilli function, $$\psi _{(2,2)}(0,r)=4\sqrt{\frac{2\pi }{15}}ML^2\frac{r(7\sqrt{r}+5\sqrt{r2M})}{(2r+3M)(\sqrt{r}+\sqrt{r2M})^5}$$ (25) and $$\dot{\psi }_{(2,2)}(0,r)=\frac{i\sqrt{30\pi }}{5}PL\frac{(4r+3M)\sqrt{r2M}}{r^{5/2}(2r+3M)}$$ (26) and the Zerilli function for $`(\mathrm{}=2,m=2)`$ is the complex conjugate of $`\psi _{(2,2)}(t,r)`$. The initial data for the $`\mathrm{}=2,m=0`$ Zerilli function is $`\psi _{(2,0)}(0,r)`$ $`=`$ $`{\displaystyle \frac{8}{3}}\sqrt{{\displaystyle \frac{\pi }{5}}}{\displaystyle \frac{ML^2r(7\sqrt{r}+5\sqrt{r2M})}{(2r+3M)(\sqrt{r}+\sqrt{r2M})^5}}`$ (27) $`\dot{\psi }_{(2,0)}(0,r)`$ $`=`$ $`0.`$ (28) ### B Evolution of the Zerilli function and computation of physical quantities Given the Cauchy data from the last section, the time evolution is obtained from the Zerilli equation , $$\frac{^2\psi _{(\mathrm{},m)}}{t^2}+\frac{^2\psi _{(\mathrm{},m)}}{r_{}^2}+V(r_{})\psi _{(\mathrm{},m)}=0,$$ (29) where $`V(r_{})`$ is the ($`m`$-independent) Zerilli potential, $$V(r_{})=2\left(1\frac{2M}{r}\right)\frac{\lambda ^2r^2\left[(\lambda +1)r+3M\right]+9M^2(\lambda r+M)}{r^3(\lambda r+3M)^2},$$ (30) where $`\lambda =(\mathrm{}1)(\mathrm{}+2)/2`$ and $`r_{}=r+2M\mathrm{ln}(r/2M1)`$. We need to establish a convenient formula for the radiated energy, similar to that present in but applied to the non-axisymmetric case. We start from the expression of the radiated energy computed via the Landau–Lifshitz pseudo-tensor following the notation and derivations of , $$\frac{d\mathrm{Power}}{d\mathrm{\Omega }}=\underset{r\mathrm{}}{lim}\frac{1}{16\pi r^2}\left[\left(\frac{\dot{h}_{\theta \varphi }}{\mathrm{sin}\theta }\right)^2+\frac{1}{4}\left(\dot{h}_{\theta \theta }\frac{1}{\mathrm{sin}^2\theta }\dot{h}_{\varphi \varphi }\right)^2\right],$$ (31) and translating to the Regge–Wheeler notation and integrating on solid angles we get, $$\mathrm{Power}=\frac{3}{16\pi }\left[2\left\|\dot{\psi }_{(2,0)}\right\|^2+4\left\|\dot{\psi }_{(2,2)}\right\|^2\right]$$ (32) and one can obtain the radiated energy integrating over time. The power naturally comes out in units of the mass of the background spacetime. To compute the radiated angular momentum one could also start by considering the Landau–Lifshitz pseudo-tensor and construct and asymptotic expression for angular momentum flux. This approach was pursued, for instance, in to compute expressions for the radiation of angular momentum in terms of multipoles. An alternative approach is to simply compute the change in the angular momentum of the spacetime, which we characterize to linear order in perturbation theory through the function, $$r_{}^{2}{}_{}{}^{\text{odd}}h_0,_r(r,t)2r^{\text{odd}}h_0(r,t)r^2^{\text{odd}}h_1,_t(r,t).$$ (33) This is a first order gauge invariant if $`^{\text{odd}}h_0`$, and $`^{\text{odd}}h_1`$ are first order perturbations. Moreover, for $`\mathrm{}=1,m=0`$, this gauge invariant is constant, equal to $`4\sqrt{3\pi }J`$, where $`J`$ is the total angular momentum, if the perturbations are axially symmetric. If we look at second order perturbations we find $$\frac{}{t}[r_{}^{2}{}_{}{}^{\text{odd}}h_0,_r(r,t)2r^{\text{odd}}h_0(r,t)r_{}^{2}{}_{}{}^{\text{odd}}h_1,_t(r,t)]=𝒮_{\text{Jdot}}$$ (34) where $`𝒮_{\text{Jdot}}`$ is a ‘source’, quadratic in first order perturbations. Therefore the change in angular momentum, due to radiation may be obtained by integrating $`𝒮_{\text{Jdot}}`$ for all $`t`$ (or from $`t=0`$ to $`t=\mathrm{}`$, it makes no difference), in the limit $`r\mathrm{}`$. After several simplifications and cancelling terms that result from integration by parts, we end up with $$\mathrm{\Delta }J=\frac{3i}{4\pi }\underset{r\mathrm{}}{lim}_0^{\mathrm{}}\left[\psi _{(2,2)}(r,t)\frac{\psi _{(2,2)}(r,t)}{t}\psi _{(2,2)}(r,t)\frac{\psi _{(2,2)}(r,t)}{t}\right]𝑑t$$ (35) If we write $$\psi _{(2,2)}(r,t)=\text{Re}(\psi )+i\text{Im}(\psi )$$ (36) we have $$\psi _{(2,2)}(r,t)=\text{Re}(\psi )i\text{Im}(\psi )$$ (37) and we find $$\mathrm{\Delta }J=\frac{3}{2\pi }\underset{r\mathrm{}}{lim}_0^{\mathrm{}}\left[\text{Im}(\psi )\frac{\text{Re}(\psi )}{t}\text{Re}(\psi )\frac{\text{Im}(\psi )}{t}\right]𝑑t.$$ (38) We have checked by explicit substitution that this form coincides with the results from the flux formulas of Thorne . It reassures our confidence in the consistency of the Regge–Wheeler–Zerilli perturbative formalism to notice that the changes to second order are in accordance with the first order flux. ## IV Evolution as a perturbation of a Kerr black hole To treat the problem as a perturbation of a Kerr black hole we need to set up initial data and evolve the Teukolsky equation. The formalism for setting up initial data in terms of Cauchy metric data was developed in , we only give a brief sketch here and refer the reader to that paper for further details. The relevant Weyl scalar for gravitational radiation is $$\psi _4=C_{\alpha \beta \gamma \delta }n^\alpha \overline{m}^\beta n^\gamma \overline{m}^\delta ,$$ (39) since it is directly related to outgoing gravitational waves. We can rewrite this as $$\psi _4=R_{ijkl}n^i\overline{m}^jn^k\overline{m}^l+4R_{0jkl}n^{[0}\overline{m}^{j]}n^k\overline{m}^l+4R_{0j0l}n^{[0}\overline{m}^{j]}n^{[0}\overline{m}^{l]},$$ (40) which in turn can be written in terms of hypersurface quantities $`g_{ij}`$ and $`K_{ij}`$. For the last term in this expression, we can use vacuum Einstein equations to eliminate terms that have time derivatives of $`K_{ij}`$. Also, we are interested merely in the first order perturbations of this scalar. Putting all this together, the final result for the first order expansion of the Weyl scalar is , $`\psi _4`$ $`=`$ $`\left[R_{ijkl}+2K_{i[k}K_{l]j}\right]_{(1)}n^i\overline{m}^jn^k\overline{m}^l4N_{(0)}\left[K_{j[k,l]}+{}_{j[k}{}^{p}K_{l]p}^{}\right]_{(1)}n^{[0}\overline{m}^{j]}n^k\overline{m}^l`$ (42) $`+4N_{(0)}^2\left[R_{jl}K_{jp}K_l^p+KK_{jl}\right]_{(1)}n^{[0}\overline{m}^{j]}n^{[0}\overline{m}^{l]}`$ where $`N_{(0)}=(g_{\text{kerr}}^{tt})^{1/2}`$ is the zeroth order lapse, $`n^i,\overline{m}^j`$ are two of the null vectors of the (zeroth order) tetrad, Latin indices run from 1 to 3, and the brackets are computed to only first order (zeroth order excluded). This expression can be used, to obtain the time derivative of the Weyl scalar too. We simply replace the first order quantities above by their time derivatives (which can be obtained via the Einstein equations). In our treatment, the extrinsic curvature and the metric, from the last section, shall be treated as a perturbation of the corresponding Kerr hypersurface quantities. Since we attempt calculations only to first order in $`PL`$ (which we identify with $`Ma`$, where $`M`$ is the mass of the background Kerr black hole and $`a`$ its angular momentum parameter), the Kerr 3-metric is (in this approximation) conformally flat. Hence we justify using the Bowen York recipe for constructing initial data for the inspiral problem. ### A Initial Data for the Teukolsky function: Using the methodology and expressions we just discussed, the initial data for the Teukolsky function, $`\mathrm{\Psi }=\rho ^4\psi _4`$, where $`\rho =1/(ria\mathrm{cos}\theta )`$, is: For the azimuthal modes, $`m=\pm 2`$ $$\frac{\mathrm{\Psi }}{\sqrt{2\pi }}=\left[\frac{3rM(2Mr)L^2}{32R^2(2R+M)}\pm i\frac{3}{8}Ma\left(1\frac{2M}{r}\right)^{\frac{3}{2}}\right](\mathrm{cos}\theta \pm 1)^2$$ (43) $$\frac{\dot{\mathrm{\Psi }}}{\sqrt{2\pi }}=\left[\frac{3(2Mr)M^2L^2}{16rR^2(2R+M)}\pm i\frac{3Ma}{16r^2}(2r21M)\left(1\frac{2M}{r}\right)^{\frac{3}{2}}\right](\mathrm{cos}\theta \pm 1)^2$$ (44) And for the azimuthal mode, $`m=0`$ $$\frac{\mathrm{\Psi }}{\sqrt{2\pi }}=\frac{3rM(r2M)L^2}{16R^2(2R+M)}\mathrm{sin}^2\theta $$ (45) $$\dot{\mathrm{\Psi }}=\frac{2M}{r^2}\mathrm{\Psi }$$ (46) Here, $`R`$ is the Schwarzschild isotropic radial coordinate. ### B Evolution of the Data using the Teukolsky equation Given the Cauchy data from the last section, the time evolution is obtained from the Teukolsky equation , $`\left\{\right[a^2\mathrm{sin}^2\theta {\displaystyle \frac{(r^2+a^2)^2}{\mathrm{\Delta }}}]_{tt}{\displaystyle \frac{4Mar}{\mathrm{\Delta }}}_{t\phi }+4[r+ia\mathrm{cos}\theta {\displaystyle \frac{M(r^2a^2)}{\mathrm{\Delta }}}]_t`$ (47) $`+\mathrm{\Delta }^2_r\left(\mathrm{\Delta }^1_r\right)+{\displaystyle \frac{1}{\mathrm{sin}\theta }}_\theta \left(\mathrm{sin}\theta _\theta \right)+\left[{\displaystyle \frac{1}{\mathrm{sin}^2\theta }}{\displaystyle \frac{a^2}{\mathrm{\Delta }}}\right]_{\phi \phi }`$ (48) $`\mathrm{\hspace{0.17em}4}[{\displaystyle \frac{a(rM)}{\mathrm{\Delta }}}+{\displaystyle \frac{i\mathrm{cos}\theta }{\mathrm{sin}^2\theta }}]_\phi (4\mathrm{cot}^2\theta +2)\}\mathrm{\Psi }=0,`$ (49) where $`M`$ is the mass of the black hole, $`a`$ its angular momentum per unit mass, $`\mathrm{\Sigma }r^2+a^2\mathrm{cos}^2\theta `$, and $`\mathrm{\Delta }r^22Mr+a^2`$. The radiated energy is given by , $$\frac{dE}{dt}=\underset{r\mathrm{}}{lim}\left\{\frac{1}{4\pi r^6}_\mathrm{\Omega }𝑑\mathrm{\Omega }\left\|_{\mathrm{}}^t𝑑\stackrel{~}{t}\mathrm{\Psi }(\stackrel{~}{t},r,\theta ,\phi )\right\|^2\right\},d\mathrm{\Omega }=\mathrm{sin}\theta d\vartheta d\phi ,$$ (50) and the angular momentum carried away by the waves can be obtained from , $$\frac{dJ_z}{dt}=\underset{r\mathrm{}}{lim}\left\{\frac{1}{4\pi r^6}\mathrm{Re}\left[_\mathrm{\Omega }𝑑\mathrm{\Omega }\left(_\phi _{\mathrm{}}^t𝑑\stackrel{~}{t}\mathrm{\Psi }(\stackrel{~}{t},r,\theta ,\phi )\right)\left(_{\mathrm{}}^t𝑑t^{}_{\mathrm{}}^t^{}𝑑\stackrel{~}{t}\overline{\mathrm{\Psi }}(\stackrel{~}{t},r,\theta ,\phi )\right)\right]\right\}.$$ (51) ## V Results of the evolutions We have evolved the Zerilli and Teukolsky equations using codes that have already been tested in other situations , . Figure 1 shows the amplitude of the waves, depicting the “+” component of the polarization, defined in terms of the Zerilli function as, $$h_+=\frac{1}{r^2}\left(h_{\theta \theta }\frac{1}{\mathrm{sin}^2\theta }h_{\varphi \varphi }\right)=\frac{\sqrt{30}}{4\sqrt{\pi }r}\left(\sqrt{6}\psi _{(2,0)}+(2\mathrm{sin}^2\theta )\left[\psi _{(2,2)}e^{2i\varphi }+\psi _{(2,2)}e^{2i\varphi }\right]\right).$$ (52) In figure 2 we give a spatial visualization of the waves, by plotting the “$`\times `$” polarization of the waves, defined as, $$h_\times =\frac{1}{r^2}h_{\theta \varphi }=\frac{i\sqrt{30}}{4r\sqrt{\pi }}\mathrm{cos}\theta \mathrm{sin}\theta \left(\psi _{(2,2)}e^{2i\varphi }\psi _{(2,2)}e^{2i\varphi }\right),$$ (53) The figure suggests a rotation pattern, but as can be seen in the accompanying movie, the shown patterns just propagate outward. Let us turn now to the evaluation of the radiated energies and angular momentum. Figure 3 shows the radiated energy as a function of the initial angular momentum, for a fixed separation of the holes. The figure compares the Regge–Wheeler–Zerilli (Z) and Teukolsky (T) calculations. As expected, they differ for large values of the angular momentum, since the Teukolsky calculation contains terms higher than linear in the angular momentum. As we explained before, one is not keeping consistently these higher order terms so one cannot argue that the Teukolsky result is “better”. A conservative view that can be taken should be that both results disagree when higher order terms start to be important, and this gives us a rough measure of the error in the Zerilli calculation. We therefore conclude that for the separation in question, one should not trust first order perturbation theory beyond $`a=0.5`$. One should stress that this view can be somewhat overconservative, our experience with explicit second order calculations for the head-on collisions shows that one should include all second order terms to have a consistent formulation and a reliable set of “error bars”. This is not accomplished by the first order Teukolsky formalism in this context. In this respect, second order Teukolsky results for this problem will be quite welcome . The second order Zerilli calculations appear as quite prohibitive in complexity. The separation of the holes quoted in figure 3 requires some explanation. The simulations start with the construction of the initial data by the Bowen-York procedure we described in section II A. As discussed there, the construction starts with the introduction of a fiducial conformal space. In such a space the separation is $`0.91M`$ where $`M`$ is the ADM mass of the spacetime. The radius of each hole (if they were non-moving, the momentum slightly changes the shape of the horizon and the radius, see ) is approximately $`M/4`$, from there the separation of $`3.64M`$ quoted in the caption. To translate to more commonly used terms, one could convert the number to the $`\mu _0`$ parameter in the Misner solution, which for our case is $`\mu _0=1.5`$. Finally, another commonly used measure of the separation is the length of the geodesic threading the throat in the Misner geometry. In terms of such a parameter, it is equivalent to $`2.75`$ times the ADM mass of the spacetime or approximately $`5.5`$ times the mass of each individual hole. A remarkable aspect of figure 3 is that linear perturbation theory has a tendency to overestimate the radiated energies for large values of the perturbative parameter, at least from our experience with head-on collisions of non-boosted , boosted and some preliminary unpublished results we have for spinning holes. This would suggest that, if the same behavior takes place for the inspiralling collisions, “reality” should lie below the curve corresponding to the Regge–Wheeler–Zerilli formalism (Z). This would indicate that the estimation obtained using the Teukolsky formalism is actually worse for the particular kind of collision under consideration. This is what we were alluding to when we warned in the introduction that it was not obvious that representing the spacetime as a perturbation of a non-rotating hole was a worse choice than of that of a Kerr hole. We now turn to the evaluation of the radiated angular momentum. This is depicted in figure 4. The two curves shown in 4 disagree significantly. They do not even agree for very small values of the angular momentum. We have checked that there is no numerical error: if in the Teukolsky evolution one keeps the initial data intact but “turns off” the $`a`$-dependent terms in the evolution equation, the RWZ straight line is reproduced. It should be noticed that the radiated angular momentum is a qualitatively different quantity insofar as its computation than the energy. The energy is roughly obtained by squaring and integrating the waveforms. The angular momentum depends on subtle phase differences. It is much more easy to disturb the calculation of the radiated angular momentum than that of the radiated energy. This, in particular, points out to the potential difficulty of estimating this kind of quantity in full numerical simulations, where phase lags in the waveforms due to grid stretching and other problems are well known. In our approach it appears that the potentially inconsistent higher order terms in the angular momentum we introduce when considering a rotating background are causing problems in the computation of radiated angular momentum. If one wishes to be ultra-conservative, one could simply conclude that both calculations only predict the correct result for zero angular momentum. Otherwise, one could conclude that for this family of initial data the Teukolsky approach really only works for non-rotating black holes, something suggested by the fact that the background spacetime is only recovered in the close limit with vanishing angular momentum. At the moment we can only say that the accurate computation of the radiated angular momentum for this problem is an open problem. It is likely that the RWZ estimate is correct, but we do not have “error bars” (even rough ones) to validate this prediction. ## VI Conclusions We have used the “close limit” to estimate the radiation in the collision at the end of the inspiral of two equal mass nonrotating black holes. The assumptions and restrictions were: (i) only the “ringdown” radiation was computed; (ii) we assumed that a simple initial data set gave an adequate representation of appropriate astrophysical conditions; (iii) we assumed that the final hole is not near the extreme Kerr limit; (iv) we used close limit estimates of the evolution. Our main conclusion is that the energy radiated in ringdown is probably not more than 1% of the total mass of the system, and the angular momentum radiated is not more than 0.1% of the initial angular momentum. The most serious uncertainty in this result is the possibility that the radiation from the early merger stage of coalescence is very much larger than the ringdown radiation. With our 1%$`Mc^2`$ estimate, collisions of black holes of $`100M_{}`$ would be detectable with signal to noise of 6 out to distances on the order of 200Mpc by the initial LIGO configuration and to distances of 4Gpc with the advanced LIGO detector. ## VII Acknowledgements We wish to thank two anonymous referees for many constructive criticisms on the initial submitted version of the paper. This work was supported in part by grants NSF-INT-9512894, NSF-PHY-9357219, NSF-PHY-9423950, NSF-PHY-9734871, NSF-PHY-9800973, NSF-PHY-9407194, by funds of the University of Córdoba, Utah, and Penn State. We also acknowledge support of CONICET and CONICOR (Argentina). JP acknowledges support from the the John S. Guggenheim foundation and hospitality from ITP at UC Santa Barbara during completion of the manuscript. RJG is a member of CONICET.
warning/0003/cs0003051.html	ar5iv	text	# Local Diagnosis ## Introduction In (?), we have shown how consistency-based diagnosis relates to belief revision and how Reiter’s algorithm can be used for kernel operations of belief change. In (?), we presented the idea of local change and characterized operations of belief change that only affect the relevant part of a belief base. In (?), we presented an algorithm for retrieving the relevant part of a belief base which can be used for implementing local change. In the present paper, we close the circle by showing how local change can be used for focusing the diagnosis process on the relevant part of the domain. We will show how a diagnosis problem can be translated into an operation of kernel semi-revision. Kernel semi-revision (?) consists in adding new information to a database and restoring consistency if necessary. To restore consistency, the expanded database is contracted by $``$. Then we will show how to use information about the structure of the device being examined in order to obtain more efficient methods of diagnosis. For this, we will use the operation of local kernel semi-revision, presented in (?), that considers only the relevant part of the database. In (?), we have presented a simple method for extracting the relevant part of a structured database, which will be used in this paper. This paper proceeds as follows: in the next two sections we give a brief introduction to kernel operations of belief change and consistency-based diagnosis. Then we show the relation between kernel semi-revision and Reiter diagnosis. Using this relation, we show how to use information about the system to focus on its relevant part during the process of diagnosis. In the rest of this paper we consider $`L`$ to be a propositional language closed under the usual truth-functional connectives and containing a constant $``$ denoting falsum. ## Kernel Semi-Revision Hansson introduced a construction for contraction operators, called kernel contraction (?), which is a generalization of the operation of safe contraction defined in (?). The idea behind kernel contraction is that, if we remove from the belief base $`B`$ at least one element of each $`\alpha `$-kernel (minimal subset of $`B`$ that implies $`\alpha `$), then we obtain a belief base that does not imply $`\alpha `$ (?). To perform these removals of elements, we use an incision function, i.e., a function that selects at least one sentence from each kernel. ###### Definition 1 (?) The kernel operation $``$ is the operation such that for every set $`B`$ of formulas and every formula $`\alpha `$, $`XB\alpha `$ if and only if: 1. $`XB`$ 2. $`\alpha Cn(X)`$ 3. for all $`Y`$, if $`YX`$ then $`\alpha Cn(Y)`$ The elements of $`B\alpha `$ are called $`\alpha `$-kernels. ###### Definition 2 (?) An incision function for $`B`$ is any function $`\sigma `$ such that for any formula $`\alpha `$: 1. $`\sigma (B\alpha )(B\alpha )`$, and 2. If $`\mathrm{}XB\alpha `$, then $`X\sigma (B\alpha )\mathrm{}`$. Semi-revision consists in adding new information to a database and restoring consistency if necessary. To restore consistency, the expanded database is contracted by $``$. Semi-revision consists of two steps: first the belief $`\alpha `$ is added to the base, and then the resulting base is consolidated, i.e., contracted by $``$. Kernel semi-revision uses kernel contraction for the second step. ###### Definition 3 (?) The kernel semi-revision of $`B`$ based on an incision function $`\sigma `$ is the operator $`\mathrm{?}_\sigma `$ such that for all sentences $`\alpha `$: $`B\mathrm{?}_\sigma \alpha =(B\{\alpha \})\sigma ((B\{\alpha \}))`$ ## Consistency-Based Diagnosis Diagnosis is a very active area within the artificial intelligence community. The problem of diagnosis consists in, given an observation of an abnormal behavior, finding the components of the system that may have caused the abnormality (?). In the area known as model-based diagnosis (?), a model of the device to be diagnosed is given in some formal language. In this paper, we will concentrate on model-based diagnosis methods that work by trying to restore the consistency of the system description and the observations. Although Reiter’s framework is based on first-order logic, most of the problems studied in the literature do not make use of full first-order logic and can be easily represented in a propositional language. For the sake of simplicity, we will adapt the definitions given in (?) to only mention formulas in the propositional language $`L`$. ### Basic Definitions The systems to be diagnosed will be described by a set of propositional formulas. For each component $`X`$ of the system, we use a propositional variable of the form $`okX`$ to indicate whether the component is working as it should. If there is no evidence that the system is not working, we can assume that variables of the form $`okX`$ are true. ###### Definition 4 A system is a pair (SD,ASS), where: 1. SD, the system description, is a finite set of formulas of $`L`$ and 2. ASS, the set of assumables, is a finite set of propositional variables of the form $`okX`$. An observation is a formula of $`L`$. We will sometimes represent a system by (SD,ASS,OBS), where OBS is an observation for the system (SD,ASS). The need for a diagnosis arises when an abnormal behavior is observed, i.e., when SD$``$ASS$``${OBS} is inconsistent. A diagnosis is a minimal set of assumables that must be negated in order to restore consistency. ###### Definition 5 A diagnosis for (SD,ASS,OBS) is a minimal set $`\mathrm{\Delta }`$ASS such that: SD $``$ {OBS} $``$ ASS$`\mathrm{\Delta }`$ $``$ {$`\neg okX\|okX\mathrm{\Delta }`$} is consistent. A diagnosis for a system does not always exist: ###### Proposition 1 (?) A diagnosis exists for (SD,ASS,OBS) if and only if SD$``${OBS} is consistent. Definition 5 can be simplified as follows: ###### Proposition 2 (?) The set $`\mathrm{\Delta }`$ASS is a diagnosis for (SD,ASS,OBS) if and only if $`\mathrm{\Delta }`$ is a minimal set such that SD $``$ {OBS} $``$ (ASS$`\mathrm{\Delta }`$) is consistent. ### Computing Diagnoses In this section we will present Reiter’s construction for finding diagnoses. Reiter’s method for computing diagnosis makes use of the concepts of conflict sets and hitting sets. A conflict set is a set of assumables that cannot be all true given the observation: ###### Definition 6 (?) A conflict set for (SD,ASS,OBS) is a set Conf $`=\{okX_1,`$ $`okX_2,`$ $`\mathrm{},`$ $`okX_n\}`$ $``$ ASS such that SD $``$ {OBS} $``$ Conf is inconsistent. From Proposition 2 and Definition 6 it follows that $`\mathrm{\Delta }`$ASS is a diagnosis for (SD,ASS,OBS) if and only if $`\mathrm{\Delta }`$ is a minimal set such that ASS$`\mathrm{\Delta }`$ is not a conflict set for (SD,ASS,OBS). A hitting set for a collection of sets is a set that intersects all sets of the collection: ###### Definition 7 (?) Let $`𝒞`$ be a collection of sets. A hitting set for $`𝒞`$ is a set $`H_{S𝒞}S`$ such that for every $`S𝒞`$, $`HS`$ is nonempty. A hitting set for $`𝒞`$ is minimal if and only if no proper subset of it is a hitting set for $`𝒞`$. The following theorem presents a constructive approach for finding diagnoses: ###### Theorem 1 (?) $`\mathrm{\Delta }`$ASS is a diagnosis for (SD,ASS,OBS) if and only if $`\mathrm{\Delta }`$ is a minimal hitting set for the collection of minimal conflict sets for (SD,ASS,OBS). Consider the circuit in Figure 1. The system description of this circuit is given by (SD,ASS), where: $`\text{ASS}=\{okX,okY,okZ\}`$ $`\begin{array}{cc}\text{SD}=\hfill & \{(AB)okXD,\hfill \\ & \neg (AB)okX\neg D,\hfill \\ & CokY\neg E,\hfill \\ & \neg CokYE,\hfill \\ & (DE)okZF,\hfill \\ & \neg (DE)okZ\neg F\}\hfill \end{array}`$ Suppose we have OBS=$`\neg C\neg F`$. This observation is inconsistent with $`\text{SD}\text{ASS}`$. There is only one minimal conflict set for (SD,ASS,OBS): {$`okY,okZ`$}. There are three possible hitting sets: {$`okY`$},{$`okZ`$}, and {$`okY,okZ`$}. Reiter considers only minimal hitting sets as diagnoses, that is, either Y or Z is not working well. ## Diagnosis via Kernel Semi-Revision In (?), we have shown that the standard method for finding consistency-based diagnosis, due to Reiter (?), is very similar to the construction of kernel semi-revision, except for the fact that Reiter only considers minimal diagnosis, which correspond to minimal values for incision functions. In this section we summarize these results. Recall that kernel operations are based on two concepts: kernels and incision functions. The kernels are the minimal subsets of a belief base implying some sentence, while the incision functions are used to decide which elements of the kernels should be given up. Let (SD,ASS,OBS) be a system. The belief base that we are going to semi-revise corresponds to SD$``$ASS and the input sentence is OBS. The conflict sets are the assumables in the inconsistent kernels of SD$``$ASS$``${OBS}. So, if $`B`$=SD$``$ASS, the conflict sets are given by {$`X`$ASS$`\|X`$ ($`B+`$OBS)$``$$``$}. Incision functions correspond loosely to hitting sets, the minimal hitting sets being the values of minimal incisions that return only assumables. Note that there is a difference in the status of formulas in SD and those in ASS: formulas in ASS represent expectations and are more easily retracted than those in SD (cf. Definition 8). We can model the diagnosis problem as a kernel semi-revision by the observation. Semi-revision can be divided in two steps. First the observation is added to the system description together with the assumables. In case the observation is consistent with the system description together with the assumables, no formula has to be given up. Otherwise, we take the inconsistent kernels and use an incision function to choose which elements of the kernels should be given up. In the case of diagnosis, we do not wish to give up sentences belonging to the system description or the observation. We prefer to give up the formulas of the form $`okX`$, where $`X`$ is a component of the system. Moreover, we are interested in minimal diagnosis, so the incision should be minimal. For this, we use a special variant of incision function. We modify Definition 2 so that incisions are minimal and elements of a given set $`A`$ are prefered over the others: ###### Definition 8 Given a set $`A`$, an $`A`$-minimal incision function is any function $`\sigma _A`$ from sets of sets of formulas into sets of formulas such that for any set $`S`$ of sets of formulas: 1. $`\sigma _A(S)S`$, 2. If $`\mathrm{}XS`$, then $`X\sigma _A(S)\mathrm{}`$, 3. If for all $`XS`$, $`XA\mathrm{}`$, then $`\sigma _A(S)A`$, and 4. $`\sigma _A(S)`$ is a minimal set satisfying 1,2, and 3. If we take $`A`$ to be the set of assumables, we obtain an incision function that prefers to select formulas of the form $`okX`$ over the others. We can show that for (SD,ASS,OBS), whenever a diagnosis exists, an ASS-minimal incision function will select only elements of ASS: ###### Proposition 3 Let (SD,ASS,OBS) be a system with an observation and $`\sigma _{ASS}`$ an ASS-minimal incision function. If a diagnosis exists, then $`\sigma _{ASS}((\text{SD}\text{ASS}\text{OBS}))`$ASS. ###### Lemma 1 The assumables that occur in an inconsistent kernel of the set SD$``$ASS$``$OBS form a conflict set for (SD,ASS,OBS) and all minimal conflict sets can be obtained in this way, i.e.: (i) For every $`X(`$SD$``$ASS$``$OBS$`)`$, $`X`$ASS is a conflict set, and (ii) For every minimal conflict set $`Y`$, there is some $`X(`$SD$``$ASS$``$OBS$`)`$ such that $`X`$ASS$`=Y`$. Note that not every inconsistent kernel determines a minimal conflict set, since for conflict sets only the elements of ASS matter, i.e., there may be two inconsistent kernels $`X_1`$ and $`X_2`$ such that $`X_1`$ASS is a proper subset of $`X_2`$ASS. Recall that given an incision function $`\sigma `$, the semi-revision of a set $`B`$ by a formula $`\alpha `$ was given by $`B\mathrm{?}_\sigma \alpha =(B+\alpha )\sigma ((B+\alpha ))`$. A diagnosis is given by the elements of ASS that are given up in a kernel semi-revision by the observation. ###### Proposition 4 Let S=(SD,ASS,OBS) be a system and $`\sigma _{ASS}`$ an ASS-minimal incision function. (SD $``$ASS)$``$((SD $``$ASS)$`\mathrm{?}_{\sigma _{ASS}}`$OBS) = $`\sigma _{ASS}`$((SD $``$ASS $``$OBS)$``$$``$) is a diagnosis. ## Using System Structure Suppose that instead of the circuit depicted in Figure 1, we have the circuit in Figure 2. Suppose also that we get the same observation, i.e., OBS$`=\neg C\neg F`$. Intuitively, only a small part of the circuit (roughly the sub-circuit at figure 1) has to be considered in order to arrive to a diagnosis. In (?), we have extended the definition of kernel semi-revision to an operation that considers only the relevant part of a database, local kernel semi-revision. In (?), we have shown how to use structure present in a database in order to find compartments and implement local kernel operations more efficiently. The key idea of the method described is to use a relation of relatedness between formulas of the belief base. In some applications, as we will see, such a relation is given with the problem. In the case of the circuit shown in Figure 2, there is a very natural dependence relation. The output of each of the components depends on the input and on whether the component is working well. The only observation we have is $`\neg C\neg F`$. Since this observation is inconsistent with the system description together with the assumption that all components are working well, there must be some faulty component. Moreover, the faulty component must be in the path between $`C`$ and $`F`$ (of course, there may be other faulty components, but we are only searching for the abnormality that explains the observation). We only need to consider the descriptions of components $`y`$ and $`z`$ in order to find the diagnosis. In the next section we will show how to use the framework described in (?) in order to find diagnoses without having to check the entire system description for consistency. ## Local Kernel Diagnosis As we have seen, diagnosis problems fit very well in the framework for local change that we proposed in (?) and (?). Besides the fact that the traditional method for finding diagnosis based on the notion of consistency is almost identical to the construction of kernel semi-revision, in most diagnosis problems there is a very natural notion of relatedness that can be used to structure the belief base so that the search for diagnoses becomes more efficient. In this section we formalize the example in Figure 2 in order to show how to derive a concrete relatedness relation from the given database. We will use a relatedness relation between atoms, as illustrated in Figure 3. The relation is not symmetric. We can easily adapt the definitions presented in (?) to deal with a directed graph. The basic algorithm is as follows: we start from the propositional variables that occur in the observation and spread the activation in the graph, following the direction of the arcs. The spreading finishes either when the end of the paths are reached or when we run out of resources (time or memory). This is done by the algorithm Retrieve below, an adaptation of the algorithm given in (?). The algorithm uses the function Adjacent to collect all nodes related to a given node, i.e., given a relatedness relation $`R`$, Adjacent($`x`$)=$`\{y`$Var(SD)$``$ASS$`\|R(x,y)\}`$, where Var($`X`$) is the set of propositional variables ocurring in the formulas of set $`X`$. For a set $`Y`$ of propositional variables, Adjacent($`Y`$)= $`\{`$Adjacent($`y`$)$`\|yY\}`$. Retrieve(OBS,ASS,Relevant): 1. For all $`p`$ Var(OBS), mark($`p`$) 2. $`\mathrm{\Delta }^1`$(OBS) := Adjacent(Var(OBS)) 3. Relevant := Var(OBS)$``$ASS 4. $`i`$ := 1; stop := false 5. While not stop do 5.1. For all $`p\mathrm{\Delta }^i`$(OBS), mark($`p`$) If $`p`$ASS, then Relevant := Relevant $`\{p\}`$ 5.2. i := i+1; $`\mathrm{\Delta }^i`$(OBS)=$`\mathrm{}`$ 5.3 For all $`p\mathrm{\Delta }^{i1}`$(OBS), $`\mathrm{\Delta }^i`$(OBS) := $`\mathrm{\Delta }^i`$(OBS) $`\{q`$Adjacent($`p`$) s.t. not marked$`(q)\}`$ 5.4 If $`\mathrm{\Delta }^i`$(OBS)$`=\mathrm{}`$, then stop := true After we have retrieved the relevant assumables, the relevant compartment is taken to be the observation together with all formulas in SD$``$ASS which mention the relevant assumables. Compartment(OBS,SD,ASS,Comp): 1. Retrieve(OBS,ASS,Relevant) 2. Comp=OBS 3. For all $`p`$Relevant, Comp:= Comp$`\{\alpha `$SD$``$ASS$`\|p`$Var($`\alpha `$)$`\}`$. As we have seen in (?), the algorithm for Retrieve is an anytime algorithm. The algorithm for Compartment is not, at least in principle. But if one keeps the order in which the relevant atoms are retrieved and uses them in this order in line 3 of algorithm Compartment, one can be sure that the description of the most relevant components will be retrieved first. For the circuit in Figure 2, we have: SD={$`(AB)okXD`$, $`\neg (AB)okX\neg D`$, $`CokY\neg E`$, $`\neg CokYE`$, $`(DE)okZF`$, $`\neg (DE)okZ\neg F`$, $`G1okW1\neg A`$, $`\neg G1okW1A`$, $`(G2G3)okW2B`$, $`\neg (G2G3)okW2\neg B`$, $`(G4G5)okW3C`$, $`\neg (G4G5)okW3\neg C`$, $`G6okW4\neg G9`$, $`\neg G6okW4G9`$, $`(G7G8)okW5G10`$, $`\neg (G7G8)okW5\neg G10`$, $`(G9G10)okW6G11`$, $`\neg (G9G10)okW6\neg G11`$, $`G11okW7G12`$, $`\neg G11okW7\neg G12`$, $`(FG12)okW8G13`$, $`\neg (FG12)okW8\neg G13`$} ASS ={$`okX,okY,okZ,okW1,okW2,okW3,okW4,`$ $`okW5,okW6,okW7,okW8`$} If we apply the algorithm Retrieve($`\neg C\neg F`$,ASS,Relevant) to the graph depicted in Figure 3, we get Relevant={okY,okZ,okW8}. For Compartment(OBS, SD, ASS, Comp) we get Comp=$`\{\neg C\neg F,CokY\neg E,`$ $`\neg CokYE,(DE)okZF,`$ $`\neg (DE)okZ\neg F,`$ $`(FG12)okW8G13,`$ $`\neg (FG12)okW8\neg G13,okY,okZ,okW8\}`$. The diagnosis can be searched using only the formulas in Comp. Note that the component w8 was not really relevant for the diagnosis but, nevertheless, we have reduced the set to be semi-revised. This is a very general method for focusing on a small part of the system description. One can add to it some domain specific heuristics to improve its efficiency. The system IDEA (?), used by FIAT repair centers works on dependence graphs that show graphically the relation between the several components of a device. In (?) we have shown that Reiter’s algorithm for consistency-based diagnosis can be used for kernel semi-revision. The algorithm for kernel operations can be easily combined with the algorithm Compartment presented in this section. Applying Reiter’s algorithm to Comp, given the observation $`\neg C\neg F`$, we get as possible diagnoses: $`\{okY\}`$ and $`\{okZ\}`$. ## Conclusions In this paper we have shown how to combine Reiter’s algorithm for consistency-based diagnosis with the algorithm for finding the relevant compartment of a database. The result is a method for finding diagnosis which focuses on the relevant part of the system description. Making clear the similarities between diagnosis and belief revision can be very profitable for both areas of research. As shown in (?), the computational tools developed in the field of diagnosis can be adapted to be used for belief revision. And as we show in this paper, theories developed for belief revision can be applied on diagnosis for obtaining more efficient methods. Future work includes the study of other approaches to diagnosis as well as the study of the computational complexity of the method proposed.
warning/0003/hep-th0003255.html	ar5iv	text	# Dimensional Transmutation and Dimensional Regularization in Quantum Mechanics I. General Theory ## I INTRODUCTION It is well known that, for various models of quantum field theory, a mass scale emerges spontaneously through the renormalization procedure, even when the original theory has no dimensional parameters. This phenomenon, called dimensional transmutation, was first analyzed in the 1973 seminal work of Coleman and Weinberg , where the scalar field of massless scalar electrodynamics was shown to develop a nonzero but arbitrary expectation value; as a consequence, the particles of the theory acquire nonzero physical masses . In short, the Coleman-Weinberg mechanism induces radiative corrections to the Higgs potential, thereby suggesting the relevance of dimensional transmutation for the generation of particle masses . The main goal of this paper is to present a thorough investigation of dimensional transmutation in nonrelativistic quantum mechanics, with a threefold purpose in mind: (i) at the conceptual level, to show that quantum field theory is not a prerequisite for its existence; (ii) mathematically, to characterize the class of scale-invariant potentials, as well as the subclass of potentials that display dimensional transmutation; and (iii) at the practical level, to develop useful tools for the treatment of a certain class of singular quantum-mechanical potentials. In particular, our work offers additional insight into two problems that have been extensively studied in the literature: the two-dimensional delta-function potential and the inverse square potential . Parenthetically, the family of delta-function potentials that is discussed in this article is actually included in a larger class of singular potentials whose apparent phenomenological usefulness has been recognized for a long time, since the introduction of pseudopotentials in the early days of quantum mechanics and with subsequent applications of the zero-range potential in nuclear physics , condensed matter physics , statistical mechanics , atomic physics , and particle physics . Likewise, the inverse square potential is related to the dipole potential, which has found applications in molecular physics . Even though earlier research on the subject had relied solely on traditional quantum-mechanical techniques, this situation has changed in recent years, with the introduction of numerous applications of quantum field-theoretic tools and renormalization theory to the same problems. Specifically, among the many applications related directly or indirectly to our singular potentials, the following are worth mentioning: (i) the mathematical formulation of the theory of pseudopotentials , which started with the works of Ref. and includes the technique of self-adjoint extensions ; (ii) the study of the nonrelativistic limit of the $`\varphi ^4`$ theory, with the concomitant question of its triviality ; (iii) the basic conceptual understanding of quantum field theory mechanisms in the simpler framework of quantum mechanics, including standard regularization and renormalization of the singular potentials mentioned above , anomalies , renormalization group analysis , and effective field theory approach ; (iv) the analysis of (2+1)-dimensional theories, including gravity , as well as Chern-Simons theory , the Aharonov-Bohm effect , and the dynamics of anyons ; (v) new applications of contact potentials in condensed matter physics, e.g., for the quantum Hall effect ; and (vi) the modern formulation, using effective field theory , of the nucleon-nucleon potential , which has led to a plethora of quantum-mechanical pseudopotentials . Our paper naturally follows the trend set by this extensive bibliography. The main focus of our work will be on the concept of dimension, a term that has been extensively used in the physics literature to describe two conceptually distinct ideas. The meaning to which dimensional transmutation refers is that which relates to measurement and which characterizes the class of physical quantities that exhibit a certain type of power-law behavior (dimensional homogeneity) with respect to a given choice of fundamental quantities . In what follows, we will use the term dimensionality to denote the exponents of the associated homogeneous behavior for any given physical quantity . At first sight, dimensional transmutation is paradoxical, because it produces a scale in a problem devoid of dimensional parameters, in apparent violation of Buckingham’s $`\mathrm{\Pi }`$ theorem of orthodox dimensional analysis . However, this paradox can be ultimately resolved by invoking the dimensional arbitrariness intrinsic in the renormalization framework . The other widely used acceptation of dimension refers to a geometric concept, a property of the space where events take place (e.g., space-time in relativistic physics). In this paper we will extensively exploit the trivial connection between these two concepts, that is, that the dimensionality of an element of volume in a given space, expressed in terms of units of length, is equal to its geometric dimension. This connection has been largely used in the renormalization of relativistic quantum field theories, where it is further reinforced by the implementation of dimensional regularization . Correspondingly, we will use dimensional regularization as a natural technique that renders obvious the spontaneous generation of a dimensional scale for a scale-invariant theory. Our paper is organized as follows. In Section II we discuss the meaning of dimensional transmutation in terms of dimensional analysis and renormalization theory. Section III is devoted to the concept of scale-invariant potentials in nonrelativistic quantum mechanics, where generic properties related to scaling and scale symmetry are derived. Section IV establishes a general framework for the regularization of scale-invariant potentials based upon dimensional continuation; this procedure is later extended to a renormalization scheme in Section V, where dimensional transmutation is shown to arise in the strong-coupling regime. Finally, an application of the theory is illustrated in Section VI, which deals with some aspects of the two-dimensional delta-function potential. The conclusions of our analysis are presented in Section VII. The appendices deal with the necessary results regarding $`D`$-dimensional Euclidean spaces, Green’s functions, and scattering. Additional applications for rotationally invariant problems will follow in the second paper in this series . ## II DIMENSIONAL ANALYSIS, RENORMALIZATION, AND DIMENSIONAL TRANSMUTATION The description of a physical system in the context of a given theory, either in terms of a Lagrangian or of a Hamiltonian, includes a certain number of parameters. Usually, their values may all be fixed from the start by the laws of nature, in the form of “constants,” but one always enjoys the mathematical freedom to make some of them become variable parameters as needed. Then, for the discussion that follows, these parameters will arbitrarily be classified into two groups: “constant” or fundamental and “variable” or dynamical. By varying the variable parameters one introduces a whole class $`𝒞`$ of physical systems, all characterized by the same values of the “constants” . Fundamental parameters are those that are fixed constants for all the members of the given class of systems. Needless to say, they are dimensional because “fundamental” dimensionless parameters amount to plain numerical constants. Typical fundamental parameters of choice are the dimensional universal constants of nature, for example, $`\mathrm{}`$ and $`c`$. The second group is composed of the dynamical parameters that acquire different values for the different members of the class $`𝒞`$. As an example we could mention the masses of particles or the coupling constants of the interactions. The reduction in the number of dimensionally independent quantities can be accomplished by arbitrarily assigning particular numerical values to a subset of the fundamental parameters. This procedure amounts to the selection of a generalized natural system of units, in which the number of fundamental dimensions is reduced. For example, in relativistic quantum field theory, it is customary to choose $`c=1`$ and $`\mathrm{}=1`$, so that the theory is described in terms of a single fundamental dimension—usually taken as inverse length $`\mathrm{\Lambda }=L^1`$, which is equivalent to mass, momentum, and energy. Even though all the physical dimensions can be restored easily at any stage of the calculation, it is clear that great simplification is achieved in the dimensional analysis of various physical quantities. Similarly, in one-particle nonrelativistic quantum mechanics, one has the freedom to use $`\mathrm{}`$ and $`m`$ (where $`m`$ is the particle’s mass) as fundamental parameters that define a particular generalized natural system of units; in this paper, we will choose $`\mathrm{}=1`$ and $`m=1/2`$. Then, we will be left with only one fundamental dimension, which we will take again as inverse length $`\mathrm{\Lambda }=L^1`$ or momentum. Consequently, in what follows, we will define the inverse-length dimensionality $`q=\mathrm{dim}\left[Q\right]`$ of a physical quantity $`Q`$ as the exponent that expresses its physical dimension $`\mathrm{\Lambda }^q`$ in terms of inverse length, that is, $$q=\mathrm{dim}\left[Q\right]=\frac{\mathrm{\Lambda }}{\left[Q\right]}\frac{\left[Q\right]}{\mathrm{\Lambda }}.$$ (1) For nonrelativistic quantum mechanics, Table I summarizes the dimensionalities of the most common physical quantities. ###### TABLE I . Physical Dimensions of Various Physical Quantities. \| Physical quantity \| Ordinary dimensions \| “Natural” dimensions \| Dimensionality \| \| --- \| --- \| --- \| --- \| \| Length \| $`L`$ \| $`L`$ \| -1 \| \| Time \| $`T`$ \| $`L^2`$ \| -2 \| \| Velocity \| $`LT^1`$ \| $`L^1`$ \| 1 \| \| Linear momentum \| $`MLT^1`$ \| $`L^1`$ \| 1 \| \| Angular momentum \| $`ML^2T^1`$ \| 1 \| 0 \| \| Energy \| $`ML^2T^2`$ \| $`L^2`$ \| 2 \| \| Cross section \| $`L^{D1}`$ \| $`L^{D1}`$ \| $`\left(D1\right)`$ \| \| Wave function (normalized) \| $`L^{D/2}`$ \| $`L^{D/2}`$ \| $`D/2`$ \| Note. “Natural” dimensions are defined by the choice $`\mathrm{}=1`$ and $`2m=1`$. The geometric dimension of position space is $`D`$. Let us now explore the consequences of the possible existence of dimensional parameters. For a given physical system, characterized by a Lagrangian or a Hamiltonian, oftentimes there exists at least one dimensional parameter, which can be used to define a system-specific or intrinsic scale. To illustrate how this is done, let us consider the nonrelativistic quantum-mechanical dynamics of a single particle in one dimension, such that the external-interaction potential contributes only one dimensional parameter $`\lambda `$; for example, for an attractive power-law potential, one may study the possible existence of bound states through the Schrödinger equation $$\left[\frac{d^2}{dx^2}+\mathrm{sgn}\left(\beta \right)\lambda \|x\|^\beta \right]\mathrm{\Psi }(x)=E\mathrm{\Psi }(x).$$ (2) Dimensional analysis shows that $`\mathrm{dim}\left[\lambda \right]=\mathrm{}=2+\beta `$; then $`\lambda ^{1/\mathrm{}}`$ will define a basic unit of inverse length or momentum. Any dimensional quantity Q of dimensionality q will then be equal to $`\lambda ^{q/\mathrm{}}`$, up to a numerical factor; similarly, a function $`Q(x)`$ of position (or $`Q(p)`$ of momentum), with dimensionality q, will then be equal to $`\lambda ^{q/\mathrm{}}`$ times a dimensionless function of $`\lambda ^{1/\mathrm{}}x`$ (or of $`\lambda ^{1/\mathrm{}}p`$). In particular, a characteristic ground-state energy may be estimated as $`\lambda ^{2/(2+\beta )}`$. In other words, dimensional analysis gives nontrivial information about the system. The obvious statements of the previous paragraph can be summarized in the $`\mathrm{\Pi }`$ theorem of dimensional analysis , which we state here without proof, with the intention of generalizing it later in this section. Consider a physical phenomenon described by $`M`$ dimensional characteristic parameters $`a_1,\mathrm{},a_M`$, such that $`R`$ of them are dimensionally independent. Then, given an equation $$F(a_1,\mathrm{},a_M)=0$$ (3) involving these $`M`$ parameters, there exist $`N`$ independent dimensionless power products $`\mathrm{\Pi }_1,\mathrm{},\mathrm{\Pi }_N`$ of $`a_1,\mathrm{},a_M`$, such that Eq. (3) is equivalent to $$\mathrm{\Phi }(\mathrm{\Pi }_1,\mathrm{},\mathrm{\Pi }_N)=0,$$ (4) with $$N=MR.$$ (5) For example, working in a natural system of units with only one independent dimension, it follows that $`N=M1`$, which describes the situation of the previous paragraph. But what happens if the system exhibits no explicit dimensional dynamical parameter at the level of the Lagrangian or Hamiltonian? As we will see in Section III, such a system is scale-invariant. An example is the power-law potential $`\lambda \|x\|^2`$ ($`\beta =2)`$ because its coupling $`\lambda `$ is dimensionless. Then, naive dimensional analysis is at a loss to make any meaningful predictions. In this case, if a new scale arises (for example, a bound state under a scale-invariant potential), dimensional analysis implies that it has the following properties: * It is spontaneously generated, in the sense that it characterizes the solution of a theory that is scale-invariant at the level of the classical Lagrangian or Hamiltonian. This amounts to an instance of quantum-mechanical breaking of classical scale symmetry–also called scale anomaly . * It is totally arbitrary because no privileged value is defined a priori within the theory. If it were not arbitrary, it would violate the $`\mathrm{\Pi }`$ theorem in an irreconcilable way. This manifestation of an arbitrary and spontaneously generated scale in a scale-invariant theory is known as dimensional transmutation . In short, in the solution to a well-posed question within the scale-invariant theory, a dimensionally transmuted scale $`B`$ may appear spontaneously. How does it come into existence in apparent violation of naive dimensional analysis? Our goal is to disentangle the mechanism that leads to this transmutation. This will be implemented by means of a regularization-renormalization procedure. The regularization technique introduces a dimensional parameter $`\mu `$, in terms of which the scale $`B`$ is expressed. Thus, a dimensional transfer takes place, whereby a dimensionless parameter $`\lambda `$ is “transmuted” into or traded for a dimensional scale $`B`$. This simple process can be represented diagrammatically in the form \| Initial problem \| Technique \| Solution \| \| --- \| --- \| --- \| \| Lagrangian/Hamiltonian \| Regularization/renormalization \| Physical quantity \| \| dimensionless \| arbitrary \| measurable \| \| coupling $`\lambda `$ \| $`\genfrac{}{}{0pt}{}{\text{ }\text{ }}{\text{dimensional scale}\mu }`$ \| dimensional scale $`B`$ \| We conclude this section by stating the modification of orthodox dimensional analysis needed to encompass this anomalous behavior. As discussed in Ref. , the usual assumption underlying the $`\mathrm{\Pi }`$ theorem is that the function $`F(a_1,\mathrm{},a_M)`$ of Eq. (3) is uniquely defined, an assumption that breaks down when the Lagrangian does not describe a single theory but a class of theories parametrized with renormalization parameters. This manifests itself in a theory that is ill-defined or exhibits singularities of some sort, in which case the Lagrangian or Hamiltonian cannot represent a complete description of the physics; thus, renormalization is needed. When the number of independent sliding scales or renormalization parameters is $`\sigma `$, the required modification is of the $`\mathrm{\Pi }`$ theorem is obviously $$N=M+\sigma R.$$ (6) Equation (6) states that the number of “available” variables is $`M^{}=M+\sigma `$, rather than $`M`$; in particular, it provides the necessary freedom to permit the emergence of dimensional transmutation. The framework for deriving conclusions directly from Eq. (6) will be referred to as generalized dimensional analysis. ## III CHARACTERIZATION OF SCALE-INVARIANT POTENTIALS ### A Scale Symmetry and Homogeneity In this section we set out to define and characterize mathematically the class of scale-invariant potentials $`V(𝐫)`$ in one-particle nonrelativistic quantum physics. In a strict sense, we are referring to a physical system whose classical action $$S=\left[\frac{1}{2}m\stackrel{2}{\stackrel{.}{𝐫}}V(𝐫)\right]𝑑t$$ (7) is invariant under the scale transformations $`𝐫\varrho 𝐫`$, $`t\tau t`$ (with $`\varrho >0`$ and $`\tau >0`$). This scale symmetry is satisfied if and only if each one of the two terms of the nonrelativistic action—the kinetic-energy term $`dtm\stackrel{2}{\stackrel{.}{𝐫}}/2`$ and the potential-energy term $`𝑑tV(𝐫)`$—are left unchanged. Due to the spatial and time dependence of the nonrelativistic kinetic-energy term, this invariance condition is satisfied only when $`\varrho ^2=\tau `$ (obviously consistent with the dimensional analysis of Table I), while the invariance of the potential-energy term requires that $$V(\varrho 𝐫)=\varrho ^2V(𝐫).$$ (8) As Eq. (8) is valid for all $`\varrho >0`$, the class of scale-invariant potentials is identical to that of homogeneous potentials of degree -2. As we will see next, this is the same condition to be satisfied when the potential does not exhibit any explicit dimensional scale. ### B Dimensional Scaling in Nonrelativistic Quantum Mechanics One-particle nonrelativistic quantum mechanics in the presence of a stationary potential $`V(𝐫)`$ is described in the $`D`$-dimensional position-space representation of the Schrödinger picture via the solutions of the time-independent Schrödinger equation, which in natural units reads $$\left[^2+V(𝐫)\right]\mathrm{\Psi }(𝐫)=E\mathrm{\Psi }(𝐫).$$ (9) The analysis and interpretation of the solutions to Eq. (9) will become more transparent when the transition to its dimensionless version is carried out. This can be accomplished by rescaling all quantities appearing in Eq. (9) by means of a dimensional parameter $`\mu `$, which we will assume to represent an inverse-length standard; then, $`\mu `$ will satisfy the properties: (i) inverse-length dimensionality, $$\mathrm{dim}\left[\mu \right]=1;$$ (10) (ii) positivity, $$\mu >0.$$ (11) In general, there are many possible characteristic scales that may serve as $`\mu `$: they could either be intrinsic to the system or arbitrary scales introduced via regularization. In any case, we will not be concerned with the multi-scale case, because our ultimate goal is to analyze the extreme scenario where there is no intrinsic dimensional parameter, but an arbitrary sliding scale $`\mu `$ is introduced by the regularization procedure. Then any physical quantity $`Q`$ of dimensionality $`q`$ will be equal to a numerical coefficient times $`\mu ^q`$, whence its dimensionless counterpart will be defined as $`\mu ^qQ`$. If, in addition, the quantity is a function of either position or momentum, it will be of the form $`\mu ^q`$ times a dimensionless function of the dimensionless position $$𝝃=\mu 𝐫,$$ (12) or of the dimensionless momentum $$𝝅=\mu ^1𝐩.$$ (13) Correspondingly, dimensional analysis predicts that the potential energy function (whose dimensionality is 2) should have a dependence on the dimensional parameter $`\mu `$ given by $$V(𝐫,\mu )=\mu ^2𝒱(\mu 𝐫),$$ (14) for arbitrary $`\mu `$. In particular, one can define the reduced function $`𝒱(𝝃)`$, via Eq. (14), in the form $$𝒱(𝝃)=\mu ^2V(\mu ^1𝝃,\mu ),$$ (15) or straightforwardly by specializing to the unit value of the dimensional parameter, i.e., $$𝒱(𝝃)=V(𝝃,\mu =1).$$ (16) The right-hand side of Eq. (14) displays the two sources of possible scale dependence of the potential: the one associated with the dimensionality of $`V`$ as a potential energy (i.e., $`\mu ^2`$) and the one associated with the functional form of the potential (as described by $`𝒱(\mu 𝐫)`$). Equation (14) implies that $$\mu \frac{\|V(𝐫,\mu )\|}{\mu }=2\|V(𝐫,\mu )\|+𝐫\mathbf{}\|V(𝐫,\mu )\|,$$ (17) where the first term on the right-hand side is the dimensionality of the potential energy and the second term represents the degree of the functional dependence of the potential energy with respect to the given scale. For example, for a power-law potential $`V(𝐫)r^\beta `$, the functional dependence amounts to $`𝒱(𝝃)\xi ^\beta `$, whence $`V(𝐫,\mu )\mu ^{2+\beta }`$, which describes the total scale dependence of the potential energy function under arbitrary rescaling. Notice that, for $`\beta =2`$, $`V(𝐫,\mu )`$ is independent of $`\mu `$; i.e., it is scale-independent (see the next section). Table II gives a list of the various dimensionless quantities of interest. ###### TABLE II . Dimensionless Counterparts of Various Physical Quantities in Nonrelativistic Quantum Mechanics. \| Physical quantity \| Symbol \| Dimensionality \| Dimensionless form \| \| --- \| --- \| --- \| --- \| \| Position \| $`𝐫`$ \| -1 \| $`𝝃=\mu 𝐫`$ \| \| Linear momentum \| $`𝐩`$ \| 1 \| $`𝝅=\mu ^1𝐩`$ \| \| Kinetic energy \| $`_𝐫^2`$ \| 2 \| $`_𝝃^2=\mu ^2_𝐫^2`$ \| \| Potential energy \| $`V(𝐫)`$ \| 2 \| $`𝒱(𝝃)`$ = $`\mu ^2V(\mu ^1𝝃,\mu )`$ \| \| Energy \| $`E`$ \| 2 \| $`\eta =\mu ^2E`$ \| \| Wave function (normalized) \| $`\mathrm{\Psi }(𝐫)`$ \| $`D/2`$ \| $`\mathrm{\Phi }(𝝃)=\mu ^{D/2}\mathrm{\Psi }(\mu ^1𝝃)`$ \| Rescaling of Eq. (9) with the parameter $`\mu `$ yields $$\left[_𝝃^2+𝒱(𝝃)\right]\mathrm{\Phi }(𝝃)=\eta \mathrm{\Phi }(𝝃),$$ (18) which describes an eigenvalue problem for the dimensionless eigenfunctions $`\mathrm{\Phi }(𝝃)`$, with dimensionless eigenvalues $$\eta =\mu ^2E,$$ (19) in a space of arbitrary dimension $`D`$. In addition, it is convenient to normalize the wave function $`\mathrm{\Phi }(𝝃)`$ with respect to its dimensionless argument $`𝝃`$, that is, $$\|\mathrm{\Phi }(𝝃)\|^2d^D\xi =\|\mathrm{\Psi }(\mu ^1𝝃)\|^2d^Dr,$$ (20) a condition that yields the rescaling $$\mathrm{\Phi }(𝝃)=\mu ^{D/2}\mathrm{\Psi }(\mu ^1𝝃).$$ (21) As usual, the solutions to Eq. (18) should be separately obtained and interpreted for the bound-state and scattering sectors. Additional conclusions about these specific problems will be drawn in Subsections V B, V C, and V D. ### C Absence of Explicit Dimensional Scales and Homogeneity Revisited Let us now characterize the class of potentials that do not exhibit any dimensional scale, in a space of arbitrary dimension $`D`$. We will resort to the general framework developed in Subsection III B, where we considered potentials that depend upon only one dimensional parameter or none (once the choice $`\mathrm{}=1`$ and $`2m=1`$ has been made); from Eq. (14), the position and dimensional dependence of the potential energy are such that $$V(\mu 𝐫,\mu =1)=\mu ^2V(𝐫,\mu ),$$ (22) for arbitrary $`\mu `$. In general, the function $`V(𝐫,\mu )`$ may have a nontrivial (i.e., not quadratic) dependence with respect to the parameter $`\mu `$, as displayed in Eq. (17). However, those potentials devoid of explicit dimensional scales are independent of any dimensional parameter $`\mu `$; then their functional dependence is of the form $`V=V(𝐫)`$ rather than of the form $`V=V(𝐫,\mu )`$. We can also say that the second argument in $`V=V(𝐫,\mu )`$ is actually spurious. An explicit mathematical statement of this independence is $$\frac{}{\mu }V(𝐫,\mu )=0.$$ (23) Even though this may be taken as the primary definition, it is convenient to derive a more illuminating form by just eliminating the spurious $`\mu `$ dependence in Eq. (22), i.e., $$V(\mu 𝐫)=\mu ^2V(𝐫).$$ (24) Equation (24), valid for $`\mu >0`$ (Eq. (11)), is identical to our earlier homogeneous property (8). An alternative derivation of this remarkably simple property can be obtained directly from Eq. (17) or by differentiation of Eq. (22) with respect to $`\mu `$; then (after setting the arbitrary scale equal to unity), $$𝐫\mathbf{}V(𝐫)=2V(𝐫),$$ (25) a relation that amounts to Euler’s theorem for a homogeneous function of degree $`2`$. In conclusion, the classes of scale-invariant potentials and potentials without any explicit dimensional scale are identical. So far our discussion has only focused on the dimensions but not on the magnitude of the potentials. It is now due time to introduce a dynamical coupling parameter $`\lambda `$ to characterize the “strength” of a given potential, according to $$V(𝐫)=\lambda W(𝐫),$$ (26) where $`W(𝐫)`$ is a homogeneous function of degree $`2`$. The homogeneity displayed by Eq. (24) has a straightforward consequence on the position dependence of scale-invariant potentials. In effect, writing $`𝐫=r\widehat{𝐫}`$, where $`\widehat{𝐫}`$ is the unit position vector, we conclude that $$V(𝐫)=r^2V(\widehat{𝐫}),$$ (27) so that the general scale-invariant potential has either one of the following two forms: (i) A generalized inverse square potential in any number $`D`$ of dimensions, $$V(𝐫)=\lambda \frac{v(\mathrm{\Omega }^{(D)})}{r^2},$$ (28) where the dimensionless function $`v(\mathrm{\Omega }^{(D)})`$ explicitly depends upon the $`D`$-dimensional solid angle $`\mathrm{\Omega }^{(D)}\widehat{𝐫}`$. In particular, $`v=1`$ corresponds to the ordinary inverse square potential , $$V(𝐫)=\frac{\lambda }{r^2},$$ (29) and $`v=\mathrm{cos}\theta `$ amounts to the dipole potential , $$V(𝐫)=\lambda \frac{\mathrm{cos}\theta }{r^2},$$ (30) where $`\theta `$ is the polar angle (measured from the orientation of the dipole moment) and $`\lambda `$ is proportional to the magnitude of the dipole moment. (ii) A homogeneous pseudopotential of degree $`2`$, for which the most general form is $$V(𝐫)=\lambda r^{D2}\delta ^{(D)}(𝐫).$$ (31) It should be noticed that Eq. (31) can be transformed into a number of alternative forms involving the radial delta function; in particular, it is equivalent to $`\delta (r)/r`$ and $`\delta ^{}(r)`$ . In this class, the two-dimensional delta-function potential $$V(𝐫)=\lambda \delta ^{(2)}(𝐫)$$ (32) is the best known example, which we will analyze in Section VI in order to illustrate our general theory. Equation (27) provides the limiting form of the scale-invariant potential at infinity, $$\underset{r\mathrm{}}{lim}V(𝐫)=0$$ (33) (which is identically true for pseudopotentials, for any $`r0`$). Equation (33) then implies that all states with $`E0`$ are scattering states for any scale-invariant potential. The sign of $`\lambda `$ in Eqs. (26)–(32) has been chosen so that $`\lambda >0`$ corresponds to an attractive potential, wherever this concept is applicable. ### D Scale Invariance and Eigenvalue Spectrum The Hamiltonian $`H_𝐫`$ associated with a scale-invariant potential, $$H_𝐫=^2\lambda W(𝐫),$$ (34) is homogeneous of degree $`2`$ with respect to the position vector, that is, $$H_{\zeta 𝐫}=\zeta ^2H_𝐫$$ (35) for arbitrary $`\zeta `$, because both the kinetic and potential energies have the same property. From our analysis of the last section, this is another way of saying that $$_𝝃=\mu ^2H_{𝝃/\mu }=_{\text{ }𝝃}^2\lambda 𝒲(𝝃),$$ (36) where $$𝒲(𝝃)=\mu ^2W(𝝃/\mu )$$ (37) (cf. Eq. (15)), is a dimensionless Hamiltonian totally independent of any explicit dimensional scale. A straightforward consequence of the scale invariance implied by Eq. (35) is the breaking of the discrete character of the bound-state energy spectrum for attractive potentials. Let us now show this. First, the Hamiltonian (36) associated with a scale-invariant potential is an example of a local operator. As is well known, the representative $`A_𝐫`$ of an abstract local operator $`A`$ is defined via $$A_𝐫\delta ^{(D)}(𝐫𝐫^{})=𝐫\|A\|𝐫^{},$$ (38) which implies that $$A_𝐫\psi (𝐫)=𝐫\|A\|\psi .$$ (39) Equation (35) refers to a particular case of a local operator with scale dimension $`s`$, defined via $$A_{\zeta 𝐫}=\zeta ^sA_𝐫,$$ (40) for some exponent $`s`$ and for arbitrary scalar $`\zeta >0`$. With arguments similar to those of Subsection III C, it is easy to see that $`s=\mathrm{dim}\left[A_𝐫\right]`$. Thus, the Hamiltonian of scale-invariant potentials is a local operator with scale dimension $`s=2`$. The statement we wish to prove is that, for any local operator with scale dimension $`s`$, the spectrum can only be continuous. In effect, if $`\psi _a(𝐫)`$ is an eigenfunction with eigenvalue $`a`$, $$A_𝐫\psi _a(𝐫)=a\psi _a(𝐫),$$ (41) then the rescaling $$𝐫^{}=\zeta 𝐫$$ (42) implies that $$A_𝐫\psi _a(\zeta 𝐫)=\zeta ^sa\psi _a(\zeta 𝐫),$$ (43) whence $`\psi _a(\zeta 𝐫)`$ is an eigenfunction of the same operator with eigenvalue $`\zeta ^sa`$. As this should be so for arbitrary $`\zeta `$, one concludes that: (i) If $`a`$ is a finite eigenvalue, then all real numbers of the same sign are eigenvalues. (ii) The eigenvalue spectrum is continuous. Of course, this is what should be expected on intuitive grounds, because if the spectrum were discrete then one would be able to identify preferential scales where none is defined a priori. For example, the rescaling $`𝐫^{}=\zeta 𝐫`$ for the plane-wave eigenstates $`e^{i𝐤𝐫}`$ of the scale-invariant free-particle Hamiltonian amounts to a rescaling of the corresponding momentum $`𝐤^{}=\zeta 𝐤`$, whence all positive energies $`E=k^2=\zeta ^2k^2`$ can be reached by continuously varying the parameter $`\zeta `$. Therefore, for the Hamiltonian of attractive scale-invariant potentials, the corresponding implications are (i) The energy spectrum is either not bounded from below or, if it is, it can only start at $`E=0`$ and be of the scattering type. (ii) The bound-state energy spectrum, if it exists, is continuous. Thus, it is clear that, for a given unregularized scale-invariant potential, in addition to the continuous scattering spectrum with energies $`0E<\mathrm{}`$ (as required by Eq. (33)), there are only three possibilities for the bound-state spectrum: 1. Spectrum devoid of bound states. 2. Continuous bound-state spectrum with energies between $`E=\mathrm{}`$ and $`E=0`$. 3. Singular bound-state spectrum with a unique energy level at $`E=\mathrm{}`$. The fact that no other cases are possible is just a consequence of the scale symmetry. In other words, in the first category the potential is so “weak” that if fails to generate any bound states; in fact, this situation is familiar: it is characteristic of any repulsive potential. The subtlety lies in that a “weak” attractive potential behaves in all respects as a repulsive one: it only has a continuous scattering spectrum extending from $`E=0`$ to $`E=\mathrm{}`$. As an alternative, the potential may be so “strong” that it produces an example of the second category, where it breaks down both the existence of a lower bound and the discrete nature of the bound-state spectrum. It is well known that both cases (categories 1 and 2) are realized by the inverse square potential with weak and strong coupling, respectively . Yet, the singular case of category 3 provides an alternative for a “strong” potential—a behavior exhibited by the two-dimensional delta-function potential . With regard to scattering, one may use the fact that the theory is scale-invariant and that the poles of the scattering matrix on the imaginary energy axis correspond to the bound states. These facts imply that: (i) When there are no bound states, the scattering matrix has no poles and is manifestly scale-invariant, i.e., independent of the energy. (ii) When the spectrum is singular (either categories 2 or 3), the scattering matrix exhibits the corresponding singular behavior. A remark about the need for renormalization is in order. A theory that produces no bound states and a scale-invariant S-matrix (category 1) needs no regularization. In effect, such theory displays no divergence whatsoever and, as we have seen, its spectrum is identical in every respect to that of repulsive potentials, so that scale-invariance is maintained even in the quantum-mechanical theory. Instead, regularization and renormalization are needed for cases 2 and 3 above, an issue to which we now turn our attention. ## IV DIMENSIONAL TRANSMUTATION VIA DIMENSIONAL REGULARIZATION ### A Dimensional Regularization of Scale-Invariant Potentials Let us consider a scale-invariant potential $`V(𝐫)`$ in $`D_0`$ dimensions. The corresponding $`D_0`$-dimensional Schrödinger equation is $$\left[_{𝐫,D_0}^2+V(𝐫)\right]\mathrm{\Psi }(𝐫)=E\mathrm{\Psi }(𝐫),$$ (44) where $`_{𝐫,D_0}^2`$ is the $`D_0`$-dimensional Laplacian. We have seen in Subsection III D that if $`V(𝐫)`$ is of the scale-invariant type, then its unregularized bound-state spectrum is either nonexistent or not bounded from below. Therefore, the difficulty here resides in that, in the initial dimension $`D_0`$, the problem is singular and has to be regularized. In this paper we will use dimensional regularization, a technique originally developed for quantum field theory and which we now adapt to nonrelativistic quantum mechanics. The $`D`$-dimensional generalization of Eq. (44) is of the form $$\left[_{𝐫,D}^2+V^{(D)}(𝐫)\right]\mathrm{\Psi }(𝐫)=E\mathrm{\Psi }(𝐫),$$ (45) where $$V^{(D)}(𝐫)=\lambda _BW^{(D)}(𝐫),$$ (46) with $`\lambda _B`$ being the bare coupling constant (see Subsection IV B), while $`W^{(D)}(𝐫)`$ is an appropriate generalization to $`D`$ dimensions of the original $`D_0`$-dimensional potential, with the only constraint $$\underset{DD_0}{lim}W^{(D)}(𝐫)=W(𝐫),$$ (47) and with $`W(𝐫)`$ defined by Eq. (26). Of course, the whole purpose of this regularization is to produce a scenario where the generalized potential is no longer scale-invariant in a dimension close but not equal to $`D_0`$. In other words, we require that, for arbitrary $`ϵ0`$, $`V^{(D_0ϵ)}(𝐫)`$ not be a homogeneous function of degree $`2`$. Even though the requirement above allows infinitely many possible generalizations, a particularly simple prescription can be developed by the use of Fourier transforms, just like it is done in quantum field theory. It turns out that a simple property of Fourier analysis immediately suggests the generalization: if $`f(𝐫)`$ is a $`D`$-dimensional homogeneous function of degree $`\beta `$, then its Fourier transform $`\stackrel{~}{f}(𝐤)`$ is homogeneous of degree $`(D+\beta )`$. With this property in mind, we define the Fourier transform of the original scale-invariant potential, $$\stackrel{~}{W}(𝐤)=_{(D_0)}\left\{W(𝐫)\right\}=d^{D_0}re^{i𝐤𝐫}W(𝐫),$$ (48) which we analytically continue to $`D`$ dimensions with the prescription that, in Fourier space, the $`D`$-dimensional functional form should be the same as the $`D_0`$-dimensional functional form, i.e., $$\stackrel{~}{W}^{(D)}(𝐤)=\stackrel{~}{W}(𝐤).$$ (49) Finally a $`D`$-dimensional inverse Fourier transform $`_{(D)}^1`$ provides the desired generalization in the position representation, i.e, $$W^{(D)}(𝐫)=_{(D)}^1\left\{\stackrel{~}{W}(𝐤)\right\}=\frac{d^Dk}{(2\pi )^D}e^{i𝐤𝐫}\stackrel{~}{W}(𝐤),$$ (50) where all the vectors are now $`D`$-dimensional. The process represented by Eqs. (48) and (50) involves a dimensional continuation that can be summarized in the following succinct expression for the potential $$W^{(D)}(𝐫_D)=\frac{d^Dk_D}{(2\pi )^D}e^{i𝐤_D𝐫_D}\left[d^{D_0}r_{D_0}e^{i𝐤_{D_0}𝐫_{D_0}}W(𝐫_{D_0})\right]_{𝐤_{D_0}𝐤_D},$$ (51) where the subscripts $`D`$ and $`D_0`$ explicitly indicate the dimension of the corresponding vector and the symbol $`𝐤_{D_0}𝐤_D`$ stands for the dimensional “jump” in momentum space that defines the dimensionally continued potential. The whole process can be represented by means of the commutative diagram $$\begin{array}{cccc}& \text{real space}& & \text{reciprocal space}\\ D_0\text{ dimensions:}& W(𝐫)=W^{(D_0)}(𝐫)& \stackrel{_{\left(D_0\right)}}{}& \stackrel{~}{W}(𝐤)=\stackrel{~}{W}^{(D_0)}(𝐤)\\ & 𝒟_{D_0D}& & 𝒟_{D_0D}\\ D\text{ dimensions: }& W^{(D)}(𝐫)& \stackrel{_{\left(D\right)}^1}{}& \stackrel{~}{W}^{(D)}(𝐤)=\stackrel{~}{W}(𝐤)\end{array},$$ (52) where $`𝒟_{D_0D}`$ is a shorthand for dimensional continuation from $`D_0`$ to $`D`$ dimensions. Correspondingly, the degree of homogeneity of a scale-invariant potential is transformed according to $$\begin{array}{ccc}\stackrel{W(𝐫)}{\mathrm{degree}=2}& \stackrel{_{\left(D_0\right)}}{}& \stackrel{\stackrel{~}{W}(𝐤)}{\mathrm{degree}=2D_0}\\ 𝒟_{D_0D}& & 𝒟_{D_0D}\\ \stackrel{W^{\left(D\right)}(𝐫)}{\mathrm{degree}=2+ϵ}& \stackrel{_{\left(D\right)}^1}{}& \stackrel{\stackrel{~}{W}^{\left(D\right)}(𝐤)}{\mathrm{degree}=2D_0}\end{array},$$ (53) where $$ϵ=D_0D.$$ (54) Diagram (53) explicitly shows that the $`D`$-dimensional real-space continuation of the potential is not of the scale-invariant type, because its degree of homogeneity is $`2+ϵ`$ (with $`ϵ0`$), rather than $`2`$. For example, when the criteria above are applied to the two-dimensional delta-function and inverse square potentials, one obtains the dimensional continuations summarized in Table III . ###### TABLE III . Dimensional Continuations $`W^{(D)}(𝐫)`$ for the Two-Dimensional Delta-Function and Inverse Square Potentials. \| Potential $`W(𝐫)`$ \| Dimension $`D_0`$ \| Dimensional continuation $`W^{(D)}(𝐫)`$ \| \| --- \| --- \| --- \| \| $`\delta ^{(2)}(𝐫)`$ \| 2 \| $`\delta ^{(D)}(𝐫)\|_{D=2ϵ}`$ \| \| $`r^2`$ \| Arbitrary \| $`\pi ^{ϵ/2}\mathrm{\Gamma }(1ϵ/2)r^{(2ϵ)}`$ \| ### B Dimensional Transmutation of the Coupling Parameter In the analysis above we have defined the appropriate functional dependence of the dimensionally continued potential but we have not spelled out the dimensionality change experienced by the coupling parameter. This change occurs in Eq. (51), when the dimensional “jump” $`𝐤_{D_0}𝐤_D`$ is performed, according to $$\left[W^{(D)}(𝐫_D)\right]=\mathrm{\Lambda }^{DD_0}\left[W^{(D_0)}(𝐫_{D_0})\right].$$ (55) Equation (55) amounts to the “creation” of a physical dimension $`L^{(DD_0)}=L^ϵ`$. If the physical dimensions of the potential energy are to remain unchanged, then the bare coupling constant $`\lambda _B`$ should acquire the physical dimensions $`L^ϵ`$, that is, its dimensionality should become $`\mathrm{dim}\left[\lambda _B\right]=ϵ`$. A convenient way of parametrizing this dimensionality change is by the introduction of an arbitrary inverse-length scale $`\mu `$, i.e., by the replacement $$\lambda \lambda _B=\lambda \mu ^ϵ,$$ (56) where $`\lambda `$ is dimensionless. When this dimensionality change is made explicit in Eq. (51), we obtain a dimensionally continued potential $`V^{(D)}(𝐫)`$ $`=`$ $`\lambda \mu ^ϵW^{(D)}(𝐫),`$ (57) $`=`$ $`\lambda \mu ^ϵ{\displaystyle \frac{d^Dk_D}{(2\pi )^D}e^{i𝐤_D𝐫_D}\left[d^{D_0}r_{D_0}e^{i𝐤_{D_0}𝐫_{D_0}}V(𝐫_{D_0})\right]_{𝐤_{D_0}𝐤_D}},`$ (58) i.e., the dimensional jump is performed simultaneously with the introduction of the factor $`\mu ^ϵ`$. Notice that Eqs. (51) and (58), as well as diagram (53) show that even though $$\mathrm{dim}\left[V^{(D)}(𝐫)\right]=2,$$ (59) in fact $$\mathrm{dim}\left[W^{(D)}(𝐫)\right]=2ϵ.$$ (60) Thus, one can write $$W^{(D)}(𝐫)=\mu ^{2ϵ}𝒲^{(D)}(\mu 𝐫),$$ (61) where $`𝒲^{(D)}(𝝃)`$ is the dimensionless counterpart of $`W^{(D)}(𝐫)`$; then, from Eqs. (15), (57), and (61), $$𝒱^{(D)}(𝝃)=\lambda 𝒲^{(D)}(𝝃).$$ (62) The dimensionality change represented by Eq. (56) introduces a completely arbitrary dimensional scale $`\mu `$. The replacement of a dimensionless coupling constant by an arbitrary dimensional scale is the phenomenon of dimensional transmutation seen from the dimensional-regularization viewpoint. In addition, Eq. (56) indicates that the space dimension $`D`$ plays a pivotal role in the determination of the physical dimensions of the coupling constant. Equations (45) and (57) imply that the dimensionally regularized Schrödinger equation has the explicit form $$\left[_{𝐫,D}^2\lambda \mu ^ϵW^{(D)}(𝐫)\right]\mathrm{\Psi }(𝐫)=E\mathrm{\Psi }(𝐫),$$ (63) where $`W^{(D)}(𝐫)`$ is a homogeneous function of degree $`2+ϵ`$. Ultimately, in order to make the potential less singular, it is necessary to choose $`ϵ>0`$ (that is, $`D<D_0`$), whence proper regularization is achieved in the limit $$ϵ=0^+.$$ (64) Alternatively, the dimensionally regularized Schrödinger equation (63) can be rewritten in the dimensionless form of Eq. (18), i.e., $$\left[_{𝝃,D}^2\lambda 𝒲^{(D)}(𝝃)\right]\mathrm{\Phi }(𝝃)=\eta \mathrm{\Phi }(𝝃).$$ (65) If the problem posed by Eq. (65) with the limit (64) is regular, then its solution provides regular eigenfunctions $`\mathrm{\Phi }(𝝃)`$ corresponding to the eigenvalues $`\eta `$. These eigenvalues depend upon the dimensionless parameters $`ϵ`$ and $`\lambda `$. Unlike the dependence of $`\eta `$ on $`ϵ`$, which is potential-dependent, the dependence of $`\eta `$ on $`\lambda `$ is the same for all scale-invariant potentials. This can be understood as follows: 1. The Schrödinger equation (65) is dimensionless and independent of $`\mu `$. 2. Instead, the Schrödinger equation (63) is more explicit in that it displays the dimensional scales $`E`$ and $`\mu `$. In addition, $`\lambda `$ and $`\mu `$ appear in it only in its second term and through the bare coupling constant $`\lambda _B`$, Eq. (56). This can be made more explicit by rewriting Eq. (63) in terms of $`\lambda _B`$, i.e., $$\left[_{𝐫,D}^2\lambda _BW^{(D)}(𝐫)\right]\mathrm{\Psi }(𝐫)=E\mathrm{\Psi }(𝐫).$$ (66) 3. Equation (56) is the basis for the introduction of an effective inverse length $$\widehat{\mu }=\lambda _B^{1/ϵ}=\lambda ^{1/ϵ}\mu ,$$ (67) where $`\lambda >0`$ is assumed because we only need to regularize the potential when it is attractive. Equation (67) allows for the rescaling of the energy $$\widehat{\eta }=\widehat{\mu }^2E=\lambda ^{2/ϵ}\eta ,$$ (68) of the potential energy $$\widehat{𝒲}=\widehat{\mu }^{(2ϵ)}W=\lambda ^{2/ϵ+1}𝒲,$$ (69) as well as of the dimensionless position $$\widehat{𝝃}=\widehat{\mu }𝐫=\lambda ^{1/ϵ}𝝃.$$ (70) Then, the Schrödinger equation (66) takes the form $$\left[_{\widehat{𝝃},D}^2\widehat{𝒲}^{(D)}(\widehat{𝝃})\right]\widehat{\mathrm{\Phi }}(\widehat{𝝃})=\widehat{\eta }\widehat{\mathrm{\Phi }}(\widehat{𝝃}).$$ (71) Equations (63), (65), and (71) are equivalent before taking the limit $`ϵ=0^+`$. When actually solving the Schrödinger equation for a specific potential, it is more straightforward to use (63) because it makes all the variables explicit, while (65) is just a convenient way of relating the dimensional scales from scratch. On the other hand, at the conceptual level, Eq. (71) provides the “universal” connection between $`\lambda `$ and $`\eta `$ for all scale-invariant potentials. In other words, Eq. (68) defines the relationship between the mathematical eigenvalues $`\widehat{\eta }`$ of Eq. (71) and the physical or dimensional eigenvalues $`E`$ of Eq. (63), which are explicit functions of the parameters $`\mu `$ and $`\lambda `$. A convenient form of this “universal” condition satisfied by the energy eigenvalues of Eq. (71) reads (from Eq. (68)) $$\lambda \mathrm{\Xi }(ϵ)\|\eta (ϵ)\|^{ϵ/2}=1,$$ (72) which, in terms of dimensional variables, explicitly states that $$\lambda \mu ^ϵ\mathrm{\Xi }(ϵ)\|E(ϵ)\|^{ϵ/2}=1.$$ (73) Any of the Eqs. (72) and (73) would be referred to as the “master eigenvalue equation,” which provides the required energies if the mathematical function $$\mathrm{\Xi }(ϵ)=\left\|\widehat{\eta }\right\|^{ϵ/2}0$$ (74) is known. In fact, Eq. (74) shows that $`\mathrm{\Xi }(ϵ)`$ is completely determined by the functional form of the potential $`\widehat{𝒲}^{(D)}(\widehat{𝝃})`$, through the solution of the differential equation (71). We shall refer to $`\mathrm{\Xi }(ϵ)`$ as the “energy generating function.” In this paper, we will exemplify the regularization procedure and computation of the energy generating function for the two-dimensional delta-function potential. A more detailed discussion of this potential, as well as the corresponding treatment of the inverse square potential, can be found in the second paper of this series . ## V RENORMALIZATION OF SCALE-INVARIANT POTENTIALS ### A Regularized Bound-State Sector The theory of Subsection IV B applies equally well to both bound and scattering states. In both cases, however, we may assume an attractive potential $`\lambda >0`$, which is the type that possibly requires renormalization. Let us now consider the bound-state sector, for which Eq. (65) provides a discrete sequence of energy eigenvalues $`\eta _n`$, in the regularized version of the theory. These eigenvalues explicitly depend upon the discrete set of quantum numbers $`n=(n_1,\mathrm{},n_D)`$, ordered as an increasing sequence in such a way that, for sufficiently small $`ϵ`$, $`E_n(ϵ)E_n^{}(ϵ)`$ if $`n_jn_j^{}`$ for all $`j=1,\mathrm{},D`$ (additional ordering rules are needed in the presence of degeneracies, but they are immaterial to our discussion). In particular, the ground state will be labeled with the lowest numbers of the sequence. Our goal is to extract additional information from Eq. (72), which we now rewrite $$\lambda \mathrm{\Xi }_n(ϵ)\|\eta _n(ϵ)\|^{ϵ/2}=1,$$ (75) with $`\mathrm{\Xi }_n(ϵ)=\left\|\widehat{\eta }_n\right\|^{ϵ/2}`$. This can be accomplished by defining the variables $$\lambda _n^{()}=\left[\underset{ϵ0}{lim}\mathrm{\Xi }_n(ϵ)\right]^1=\underset{ϵ0}{lim}\left\|\widehat{\eta }_n(ϵ)\right\|^{ϵ/2},$$ (76) such that $$\underset{ϵ0}{lim}\left\|\eta _n(ϵ)\right\|=\underset{ϵ0}{lim}\left[\frac{\lambda }{\lambda _n^{()}}\right]^{2/ϵ}.$$ (77) In Eq. (77) one can see that the unregularized energy is critically dependent on the ratio $`\lambda /\lambda _n^{()}`$, in the limit $`ϵ0`$. Thus, $`\lambda _n^{()}`$ acts as a critical coupling strength for the given energy level labeled by $`n`$. From now on, we will identify the following three regimes: (i) Strong coupling, characterized by $`\lambda >\lambda _n^{()}`$, for which Eq. (77) gives a bound state at $`\mathrm{}`$. (ii) Weak coupling, characterized by $`\lambda <\lambda _n^{()}`$, for which Eq. (77) gets rid of the bound state by pushing it all the way up to $`0`$. (iii) Critical coupling, characterized by $`\lambda =\lambda _n^{()}`$, for which additional analysis is needed. In fact, Eq. (77) implies the following behavior according to the values of the critical coupling: (a) $`\lambda _n^{()}=0`$ amounts to a strong coupling for all finite and positive $`\lambda `$, a condition that is manifested by the “collapse” of the given bound state, $`\eta _n\mathrm{}`$. (b) $`\lambda _n^{()}=\mathrm{}`$ amounts to a weak coupling for all finite $`\lambda >0`$, a condition that is manifested by the loss of the regularized bound state, i.e., $`\eta _n0`$. (c) $`0<\lambda _n^{()}<\mathrm{}`$ permits the existence of the three possible regimes. A few results are implied by the above analysis. First, because of the assumed ordering of quantum numbers, from Eq. (76), it follows that $`\lambda _n^{()}\lambda _n^{}^{()}`$, when $`n_jn_j^{}`$ for all $`j=1,\mathrm{},D`$. In particular, for the ground state, which we will subsequently label with $`(\mathrm{gs})`$, we define the “principal” critical coupling $`\lambda ^{()}`$, which satisfies the condition $$\lambda ^{()}=\lambda _{(\mathrm{gs})}^{()}\lambda _n^{()}.$$ (78) For example, if the coupling is weak for the ground state, it is also weak for all other states, so that the unregularized theory is completely deprived from bound states. The analysis above assumes that $`\lambda `$ is independent of $`ϵ`$ and displays the singular behavior associated with dimensional transmutation as $`ϵ0`$. Renormalization is called for in order to obtain meaningful results. ### B Renormalized Bound-State Sector Equation (75) provides a regularization of the energy levels in terms of the parameter $`ϵ`$. In this section we introduce the general strategy for renormalization and use it to reach a few general conclusions about dimensional transmutation. In order to obtain finite results, it is necessary to renormalize the energy levels by the following procedure: (i) Letting the coupling constant $`\lambda `$ be a function of the regularization parameter $`ϵ`$, i.e., $`\lambda =\lambda (ϵ)`$. (ii) Adjusting $`\lambda (ϵ)`$ by comparison with a specific bound state, which is conveniently chosen to be the ground state of the theory; notice that if bound states exist at all, the ground state is the only one that is guaranteed to exist. We will refer to this procedure as ground-state renormalization. Consequently, in the following analysis, it will prove useful to compare the values of the function $`\mathrm{\Xi }_n(ϵ)`$ with its ground-state value $`\mathrm{\Xi }_{_{(\mathrm{gs})}}(ϵ)`$, by means of the replacement $$\mathrm{\Xi }_n(ϵ)=\mathrm{\Xi }_{_{(\mathrm{gs})}}(ϵ)\left[1+\frac{ϵ}{2}_n(ϵ)\right],$$ (79) which defines a new function $`_n(ϵ)`$, with the obvious implication that $$_{_{(\mathrm{gs})}}(ϵ)=0.$$ (80) Then, the regularized dependence of the energy levels with respect to $`ϵ`$ can be derived from (75) in the limit (64), and using Eq. (79), which implies $`\|\eta _n(ϵ)\|`$ $`=`$ $`\left[\lambda (ϵ)\mathrm{\Xi }_{_{(\mathrm{gs})}}(ϵ)\right]^{2/ϵ}\left[1+{\displaystyle \frac{ϵ}{2}}_n(ϵ)\right]^{2/ϵ}`$ (81) $``$ $`\left[\lambda (ϵ)\mathrm{\Xi }_{_{(\mathrm{gs})}}(ϵ)\right]^{2/ϵ}\mathrm{exp}\left[_n(ϵ)\right].`$ (82) Furthermore, the analysis of the previous section shows that finite results follow only if the coupling constant takes a critical value. For the ground state, this requires that $$\lambda (ϵ)\stackrel{\left(ϵ0\right)}{}\left[\mathrm{\Xi }_{_{(\mathrm{gs})}}(ϵ)\right]^1.$$ (83) Even though Eq. (83) is sufficient for renormalization purposes, let us consider a more general expression $$\lambda (ϵ)=\left[\mathrm{\Xi }_{_{(\mathrm{gs})}}(ϵ)\right]^1\left[1+\frac{ϵ}{2}g(ϵ)\right].$$ (84) Equation (84) defines a residual coupling function $`g(ϵ)`$, which—according to the definition of critical coupling, Eq. (76)—should have the limiting behavior $$ϵg(ϵ)=o(1).$$ (85) As a consequence, from Eqs. (82) and (84), $$\|\eta _n(ϵ)\|=\mathrm{exp}\left[g(ϵ)+_n(ϵ)\right],$$ (86) and $$\left\|\frac{\eta _n(ϵ)}{\eta _{_{(\mathrm{gs})}}(ϵ)}\right\|=\mathrm{exp}\left[_n(ϵ)\right].$$ (87) From the form of Eqs. (86) and (87), it proves convenient to resolve both $`_n(ϵ)`$ and $`g(ϵ)`$ into their various components, i.e., $$_n(ϵ)=_n^{()}(ϵ)+_n^{(0)}+_n^{(+)}(ϵ)$$ (88) and $$g(ϵ)=g^{()}(ϵ)+g^{(0)}+g^{(+)}(ϵ),$$ where $`_n^{()}(ϵ)`$ and $`g^{()}(ϵ)`$ are the divergent pieces; $`_n^{(0)}`$ and $`g^{(0)}`$ are the limits, for $`ϵ=0^+`$, of the finite parts; and $`_n^{(+)}(ϵ),g^{(+)}(ϵ)=o(1)`$. Then Eq. (86) will again give 0 or $`\mathrm{}`$, unless $$[g^{()}(ϵ)+_n^{()}(ϵ)]+[g^{(0)}+_n^{(0)}]+[g^{(+)}(ϵ)+_n^{(+)}(ϵ)]=O(1).$$ (89) From now on, the terms $`g^{(+)}(ϵ)`$ and $`_n^{(+)}(ϵ)`$ can and will be omitted, as they are clearly irrelevant at the level of the renormalized energies. In turn, in Eq. (89), the terms $`[g^{()}(ϵ)+_n^{()}(ϵ)]`$ would give a divergent contribution unless $$g^{()}(ϵ)=_n^{()}(ϵ).$$ (90) Condition (90), in general, cannot be satisfied for all bound states, but it should be satisfied, in particular, for the ground state, so that (from Eqs. (80) and (84)) $$\lambda (ϵ)=\left[\mathrm{\Xi }_{_{(\mathrm{gs})}}(ϵ)\right]^1\left[1+\frac{ϵ}{2}g^{(0)}+o(ϵ)\right]$$ (91) and (from Eqs. (80) and (86)) $$\eta _{_{(\mathrm{gs})}}=e^{g^{(0)}}.$$ (92) Once the ground-state renormalization is established, one can analyze the existence and properties of the excited states. In this regard, Eq. (78) selects only a subset of the states of the regularized theory as bound states of the renormalized theory, once the limit $`ϵ0`$ is taken. More precisely, for any set of quantum numbers $`n`$ for which Eq. (78) is a strict inequality, the coupling becomes weak in the limit $`ϵ=0^+`$, so that the given state is merged with the continuum. Thus, the equality $$\lambda ^{()}=\lambda _n^{()}$$ (93) is a necessary condition for the state labeled by $`n`$ to survive as a bound state. Then, for any state that satisfies Eq. (93), the function $`_n(ϵ)`$ defined in (79) is constrained by the limiting form $$ϵ_n(ϵ)=o(1)$$ (94) and (from Eqs. (74) and (78)) satisfies the inequality $$_n(ϵ)0;$$ (95) in particular, $$_n^{()}(ϵ)0.$$ (96) If the inequality (96) were strict, then Eq. (87) would annihilate the state labeled by $`n`$ in the bound-state sector by exponential suppression; as a consequence, the only alternative option for the state to “survive,” as allowed by the inequality (96), is $$_n^{()}(ϵ)=0,$$ (97) in which case $$\left\|\frac{\eta _n(ϵ)}{\eta _{_{(\mathrm{gs})}}(ϵ)}\right\|=\mathrm{exp}\left[_n^{(0)}\right].$$ (98) It should be pointed out that, when Eqs. (93) and (97) are satisfied, the inequality $`_n^{(0)}<0`$ (from Eq. (95)) guarantees that $`E_n>E_{_{(\mathrm{gs})}}`$. In summary, Eqs. (93), (97), and (98) give the following conditions for the existence of bound states. An excited state labeled with the index $`n(\mathrm{gs})`$ exists if: (i) The critical coupling $`\lambda _n^{()}`$ satisfies the equality (93). (ii) The function $`_n^{()}(ϵ)`$ is identically zero (condition (97)) for the states that already satisfy Eq. (93). (iii) $`_n^{(0)}_{_{(\mathrm{gs})}}^{(0)}=0`$. These are indeed very stringent conditions, so “ordinarily” dimensional transmutation will produce a single bound state, as is the case for the two-dimensional delta-function and inverse square potentials (Ref. ). A final digression on strategy may provide a more direct path in a typical problem. If the energy generating function $`\mathrm{\Xi }_n(ϵ)`$ admits the expansion $$\mathrm{\Xi }_n(ϵ)=\left[L_n(ϵ)\right]^1\left[1+\frac{ϵ}{2}𝒢_n(ϵ)\right],$$ (99) with a power-law leading term $$L_n(ϵ)=a_nϵ^{\tau _n}$$ (100) (where $`a_n`$ and $`\tau _n`$ are constants), and $$ϵ𝒢_n(ϵ)=o(1),$$ (101) then the following results will directly apply. First, the critical coupling will be (from Eqs. (76), (99), and (101)) $$\lambda _n^{()}=\underset{ϵ0}{lim}L_n(ϵ),$$ (102) with a regularized coupling (Eq. (91)) $$\lambda (ϵ)=L_{_{(\mathrm{gs})}}(ϵ)\left\{1+\frac{ϵ}{2}\left[g^{(0)}𝒢_{_{(\mathrm{gs})}}(ϵ)\right]\right\}+o(ϵ),$$ (103) while the function $`_n(ϵ)`$ of Eq. (79) will become $$_n(ϵ)\stackrel{\left(ϵ0\right)}{}𝒢_n(ϵ)𝒢_{_{(\mathrm{gs})}}(ϵ),$$ (104) up to higher order corrections. Thus, the condition for the existence of the ground state will be $`\tau _{_{(\mathrm{gs})}}0`$, while the conditions for the existence of excited states will amount to the existence of an index $`n(\mathrm{gs})`$, such that: (i) $`a_n=a_{_{(\mathrm{gs})}}`$ and $`\tau _n=\tau _{_{(\mathrm{gs})}}`$; (ii) $`𝒢_n^{()}(ϵ)=𝒢_{_{(\mathrm{gs})}}^{()}(ϵ)`$; and (iii) $`𝒢_n^{(0)}𝒢_{_{(\mathrm{gs})}}^{(0)}`$. We now turn to the scattering problem. ### C Renormalized Scattering Sector For the scattering sector, the scattering amplitude $`f_k^{(D)}(\mathrm{\Omega }^{(D)})`$ and the differential scattering cross section $`d\sigma ^{(D)}(k,\mathrm{\Omega }^{(D)})/d\mathrm{\Omega }_D=\|f_k^{(D)}(\mathrm{\Omega }^{(D)})\|^2`$ are functions of the energy $`E=k^2`$ associated with the incident momentum $`k`$, as well as of the hyperspherical angles $`\mathrm{\Omega }^{(D)}`$ (with $`d\mathrm{\Omega }_D`$ being the corresponding element of the $`D`$-dimensional solid angle; see Appendix A). As discussed in Subsection III D and using the language developed in Subsections V A and V B, there are two distinct regimes for scattering. In the weak-coupling regime, $`\lambda <\lambda ^{()}`$, the scattering is well defined for all values of the coupling constant (consistent with the inequality defining the weak-coupling regime). In particular, this scattering is either scale-invariant (energy-independent) or trivial and needs no regularization whatsoever. For example, the inverse square potential gives scale-invariant scattering , when $`\lambda <\left(l+D/21\right)^2`$ (where $`l`$ is the angular momentum quantum number), while the two-dimensional delta potential yields no scattering for $`\lambda <0`$. These examples and issues will be analyzed in greater detail in the second paper in this series. Instead, in the strong-coupling regime, $`\lambda >\lambda ^{()}`$, the coupling constant gets renormalized according to the theory developed in Subsection V B. In particular, this implies that the coupling parameter of the regularized theory is $`ϵ`$-dependent, as displayed in Eq. (91), with a limiting critical value $`\lambda =\lambda ^{()}+0^+`$. Moreover, Eq. (72) is still applicable, as it was derived solely using scaling arguments—but the function $`\mathrm{\Xi }(ϵ)`$ will now have a different specific form, one that is no longer discrete. Then, the scattering matrix and amplitude are determined from the asymptotic form of the scattering wave function, which is an appropriate linear combination of scattering solutions with arguments $`kr=\left(\widehat{\eta }\right)^{1/2}\|\widehat{𝝃}\|=\left[\mathrm{\Xi }(ϵ)\right]^{1/ϵ}\|\widehat{𝝃}\|`$. Equation (72) then implies that the scattering depends upon $$\|\eta (ϵ)\|^{ϵ/2}=1\frac{ϵ}{2}\mathrm{ln}\left(k^2/\mu ^2\right),$$ (105) which, after taking the limit $`ϵ=0^+`$ and using Eq. (91), should provide a breakdown of the scale symmetry through the logarithmic dependence $`\mathrm{ln}\left(k^2/\mu ^2\right)`$. This suggests that the scattering amplitude should be of the form $$f_k^{(D)}(\mathrm{\Omega }^{(D)})=F\mathbf{(}\text{ }k,\mathrm{ln}\left(k^2/\mu ^2\right),\mathrm{\Omega }^{(D)}\mathbf{)}\text{ },$$ (106) where $`F`$ is a dimensional quantity. This procedure will be illustrated for the two-dimensional delta-function potential in Section VI. In Eq. (106), the variable $`k^2/\mu ^2`$ explicitly appears in a characteristic logarithmic form. However, the function $`F`$ is not dimensionless so that its form can be simplified by the $`\mathrm{\Pi }`$ theorem. We now turn to such dimensional considerations. ### D Dimensional Analysis Revisited Let us now rephrase some of the results of Subsections V B and V C in terms of dimensional variables. The dimensional bound-state energies are arranged in a spectrum $$E_n=\mu ^2\eta _n,$$ (107) a conclusion that can be directly drawn from dimensional analysis. In particular, the ground state defines a conventional characteristic scale $$E_{_{(\mathrm{gs})}}=\mu ^2\eta _{_{(\mathrm{gs})}}\mu ^2,$$ (108) where the symbol $``$ refers to the freedom to make a convenient choice of $`g^{(0)}`$, due to its arbitrariness; in this case, we selected $`g^{(0)}=0`$. This point will be further discussed and illustrated in Subsection VI B for the particular case of the two-dimensional delta-function potential. Equation (108) displays in its purest form the emergence of an energy scale from the renormalization procedure; in addition, it shows that naive generalized dimensional analysis (including renormalization parameters according to Eq. (6)) gives straightforwardly the correct result. Of course, Eq. (107) also refers to the excited states, if they exist. Again, by the generalized $`\mathrm{\Pi }`$ theorem, the only remaining information about the spectrum is conveyed by the ratios (cf. Eqs. (98) and (104)) $$\rho _n=\frac{\eta _n}{\eta _{_{(\mathrm{gs})}}}=\mathrm{exp}\left[_n(ϵ)\right]\stackrel{\left(ϵ0\right)}{}\mathrm{exp}\left[𝒢_n(ϵ)𝒢_{_{(\mathrm{gs})}}(ϵ)\right],$$ (109) which give its characteristic “structure,” with the restrictions discussed in Subsection V B. For the scattering sector, as the dimensionality of the cross section is $`(D1)`$ (“area” of hypersuface), it follows that Eq. (106) can be rewritten in the form $$f_k^{(D)}(\mathrm{\Omega }^{(D)})=k^{(D1)/2}\mathrm{\Pi }\mathbf{(}\text{ }\mathrm{ln}\left(k^2/\mu ^2\right),\mathrm{\Omega }^{(D)}\mathbf{)}\text{ },$$ (110) where $`\mathrm{\Pi }(u,\mathrm{\Omega }^{(D)})`$ is a dimensionless function of the dimensionless ratio $`u=(k/\mu )^2`$. Equation (110) will be valid, whether the system is capable of producing bound states or not; in the weak-coupling regime, the function $`\mathrm{\Pi }`$ is identically constant. On the other hand, if there is at least one bound state, the existence of a characteristic energy scale $`E_{_{(\mathrm{gs})}}`$, Eq. (108), yields an alternative form of (110), $$f_k^{(D)}(\mathrm{\Omega }^{(D)})=k^{(D1)/2}\stackrel{ˇ}{\mathrm{\Pi }}\mathbf{(}\text{ }\mathrm{ln}\left(k^2/E_{_{(\mathrm{gs})}}\right),\mathrm{\Omega }^{(D)}\mathbf{)}\text{ },$$ (111) where $`\stackrel{ˇ}{\mathrm{\Pi }}\mathbf{(}\text{ }\mathrm{ln}\left(E/E_{_{(\mathrm{gs})}}\right),\mathrm{\Omega }^{(D)}\mathbf{)}\text{ }=\mathrm{\Pi }\mathbf{(}\text{ }\mathrm{ln}\left(E/\mu ^2\right),\mathrm{\Omega }^{(D)}\mathbf{)}\text{ }`$ is another dimensionless function. In fact, when the assignment $`g^{(0)}=0`$ is made (Eq. (108)), the simple identity $`\stackrel{ˇ}{\mathrm{\Pi }}=\mathrm{\Pi }`$ takes place. Equations (107), (110), and (111) are just a consequence of generalized dimensional analysis. ## VI TWO-DIMENSIONAL DELTA-FUNCTION POTENTIAL One of the most basic properties of a quantum field theory is locality, which leads to a nonrelativistic limit with highly singular potentials of zero range, also known as pseudopotentials . The simplest pseudopotentials are plain delta functions, which display a large number of unusual features; however, in this section we will only explore those properties that relate to the dimensional transmutation produced by the two-dimensional representative of this class. As we will see, this potential displays all the characteristic fingerprints of dimensional transmutation that we discussed in previous sections. In fact, the two-dimensional delta-function potential has been extensively studied in the literature, mainly using cutoff regularization and square-well regularization . In our approach, we will exclusively use dimensional regularization within the framework defined in Sections IV and V. Our strategy will be to compare the calculations with the predictions and requirements of the general theory of Sections IV and V. However, we will use the dimensional form (63) of the Schrödinger equation from scratch, rather than any of the dimensionless equations (65) or (71). In effect, the dimensionless counterparts are most useful at a theoretical level, in establishing the relationships between all relevant parameters for our problem; yet, in practice, it is more direct to set up the “ordinary” (dimensional) dimensionally regularized Schrödinger equation. ### A Dimensional Regularization of the Two-Dimensional Delta-Function Potential The two-dimensional delta-function potential is a particular zero-range or contact interaction of the form $$V(𝐫)=\lambda \delta ^{(2)}(𝐫).$$ (112) We have already seen that this potential is scale-invariant. Using the techniques of Section IV, we now apply the dimensional-continuation prescription of Eq. (51) to $`W^{(2)}(𝐫)=\delta ^{(2)}(𝐫)`$, with $`D_0=2`$, i.e., $`W^{(D)}(𝐫_D)`$ $`=`$ $`{\displaystyle \frac{d^Dk_D}{(2\pi )^D}e^{i𝐤_D𝐫_D}\left[d^2r_2e^{i𝐤_2𝐫_2}\delta ^{(2)}(𝐫_2)\right]_{𝐤_2𝐤_D}},`$ (113) $`=`$ $`{\displaystyle \frac{d^Dk_D}{(2\pi )^D}e^{i𝐤_D𝐫_D}\left[1\right]_{𝐤_2𝐤_D}}`$ (114) $`=`$ $`{\displaystyle \frac{d^Dk_D}{(2\pi )^D}e^{i𝐤_D𝐫_D}}=\delta ^{(D)}(𝐫),`$ (115) which is the obvious dimensional extension of the original delta-function potential. Thus, in what follows, we will consider the dimensionally regularized problem $$[_{𝐫,D}^2\lambda \mu ^ϵ\delta ^{(D)}(𝐫)]\mathrm{\Psi }(𝐫)=E\mathrm{\Psi }(𝐫).$$ (116) Straightforward solution of Eq. (116) in this context should not be interpreted as a way of drawing conclusions about the $`D`$-dimensional delta-function potential. Instead, it is just the means to regularize the $`D_0=2`$ case. Of course, one could adjust the regularization to be applied around a value $`D_02`$; however, we will not attempt such modification in this paper, as it is not directly relevant to dimensional transmutation. ### B Bound-State Sector for a Two-Dimensional Delta-Function Potential Equation (116) can be easily solved in momentum space; for the bound-state sector, $$\stackrel{~}{\mathrm{\Psi }}(𝐪)=\lambda \mu ^ϵ\frac{\mathrm{\Psi }(\mathrm{𝟎})}{q^2E},$$ (117) which, via the inverse Fourier transform, yields the position-space eigenfunctions. However, if we are only interested in the eigenvalue equation, it suffices to consider the value of the wave function at the origin, $$\mathrm{\Psi }(\mathrm{𝟎})=\frac{d^Dq}{(2\pi )^D}\stackrel{~}{\mathrm{\Psi }}(𝐪),$$ (118) so that Eq. (117) gives $$\frac{\lambda \mu ^ϵ}{(2\pi )^D}\frac{d^Dq}{q^2+\|E\|}=1,$$ (119) where $`E=\|E\|`$ ($`E<0`$ for the possible bound states). Equation (119) can be straightforwardly integrated using Eq. (A19), which implies that $$\frac{d^Dq}{q^2+\|E\|}=\pi ^{D/2}\|E\|^{D/21}\mathrm{\Gamma }\left(1\frac{D}{2}\right),$$ (120) and the regularized eigenvalue equation takes the form $$\frac{\lambda \mu ^ϵ}{4\pi }\left(\frac{\|E\|}{4\pi }\right)^{D/21}\mathrm{\Gamma }\left(1\frac{D}{2}\right)=1.$$ (121) It is a simple exercise to verify that Eq. (121) reduces to the familiar textbook result $`\kappa =\sqrt{\|E\|}=\lambda /2`$ for $`D=1`$ . On the other hand, the left-hand side is divergent for $`D=`$ $`2,4,6,\mathrm{}`$. However, the restriction on the spatial dimension $`D`$ of regular potentials is even stronger because more stringent conditions are dictated by the eigenfunctions, as we will see next. In our subsequent analysis, both for two-dimensional delta-function and inverse square potentials, the dimension $`D`$ will usually appear in terms of the variable $$\nu =D/21,$$ (122) which will thereby simplify the form of most formulas; for example, the eigenvalue Eq. (121) reads (with $`ϵ=2\nu `$) $$\frac{\lambda \mu ^{2\nu }}{4\pi }\left(\frac{\|E\|}{4\pi }\right)^\nu \mathrm{\Gamma }\left(\nu \right)=1.$$ (123) Then the inverse Fourier transform, $$\mathrm{\Psi }(𝐫)=\lambda \mu ^ϵ\mathrm{\Psi }(\mathrm{𝟎})\frac{e^{i𝐪𝐫}}{q^2+\|E\|}\frac{d^Dq}{(2\pi )^D},$$ (124) is recognized to be proportional to the Green’s function $`𝒦_D(𝐫;\kappa )`$ for the modified Helmholtz equation (see Appendix B, in particular Eqs. (B1)–(B4)), $`\mathrm{\Psi }(𝐫)`$ $`=`$ $`\lambda \mu ^ϵ\mathrm{\Psi }(\mathrm{𝟎})𝒦_D(𝐫;\kappa )`$ (125) $`=`$ $`{\displaystyle \frac{\lambda \mu ^ϵ\mathrm{\Psi }(\mathrm{𝟎})}{2\pi }}\left({\displaystyle \frac{\kappa }{2\pi r}}\right)^\nu K_\nu (\kappa r),`$ (126) where $`\kappa =\sqrt{\|E\|}`$, $`r=\|𝐫\|`$, and $`K_\nu (z)`$ is the modified Bessel function of the second kind of order $`\nu `$. The asymptotic behavior of the wave function $`\mathrm{\Psi }(𝐫)`$ (and of the Green’s function $`𝒦_D(𝐫;\kappa )`$) is governed by that of $`K_\nu (z)`$ , $$K_\nu (z)\stackrel{\left(z\mathrm{}\right)}{}\sqrt{\frac{\pi }{2z}}e^z\left[1+O(1/z)\right],$$ (127) whence $$\mathrm{\Psi }(𝐫)\stackrel{\left(r\mathrm{}\right)}{}\frac{\lambda \mu ^ϵ\mathrm{\Psi }(\mathrm{𝟎})}{4\pi }\left(\frac{\kappa }{2\pi }\right)^{(D3)/2}\frac{e^{\kappa r}}{r^{(D1)/2}}\left[1+O(1/r)\right],$$ (128) which displays the correct behavior for a bound-state wave function at infinity. However, near the origin, the modified Bessel function has a singular behavior of the form $`K_p(z)`$ $`\stackrel{(z0)}{}`$ $`{\displaystyle \frac{1}{2}}\left[\mathrm{\Gamma }(p)\left({\displaystyle \frac{z}{2}}\right)^p+\mathrm{\Gamma }(p)\left({\displaystyle \frac{z}{2}}\right)^p\right]\left[1+O(z^2)\right],`$ (129) $`\stackrel{(z0)}{}`$ $`\{\begin{array}{cc}\frac{1}{2}\mathrm{\Gamma }(\|p\|)(2/z)^{\|p\|}\left[1+O(z^2)\right]\hfill & \mathrm{for}p0\hfill \\ \left[\mathrm{ln}(z/2)+\gamma \right]\left[1+O(p,z^2)\right]\hfill & \mathrm{for}p0\hfill \end{array},`$ (132) where $`\gamma `$ is the Euler-Mascheroni constant; then, the explicit form of the wave function is $$\mathrm{\Psi }(𝐫)\stackrel{\left(r0\right)}{}\frac{\lambda \mu ^ϵ\mathrm{\Psi }(\mathrm{𝟎})}{4\pi }\left(\frac{\kappa ^2}{4\pi }\right)^\nu \times \{\begin{array}{cc}\mathrm{\Gamma }(\nu )(\kappa r/2)^{2\nu }\hfill & \mathrm{for}D>2\hfill \\ 2\left[\mathrm{ln}(\kappa r/2)+\gamma \right]\hfill & \mathrm{for}D=2\hfill \\ \mathrm{\Gamma }(\nu )\hfill & \mathrm{for}\mathrm{\hspace{0.33em}0}<D<2\hfill \end{array},$$ (133) which shows that the nature of the solution changes around $`\nu =0`$, i.e., for $`D=2`$. This confirms the critical character of the dimension $`D=2`$ for the delta-function potential. Notice that the wave functions are singular at the origin for any dimension $`D2`$. Parenthetically, this is an example of an ultraviolet divergence: the wave function is singular at small distances or due to large momenta (cf. Eq. (119)). For $`D<2`$, we can regularize the two-dimensional delta-function potential and take the limit $`\nu 0`$ in Eq. (133), thus recovering self-consistently the eigenvalue equation (123), which, with $`D=2ϵ`$, i.e., $`\nu =ϵ/2`$, reads $$\frac{\lambda \mu ^ϵ}{4\pi }\left(\frac{\|E\|}{4\pi }\right)^{ϵ/2}\mathrm{\Gamma }\left(\frac{ϵ}{2}\right)=1,$$ (134) Alternatively, in the language of Eq. (126), this eigenvalue equation can be simply enforced by the condition $$\lambda \mu ^ϵ𝒦_D(\mathrm{𝟎};\kappa )=1.$$ (135) Having completed the exploratory analysis of the bound-state sector and found the eigenvalue equation ((121), (123), (134), or (135)), we are ready to compare these expressions with the general eigenvalue equation (75), which will now include an energy generating function $$\mathrm{\Xi }(ϵ)=\frac{1}{4\pi }\left(4\pi \right)^{ϵ/2}\mathrm{\Gamma }\left(\frac{ϵ}{2}\right).$$ (136) It is immediately apparent from Eq. (136) that the theory has only one bound state, so that there is no need for a quantum number. In order to determine whether this ground state survives the renormalization process, we should look at the $`ϵ=0^+`$ limit of Eq. (136), which can be conveniently examined through the expansion $$\mathrm{\Xi }(ϵ)=\frac{1}{2\pi ϵ}\left[1+\frac{ϵ}{2}\left(\mathrm{ln}4\pi \gamma \right)+O(ϵ^2)\right].$$ (137) From Eqs. (76) and (137), the critical coupling is found to be $$\lambda ^{()}=0,$$ (138) so that the theory looks asymptotically free but still engenders a unique bound state. It should be noticed that this is achieved by the regularization of the coupling constant through the strategy of Eq. (91) (or Eq. (103)), so that $$\lambda (ϵ)=2\pi ϵ\left\{1+\frac{ϵ}{2}\left[g^{(0)}\left(\mathrm{ln}4\pi \gamma \right)\right]\right\},$$ (139) leading to a ground state $$E_{_{(\mathrm{gs})}}=\mu ^2e^{g^{(0)}}.$$ (140) A final remark about renormalization shows additional parallels with the corresponding field-theory problems. As usual, the arbitrariness in the choice of the finite part $`g^{(0)}`$ can be used to simplify the expressions above in such a way that $`\|E_{_{(\mathrm{gs})}}\|=\mu ^2`$, as displayed in Eq. (108). On the other hand, the singular nature of the ground state has been tamed by subtracting the divergent part of Eq. (136), which amounts to the subtraction of the pole $`1/2\pi ϵ`$. However, due to the arbitrariness in the choice of $`g^{(0)}`$, at the level of the ground-state energy, we have also subtracted—along with the pole—the term $`\mathrm{ln}4\pi \gamma `$ (which is an artifact of the dimensional-regularization technique). This is recognized to be the usual modified minimal subtraction ($`\overline{\mathrm{MS}}`$) scheme . In conclusion, the unregularized problem has a singular spectrum with a unique energy level at $`\mathrm{}`$ and vanishing critical coupling. The regularization process brings this level to a finite value, which, upon renormalization, becomes the unique ground state of the two-dimensional delta-function potential. ### C Scattering Sector for a Two-Dimensional Delta-Function Potential The scattering sector of the Schrödinger equation is described by its equivalent Lippmann-Schwinger equation (C3), which, for a two-dimensional delta-function potential (116), takes the simple form $$\mathrm{\Psi }^{(+)}(𝐫)=e^{i𝐤𝐫}\lambda \mu ^ϵ𝒢_D^{(+)}(𝐫;k)\mathrm{\Psi }^{(+)}(\mathrm{𝟎}),$$ (141) where $`𝒢_D^{(+)}(𝐫;k)`$ is the Green’s function for the Helmholtz equation, with outgoing boundary conditions (see Appendix B, in particular Eqs.(B5)–(B10)). In particular, Eq. (141) implies that $$\mathrm{\Psi }^{(+)}(\mathrm{𝟎})=\left[1+\lambda \mu ^ϵ𝒢_D^{(+)}(\mathrm{𝟎};k)\right]^1.$$ (142) The asymptotic form of Eq. (141) is obtained as described in Appendix C, according to which (e.g., Eqs. (C5) and (C8)) the scattering amplitude becomes $$f_k^{(D)}(\mathrm{\Omega }^{(D)})=\mathrm{\Gamma }_D(k)\lambda \mu ^ϵ\mathrm{\Psi }^{(+)}(\mathrm{𝟎}),$$ (143) where $$\mathrm{\Gamma }_D(k)=\frac{1}{4\pi }\left(\frac{k}{2\pi }\right)^{(D3)/2}.$$ (144) Finally, Eqs. (142) and (143) provide the desired expression, $$f_k^{(D)}(\mathrm{\Omega }^{(D)})=\mathrm{\Gamma }_D(k)\left[\left(\lambda \mu ^ϵ\right)^1+𝒢_D^{(+)}(\mathrm{𝟎};k)\right]^1.$$ (145) Equation (145) is singular for $`D2`$, as can be seen from the divergent small-argument limit of Eq. (B10). However, one can use the renormalization of the bound-state sector to eliminate this divergence through the regularized coupling, Eq. (139). More precisely, for an attractive potential, we found that the coupling constant can be traded in favor of dimensional parameters, e.g., using the Green’s function $`𝒦_D(𝐫;\kappa )`$ for the bound-state sector in Eq. (135). In other words, using the renormalization for the bound-state sector, we will now obtain directly the renormalized scattering amplitude, which is explicitly given by the limit $$f_k(\mathrm{\Omega }^{(D)})=\mathrm{\Gamma }_D(k)\underset{r0}{lim}\left[𝒦_D(𝐫;\sqrt{E_{_{(\mathrm{gs})}}})𝒢_D^{(+)}(𝐫;k)\right]^1,$$ (146) where the replacement $$\kappa =\sqrt{\|E_{_{(\mathrm{gs})}}\|}$$ (147) was made. Equation (146) already displays a remarkable property of the scattering by a delta-function potential: it is isotropic, i.e., it only scatters s-waves, as it corresponds intuitively to a contact interaction. Let us now evaluate the limit in Eq. (146). First, from Eqs. (120) and (B3), $$\underset{r0}{lim}𝒦_D(𝐫;\sqrt{\|E_{_{(\mathrm{gs})}}\|})=\frac{1}{(2\pi )^D}\frac{d^Dq}{q^2+\|E_{_{(\mathrm{gs})}}\|}=\frac{1}{(4\pi )^{D/2}}\|E_{_{(\mathrm{gs})}}\|^{D/21}\mathrm{\Gamma }\left(1\frac{D}{2}\right).$$ (148) Next, $`lim_{r0}𝒢_D^{(+)}(𝐫;k)`$ can be obtained by analytic continuation $$𝒢_D^{(+)}(𝐫;k)=𝒦_D(𝐫;\kappa )\|_{\kappa ^2(k^2+i\delta )},$$ (149) where $`\delta =0^+`$, which implies that $$\underset{r0}{lim}𝒢_D^{(+)}(𝐫;k)=\frac{1}{(2\pi )^D}\frac{d^Dq}{k^2q^2+i\delta }=\frac{1}{(4\pi )^{D/2}}\left(k^2i\delta \right)^{D/21}\mathrm{\Gamma }\left(1\frac{D}{2}\right).$$ (150) Thus, $$\underset{r0}{lim}\left[𝒢_D^{(+)}(𝐫;k)𝒦_D(𝐫;\kappa )\right]=\frac{1}{(4\pi )^{D/2}}\mathrm{\Gamma }\left(1\frac{D}{2}\right)\left[\|E_{_{(\mathrm{gs})}}\|^{D/21}\left(k^2i\delta \right)^{D/21}\right],$$ (151) which can be evaluated in the limit $`ϵ0^+`$, with $`D=2ϵ`$, $$\underset{ϵ0}{lim}\underset{r0}{lim}\left[𝒢_D^{(+)}(𝐫;k)𝒦_D(𝐫;\kappa )\right]=\frac{1}{4\pi }\left(\mathrm{ln}\|E_{_{(\mathrm{gs})}}\|\mathrm{ln}k^2+i\pi \right),$$ (152) where the identity $`\mathrm{ln}[(k^2+i\delta )]=\mathrm{ln}k^2i\pi `$ was used. Finally, the scattering amplitude is obtained by replacing Eqs. (144) (with $`D=2`$) and (152) in (146), i.e., $$f_k^{(2)}(\mathrm{\Omega }^{(2)})=\sqrt{\frac{2\pi }{k}}\left[\mathrm{ln}\left(\frac{k^2}{E_{_{(\mathrm{gs})}}}\right)i\pi \right]^1.$$ (153) Equation (153) is seen to agree with the prediction of generalized dimensional analysis, Eqs. (110)–(111), with a dimensionless variable $`\mathrm{\Pi }(u)=\sqrt{2\pi }\left[\mathrm{ln}ui\pi \right]^1`$, and $`u=k^2/E_{_{(\mathrm{gs})}}`$. Finally, the differential scattering cross section $`d\sigma ^{(2)}(k,\mathrm{\Omega }^{(2)})/d\mathrm{\Omega }_2`$, from Eq. (C9), again agrees with the prediction of generalized dimensional analysis, Eqs. (110)–(111), providing a dimensionless variable $`\mathrm{\Pi }(u)=2\pi \left[(\mathrm{ln}u)^2+\pi ^2\right]^1`$, for the energy ratio $`u=k^2/E_{_{(\mathrm{gs})}}`$. ## VII CONCLUSIONS Until recently, it had been generally assumed that generic field-theoretic tools and concepts are useful only for systems with infinitely many degrees of freedom. While this perception is essentially correct for “regular” systems, it is now recognized, as discussed in our series of papers, that such techniques can be generalized and used to extract meaningful physical results for certain “singular” systems with a finite number of degrees of freedom. In this paper, we developed systematic uses of the techniques of dimensional regularization and renormalization, and of the concept of dimensional transmutation, with the purpose of gathering information about the class of scale-invariant potentials. Our discussion relied on dimensional regularization, which we argued provides a generic tool for the treatment of all members of that class, by establishing a simple link between the two meanings of the word dimension. From our fairly general investigation, we have learned that all scale-invariant potentials are homogeneous of degree $`2`$ and share a number of remarkable properties; here, without attempting to give an exhaustive list, we summarize a few of the most outstanding: (i) There exists a critical coupling $`\lambda ^{()}`$ such that, for $`\lambda <\lambda ^{()}`$ (weak coupling) the Hamiltonian is self-adjoint but produces no bound states, while for $`\lambda >\lambda ^{()}`$ (strong coupling) it loses its self-adjoint character, generating a continuum of bound states not bounded from below and requiring renormalization. (ii) Solution of the regularized theory for strong coupling yields a master eigenvalue equation (72), which condenses all the information about the given scale-invariant potential and requires proper renormalization. (iii) The ground state of a given “strong” scale-invariant potential exists provided that $`\lambda ^{()}=\left[lim_{ϵ0}\mathrm{\Xi }_n(ϵ)\right]^1`$ (Eq. (76)) be finite. (iv) Excited states exist under the demanding conditions listed after Eq. (98). Thus, “strong” scale-invariant potentials have a tendency to suppress excited states. (v) The scattering sector remains scale-invariant or trivial in the weak-coupling regime, while it displays a logarithmic dependence $`\mathrm{ln}\left(k^2/\mu ^2\right)`$, with respect to the energy $`k^2`$, in the strong-coupling regime. (vi) In short, in the strong-coupling regime, a given scale-invariant potential leads to dimensional transmutation, which manifests itself on the existence of at least one bound state and a scale-dependent scattering matrix. The dimensional transmutation exhibited for strong coupling amounts to the emergence of a scale anomaly, i.e., quantum-mechanical breaking of the classical scale symmetry. Additional progress in understanding these singular quantum-mechanical systems can best be achieved by studying specific cases. A first attempt was made in this paper by considering some aspects of the two-dimensional delta-function potential. In that regard, the second paper in this series will present a more thorough analysis of the two-dimensional delta-function potential, as well as a novel treatment of the anomalous transmuting behavior of the inverse square potential. ## A DIMENSIONAL REGULARIZATION IN $`D`$-DIMENSIONAL EUCLIDEAN SPACES Just like for the corresponding case in quantum field theory , our approach is based on the dimensional extension of mathematical expressions for a system that is assumed to be embedded in a $`D`$-dimensional Euclidean space. Then, starting from the Cartesian coordinates $`(x_1,\mathrm{},x_D)`$, one can introduce an alternative set of $`D`$-dimensional hyperspherical polar coordinates $`(q_0=r,q_1=\theta _1,\mathrm{},q_{D1}=\theta _{D1})`$ via the transformation equations $`x_1`$ $`=`$ $`r\mathrm{cos}\theta _1`$ (A1) $`x_2`$ $`=`$ $`r\mathrm{sin}\theta _1\mathrm{cos}\theta _2`$ (A2) $`\mathrm{}`$ (A3) $`x_j`$ $`=`$ $`r\left({\displaystyle \underset{k=1}{\overset{j1}{}}}\mathrm{sin}\theta _k\right)\mathrm{cos}\theta _j`$ (A4) $`\mathrm{}`$ (A5) $`x_D`$ $`=`$ $`r{\displaystyle \underset{k=1}{\overset{D1}{}}}\mathrm{sin}\theta _k,`$ (A6) where the ranges for the hyperspherical polar variables are $`0\theta _j\pi `$ for $`j=1,\mathrm{},D2`$; $`0\varphi \theta _{D1}2\pi `$; and $`0r<\mathrm{}`$. All the basic geometric quantities associated with hyperspherical coordinates can be constructed through the scale coefficients $`h_j`$ for the diagonal metric $`(g_{ij})=\mathrm{diag}(h_j^2)_{0jD1}`$ ; these coefficients are straightforwardly given by $`h_0=1`$, $`h_1=r`$, and $`h_j=r_{k=1}^{j1}\mathrm{sin}\theta _k`$ (for $`2jD1`$), while their product is $$h^{(D)}=\underset{j=0}{\overset{D1}{}}h_j=r^{D1}\underset{j=1}{\overset{D1}{}}\mathrm{sin}^{Dj1}\theta _j.$$ (A7) In our series of papers, both the Laplacian operator and the element of volume are needed. The Laplacian can be computed from $`_D^2`$ $`=`$ $`{\displaystyle \frac{1}{h^{(D)}}}{\displaystyle \underset{j=0}{\overset{D1}{}}}{\displaystyle \frac{}{q_j}}\left({\displaystyle \frac{h^{(D)}}{h_j^2}}{\displaystyle \frac{}{q_j}}\right)`$ (A8) $`=`$ $`\mathrm{\Delta }_r^{(D)}+{\displaystyle \frac{1}{r^2}}\mathrm{\Delta }_{\mathrm{\Omega }^{(D)}},`$ (A9) where its radial part is, explicitly, $$\mathrm{\Delta }_r^{(D)}=\frac{1}{r^{D1}}\frac{}{r}\left(r^{D1}\frac{}{r}\right),$$ (A10) while its angular part is $$\mathrm{\Delta }_{\mathrm{\Omega }^{(D)}}=\underset{j=1}{\overset{D1}{}}\left[\left(\underset{k=1}{\overset{j1}{}}\mathrm{sin}^2\theta _k\right)\mathrm{sin}^{Dj1}\theta _j\right]^1\frac{}{\theta _j}\left(\mathrm{sin}^{Dj1}\theta _j\frac{}{\theta _j}\right),$$ (A11) in which the notation $`\mathrm{\Omega }^{(D)}(\theta _1,\mathrm{},\theta _{D1})`$ has been introduced and it is implied that $`_{k=k_1}^{k_2}1`$ when $`k_1>k_2`$ (i.e., for $`j=1`$). Similarly, the element of the $`D`$-dimensional solid angle becomes $$d\mathrm{\Omega }_D=\underset{j=1}{\overset{D1}{}}h_jd\theta _j=h^{(D)}\underset{j=1}{\overset{D1}{}}d\theta _j=\underset{j=1}{\overset{D1}{}}\mathrm{sin}^{Dj1}\theta _jd\theta _j,$$ (A12) in terms of which the $`D`$-dimensional volume element is given by $$d^Dr=r^{D1}d\mathrm{\Omega }_Ddr.$$ (A13) Equation (A12) can be integrated to a total $`D`$-dimensional solid angle $$\mathrm{\Omega }_D=𝑑\mathrm{\Omega }_D=\left(\underset{j=1}{\overset{D2}{}}_0^\pi 𝑑\theta _j\mathrm{sin}^{Dj1}\theta _j\right)_0^{2\pi }𝑑\theta _{D1},$$ (A14) where the angular integrals can be evaluated using the beta-function identity , $$_0^{\pi /2}\mathrm{sin}^m\theta d\theta =\frac{1}{2}B\mathbf{(}\text{ }(m+1)/2,1/2\mathbf{)}\text{ }=\frac{\sqrt{\pi }\mathrm{\Gamma }\mathbf{(}\text{ }(m+1)/2\mathbf{)}\text{ }}{2\mathrm{\Gamma }\mathbf{(}\text{ }(m+2)/2\mathbf{)}\text{ }},$$ (A15) whence $$\mathrm{\Omega }_D=\frac{2\pi ^{D/2}}{\mathrm{\Gamma }(D/2)}.$$ (A16) With the given element of volume, Eq. (A13), one can compute the integral of any function $`f(r)`$ that exhibits $`D`$-dimensional central symmetry, that is, $$f(r)d^Dr=\frac{2\pi ^{D/2}}{\mathrm{\Gamma }(D/2)}_0^{\mathrm{}}r^{D1}f(r)𝑑r.$$ (A17) In particular, Eqs. (A12)–(A17) are essential for the evaluation of expressions in dimensional regularization, in conjunction with another beta-function identity , $$_0^{\mathrm{}}\frac{x^{2\alpha 1}}{(x^2+1)^{\alpha +\beta }}𝑑x=\frac{1}{2}B(\alpha ,\beta ),$$ (A18) whence $$\frac{(q^2)^n}{(q^2+a^2)^m}d^Dq=\pi ^{D/2}a^{D+2n2m}\frac{\mathrm{\Gamma }(n+D/2)\mathrm{\Gamma }(mnD/2)}{\mathrm{\Gamma }(D/2)\mathrm{\Gamma }(m)}.$$ (A19) Finally, let us consider the general $`D`$-dimensional Fourier transform $`\stackrel{~}{f}(𝐬)`$ of a function $`f(𝐮)`$, defined by $$\stackrel{~}{f}(𝐬)=\frac{n_D}{(2\pi )^{D/2}}d^Due^{i𝐬𝐮}f(𝐮),$$ (A20) with an arbitrary normalization constant $`n_D`$ (usually $`n_D=1`$ or $`n_D=(2\pi )^{\pm D/2}`$). Its computation can be simplified considerably for the particular case when the function displays central symmetry, i.e., $`f(𝐮)=f(u)`$. In effect, for the integration of $`f(u)`$, selecting coordinates according to $`𝐬𝐮=su\mathrm{cos}\theta _1`$, it follows that $`\stackrel{~}{f}(𝐬)`$ $`=`$ $`{\displaystyle \frac{n_D}{(2\pi )^{D/2}}}\left({\displaystyle \underset{j=2}{\overset{D2}{}}}{\displaystyle _0^\pi }𝑑\theta _j\mathrm{sin}^{Dj1}\theta _j\right){\displaystyle _0^{2\pi }}𝑑\theta _{D1}`$ (A21) $`\times `$ $`{\displaystyle _0^{\mathrm{}}}𝑑uu^{D1}f(u){\displaystyle _0^\pi }𝑑\theta _1\mathrm{sin}^{D2}\theta _1e^{isu\mathrm{cos}\theta _1}`$ (A22) $`=`$ $`n_D(2\pi )^{D/2}\mathrm{\Omega }_{D1}{\displaystyle _0^{\mathrm{}}}𝑑uu^{D1}f(u)(D/21,su),`$ (A23) where $$(\nu ,z)=_0^\pi e^{iz\mathrm{cos}\theta }sin^{2\nu }\theta 𝑑\theta =\frac{\mathrm{\Gamma }(\nu +1/2)\mathrm{\Gamma }(1/2)}{(z/2)^\nu }J_\nu (z),$$ (A24) which implies that ($`\nu =D/21`$) $$\stackrel{~}{f}(𝐬)=\frac{n_D}{(2\pi )^{D/2}}d^Due^{i𝐬𝐮}f(𝐮)=\frac{n_D}{s^{D/21}}_0^{\mathrm{}}f(u)J_{D/21}(su)u^{D/2}𝑑u.$$ (A25) Equation (A25), which is a Hankel transform, is sometimes referred to as Bochner’s theorem. ## B $`D`$-DIMENSIONAL GREEN’S FUNCTIONS As an example of Bochner’s theorem, we will now compute the infinite-space Green’s function for the $`D`$-dimensional Helmholtz equation. We will start with the modified Helmholtz equation, $$\left[_{𝐫,D}^2\kappa ^2\right]K_D(𝐫,𝐫^{};\kappa )=\delta ^{(D)}(𝐫𝐫^{}),$$ (B1) whose Green’s function $`K_D(𝐫,𝐫^{};\kappa )`$ can be computed by applying translational invariance, i.e., $`K_D(𝐫,𝐫^{};\kappa )=𝒦_D(𝐑;\kappa )`$, with $`𝐑=𝐫𝐫^{}`$. Its Fourier transform $`\stackrel{~}{𝒦}_D(𝐪;\kappa )=(q^2+\kappa ^2)^1`$ leads to an integral expression $`𝒦_D(𝐑;\kappa )`$ $`=`$ $`{\displaystyle \frac{d^Dq}{(2\pi )^D}\frac{e^{i𝐪𝐑}}{q^2+\kappa ^2}}`$ (B2) $`=`$ $`(2\pi )^{D/2}R^{(D/21)}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{q^{D/2}J_{D/21}(qR)}{q^2+\kappa ^2}}𝑑q,`$ (B3) in which Eq. (A25) was used. Equation (B3) can be explicitly evaluated in terms of the modified Bessel function of the second kind $`K_\nu (\kappa R)`$, of order $`\nu =D/21`$, i.e. , $$𝒦_D(𝐑;\kappa )=\frac{1}{2\pi }\left(\frac{\kappa }{2\pi R}\right)^\nu K_\nu (\kappa R),$$ (B4) where the dimensional variable $`\nu =D/21`$ (cf. Eq. (122)) has been explicitly introduced. Likewise, for the ordinary Helmholtz equation, $$\left[_{𝐫,D}^2+k^2\right]G_D(𝐫,𝐫^{};k)=\delta ^{(D)}(𝐫𝐫^{})$$ (B5) in infinite space, translational invariance implies that $`G_D(𝐫,𝐫^{};k)=𝒢_D(𝐑;k)`$, with $`𝐑=𝐫𝐫^{}`$. However, its Fourier transform $`\stackrel{~}{𝒢}_D(𝐪;k)=(k^2q^2)^1`$ leads to an ill-defined integral expression that needs to be evaluated by a prescription defining the boundary conditions at infinity; for outgoing $`(+)`$ and incoming $`()`$ boundary conditions, $`𝒢_D^{(\pm )}(𝐑;k)`$ $`=`$ $`{\displaystyle \frac{d^Dq}{(2\pi )^D}\frac{e^{i𝐪𝐑}}{k^2q^2\pm i\delta }}`$ (B6) $`=`$ $`(2\pi )^{D/2}R^{(D/21)}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{q^{D/2}J_{D/21}(qR)}{k^2q^2\pm i\delta }}𝑑q,`$ (B7) where Eq. (A25) was used and $`\delta =0^+`$. Equation (B7) can be explicitly evaluated in terms of Hankel functions of order $`\nu =D/21`$; in fact, it is easy to see that Eq. (B5) can be obtained from (B1) with the replacement $`\kappa =ik`$, so that $$𝒢_D^{(\pm )}(𝐑;k)=𝒦_D(𝐑;\kappa =ik),$$ (B8) and the choice of signs amounts to the choice of boundary conditions at infinity or the $`i\delta `$ prescription. From the identity $$K_\nu (iz)=\pm \frac{\pi i}{2}e^{\pm i\pi \nu /2}H_\nu ^{(1,2)}(z),$$ (B9) Eq. (B8) is converted into $$𝒢_D^{(\pm )}(𝐑;k)=\frac{i}{4}\left(\frac{k}{2\pi R}\right)^\nu H_\nu ^{(1,2)}(kR).$$ (B10) Equations (B4) and (B10) are well known and reduce to the familiar results in one, two, and three dimensions . ## C SCATTERING IN $`D`$ DIMENSIONS Just as in the standard 3-D scattering formalism, the $`D`$-dimensional time-independent operator Schrödinger equation $$\left(H_0E\right)\|\mathrm{\Psi }>=V\|\mathrm{\Psi }>$$ (C1) is equivalent to a Lippmann-Schwinger equation $$\|\mathrm{\Psi }^{(+)}>=\|\chi >+\left(EH_0+i\delta \right)^1V\|\mathrm{\Psi }^{(+)}>$$ (C2) in which the state vector $`\|\mathrm{\Psi }^{(+)}>`$ is explicitly resolved into an incident wave $`\|\chi >`$ (solution of the free-particle case) and a second term that represents the outgoing scattered wave (with an appropriate boundary condition summarized by the $`i\delta =i0^+`$ prescription). In what follows, we will assume that $`\|\chi >=\|\chi _𝐤>`$ represents a $`D`$-dimensional plane wave $`e^{i𝐤𝐫}`$. Equation (C2) can be converted into the integral form $$\mathrm{\Psi }^{(+)}(𝐫)=e^{i𝐤𝐫}+d^Dr^{}G_D^{(+)}(𝐫,𝐫^{};k)V(𝐫^{})\mathrm{\Psi }^{(+)}(𝐫^{}),$$ (C3) by the introduction of one of the Green’s functions computed in Appendix A, namely, $`G_D^{(+)}(𝐫,𝐫^{};k)=<𝐫\|\left(EH_0+i\delta \right)^1\|𝐫^{}>`$, which is a solution to Eq. (B5), and explicitly given by Eq. (B10). For the scattering problem, one is interested in the asymptotic form of the Green’s function, which follows from $$H_\nu ^{(1)}(z)\stackrel{\left(z\mathrm{}\right)}{}\sqrt{\frac{2}{\pi z}}e^{i(z\nu \pi /2\pi /4)},$$ (C4) whence $$G_D^{(+)}(𝐫,𝐫^{};k)=𝒢_D^{(+)}(𝐑;k)\stackrel{\left(r\mathrm{}\right)}{}\frac{1}{4\pi }\left(\frac{k}{2\pi }\right)^{(D3)/2}e^{i\gamma _D}\frac{e^{ikr}}{r^{(D1)/2}}e^{i𝐤^{}𝐫^{}},$$ (C5) where $`𝐤^{}=k𝐫/r`$ and $$\gamma _D=\left(3D\right)\frac{\pi }{4}.$$ (C6) From Eqs. (C3) and (C5), in the position representation, $$\mathrm{\Psi }^{(+)}(𝐫)\stackrel{\left(r\mathrm{}\right)}{}e^{i𝐤𝐫}+f_k^{(D)}(\mathrm{\Omega }^{(D)})e^{i\gamma _D}\frac{e^{ikr}}{r^{(D1)/2}},$$ (C7) where the scattering amplitude $$f_k^{(D)}(\mathrm{\Omega }^{(D)})=\frac{1}{4\pi }\left(\frac{k}{2\pi }\right)^{(D3)/2}d^Dr^{}e^{i𝐤^{}𝐫^{}}V(𝐫^{})\mathrm{\Psi }^{(+)}(𝐫^{})$$ (C8) leads to the usual expression for the differential scattering cross section, $$\frac{d\sigma ^{(D)}(k,\mathrm{\Omega }^{(D)})}{d\mathrm{\Omega }_D}=\|f_k^{(D)}(\mathrm{\Omega }^{(D)})\|^2.$$ (C9) ACKNOWLEDGEMENTS This research was supported in part by CONICET and ANPCyT, Argentina (L.N.E., H.F., and C.A.G.C.) and by the University of San Francisco Faculty Development Fund (H.E.C.). H.E.C. acknowledges instructive discussions with Professor Carlos R. Ordóñez and the generous hospitality of the University of Houston during the final stage of preparation of this article.
warning/0003/hep-th0003253.html	ar5iv	text	# Non-Perturbative Mass Renormalization in Quenched QED from the Worldline Variational Approach ## I Introduction Variational methods are widely used in many areas of physics but are not very prominent in field theory . This is due to the infinite number of degrees of freedom and the singular short-distance behaviour of relativistic field theories. A very successful application of variational methods in a non-relativistic field theory is provided by Feynman’s treatment of the polaron : after integrating out the phonon degrees of freedom and approximating variationally the remaining effective action by a retarded quadratic trial action one obtains the best approximation scheme which works for both small and large coupling constants. Detailed numerical investigations have shown that Feynman’s approximate solution deviates at most $`2.2\%`$ from the true ground state energy for all coupling constants. It is therefore very attractive to apply similar techniques to problems in relativistic quantum field theory where there is much need for non-perturbative methods. In previous publications we have done that in the context of a scalar, super-renormalizable model theory . In this paper we present the first results obtained by applying polaron variational methods to a realistic theory, namely Quantum Electrodynamics (QED) in the quenched approximation where electron-positron loops are neglected. While the actual coupling constant between electrons and photons $`\alpha =e^2/(4\pi )1/137`$ is small enough to apply perturbation theory in most cases, there is enough interest to study the theory at larger coupling: first, the strong coupling behaviour of any physical theory is of interest in itself, second, the possibility of chiral symmetry breaking demands an investigation at large $`\alpha `$ and, finally, bound state problems are inherently non-perturbative and involve powers of $`\mathrm{ln}1/\alpha 4.92`$ in radiative corrections. The extension of our methods to QED requires a formalism to include fermions and a treatment of the more severe singularities encountered in a renormalizable field theory rather than a super-renormalizable or non-relativistic one. We do this within the worldline technique which has recently experienced a revival. In this formulation, the degrees of freedom describing the electron are its bosonic worldline $`x_\mu (t)`$, which is the four-dimensional analogue to the polaron trajectory, as well as a Grassmannian path $`\zeta _\mu (t)`$ needed to describe the electron’s spin . Here $`t`$ is the proper time which parametrizes the paths and runs from $`0`$ to $`T`$. The dynamics of the electron in an external vector field $`A_\mu (x)`$ with field strength $`F_{\mu \nu }(x)`$ are then described by the following worldline Lagrangian $$L=\frac{\kappa _0}{2}\dot{x}^2+i\zeta \dot{\zeta }+\frac{1}{T}\dot{x}\zeta \chi e\dot{x}A(x)\frac{ie}{\kappa _0}F_{\mu \nu }(x)\zeta ^\mu \zeta ^\nu .$$ (1) Here $`\kappa _0`$ is an arbitrary parameter which may be used to reparametrize the proper time without changing the physics and $`\chi `$ is a Grassmannian (super-)partner of the proper time $`T`$. Note that the above action exhibits a well-known supersymmetry between bosonic and fermionic degrees of freedom . For further details about the application of the worldline formalism to QED we refer the reader to Ref. . The photon field $`A_\mu `$ may be integrated out exactly in complete analogy to the phonons in the polaron case, resulting in an effective action for the electron only $`S_{\mathrm{eff}}`$ $`=`$ $`S_0{\displaystyle \frac{e^2}{2}}{\displaystyle _0^T}𝑑t_1𝑑t_2{\displaystyle \frac{d^4k}{(2\pi )^4}G^{\mu \nu }(k)\left[\dot{x}_\mu (t_1)+\frac{2}{\kappa _0}\zeta _\mu (t_1)k\zeta (t_1)\right]}`$ (3) $`\left[\dot{x}_\nu (t_2){\displaystyle \frac{2}{\kappa _0}}\zeta _\nu (t_2)k\zeta (t_2)\right]e^{ik\left[x(t_1)x(t_2)\right]}.`$ Here $`S_0`$ denotes the free action $$S_0=_0^T𝑑t\left[\frac{\kappa _0}{2}\dot{x}^2(t)+i\zeta (t)\dot{\zeta }(t)+\frac{1}{T}\dot{x}(t)\zeta (t)\chi \right]$$ (4) and $`G^{\mu \nu }(k)`$ the gauge-fixed photon propagator. As described in Ref. , the electron propagator is obtained by carrying out a path integral over the degrees of freedom $`x_\mu (t)`$ and $`\zeta _\mu (t)`$, as well as a weighted integral over the proper times $`T`$ and $`\chi `$, and finally identifying the Grassmannian variable $`\mathrm{\Gamma }=\zeta (0)+\zeta (T)`$ with the Dirac matrix $`\gamma `$: $`G_2(p)`$ $`=`$ $`e^{\gamma \frac{}{\mathrm{\Gamma }}}{\displaystyle _0^{\mathrm{}}}𝑑T{\displaystyle 𝑑\chi \mathrm{exp}\left\{\frac{i}{2\kappa _0}\left[(p^2M_0^2)T+(p\mathrm{\Gamma }M_0)\chi \right]\right\}}`$ (6) $`{\displaystyle \frac{𝒟\stackrel{~}{x}𝒟\zeta \mathrm{exp}\left[ipx+\zeta (0)\zeta (T)\right]\mathrm{exp}\left(iS_{\mathrm{eff}}\right)}{𝒟\stackrel{~}{x}𝒟\zeta \mathrm{exp}\left[ipx+\zeta (0)\zeta (T)\right]\mathrm{exp}(iS_0)}}\|_{\mathrm{\Gamma }=0}.`$ Here $`M_0`$ is the bare mass and $`𝒟\stackrel{~}{x}`$ contains an integration over the endpoint $`x=x(T)`$. Note that we have divided and multiplied by the path integral for the free theory, so the bare propagator may be obtained by just ignoring the last line. For non-zero couplings, of course, the path integrals in the last line cannot be performed; these we shall approximate variationally in the next section. ## II Variational Approach Feynman’s variational principle has its root in Jensen’s inequality for convex functions applied to $`\mathrm{exp}(S_E)`$, where $`S_E`$ is a Euclidean action. In Minkowski space and/or for complex actions the variational principle remains valid, however it becomes a stationary principle rather than a minimum principle. To be more precise, the path integral over bosonic and fermionic paths obeys $$<\mathrm{exp}\left[i(SS_t)\right]>_t\stackrel{\mathrm{stat}.}{}\mathrm{exp}\left[iSS_t_t\right],$$ (7) where $`<\mathrm{}>_t`$ indicates an average involving the weight function $`e^{iS_t}`$ in the relevant functional integral and $`S_t`$ is a suitable trial action. Note that corrections to this variational approximation may be calculated in a systematic way and that, furthermore, to first order in the interaction (i.e. to order $`\alpha `$) the relation is in fact an equality if $`S_t`$ reduces to the free action for small couplings. For the trial action required in Eq. (7) we choose a general retarded quadratic action which is a two-time modification of the free action in Eq. (4) $`\stackrel{~}{S}_t`$ $`=`$ $`S_0+i\kappa _0^2{\displaystyle _0^T}dt_1dt_2[g_B(\sigma )\dot{x}(t_1)\dot{x}(t_2)+{\displaystyle \frac{2i}{\kappa _0}}g_F^{}(\sigma )\zeta (t_1)\zeta (t_2)`$ (9) $`\mathrm{\hspace{0.17em}2}{\displaystyle \frac{\sigma }{\kappa _0T}}g_{SO}^{}(\sigma )\dot{x}(t_1)\zeta (t_2)\chi ]+\lambda _1pxi\lambda _2\zeta (0)\zeta (T).`$ Here the variational parameters are contained in the retardation functions $`g_i(\sigma )`$ for bosonic, fermionic and spin-orbit interactions; these are even functions of $`\sigma =t_1t_2`$ and they become identical for a supersymmetric trial action <sup></sup><sup></sup>*We have explicitly separated out the free action in Eq. (9), which could have alternatively been added into the second term by adding $`\delta (\sigma )/(2i\kappa _0)`$ to each of the $`g_i(\sigma )`$’s. This way our retardation functions contain no distributions. Also, in the supersymmetric limit our trial action could be written in the explicitly supersymmetric notation of Ref. as $`i\kappa _0^2_0^T𝑑t_1𝑑t_2𝑑\theta _1𝑑\theta _2g(T_{12})DX_1DX_2`$, with a single retardation function.. The variational principle ‘adjusts’ these functions in order to compensate for the fact that the true effective action (3) is not quadratic in the variables $`x(t)`$, $`\zeta (t)`$. Feynman’s polaron result was obtained by taking a specific Ansatz for the retardation functions but here we leave their functional form free. This is because one expects that the correct short-time behaviour of these functions is much more important for a renormalizable theory like QED than for the polaron problem which does not exhibit any ultraviolet divergences. Indeed one finds that for small $`\sigma `$ the “best” $`g_B(\sigma )`$ behaves like $`\sqrt{\sigma }`$, $`\mathrm{ln}\sigma `$ and $`1/\sigma `$ in the polaron, super-renormalizable and QED case, respectively. The “tilde” over the trial action indicates that it includes the boundary terms already present in Eq. (6) and that we are using “momentum averaging” . These terms, involving the external momentum $`p`$ and the Grassmann variable $`\mathrm{\Gamma }`$, are multiplied by additional variational parameters $`\lambda _1`$ and $`\lambda _2`$, respectively. They provide additional freedom to modify the strength of the boundary terms. We have allowed this freedom because of our experience in scalar relativistic field theory , where the variational parameter $`\lambda _1`$ turned out to be essential for describing the instability of the Wick-Cutkosky model. Since the trial action (9) is at most quadratic in $`x(t)`$ and $`\zeta (t)`$ it is possible to evaluate the various averages required in Eq. (7) analytically. A particular simplification occurs if one restricts oneself to $`p^2=M^2`$, where $`M`$ is the physical (i.e. pole) mass: as discussed in , the divergence of the propagator on its mass shell results from a divergence of the integral over the proper time $`T`$. Indeed, the variational approximation (7) results in an electron propagator \[see Eq. (6)\] which has the form $`G_2^{\mathrm{var}}(p)`$ $`=`$ $`e^{\gamma \frac{}{\mathrm{\Gamma }}}{\displaystyle _0^{\mathrm{}}}𝑑T{\displaystyle 𝑑\chi \mathrm{exp}\left\{\frac{iT}{2\kappa _0}\left[M_0^2+p^2(2\lambda \lambda ^2)\right]\right\}}`$ (11) $`\mathrm{exp}\left\{{\displaystyle \frac{iT}{\kappa _0}}\left(\mathrm{\Omega }[A_B]\mathrm{\Omega }[A_F]+V[\mu _B^2,\mu _F^2]\right)+F(\chi ,\mathrm{\Gamma };T;p)\right\}\|_{\mathrm{\Gamma }=0},`$ where $`\lambda `$ (which is defined below), the $`\mathrm{\Omega }`$’s and $`V`$ are $`T`$-independent and the function $`F(\chi ,\mathrm{\Gamma };T;p)`$ is subleading in $`T`$. The latter therefore contains information relevant for the wavefunction renormalization of $`G_2^{\mathrm{var}}(p)`$, and not the pole structure. We leave the discussion of this function, which also contains the entire $`\chi `$ and $`\mathrm{\Gamma }`$ dependence (and hence the spin structure of the propagator), for a future publication as it is not required for our present investigation. From Eq. (11) we see that the bare and physical mass are related through $$M_0^2=M^2(2\lambda \lambda ^2)2\left(\mathrm{\Omega }[A_B]\mathrm{\Omega }[A_F]+V[\mu _B^2,\mu _F^2]\right).$$ (12) We have labeled this relationship Mano’s equation as K. Mano first applied polaron techniques to a scalar relativistic field theory . Note that, on mass shell, the variational equations resulting from Eq. (7) are equivalent to demanding stationarity of Mano’s equation. The nomenclature in Mano’s equation corresponds to that introduced in Ref. : $`\mathrm{\Omega }[A_B]`$ and $`\mathrm{\Omega }[A_F]`$ originate from contributions (bosonic and fermionic, respectively) of the terms in Eq. (7) involving $`S_0`$ and $`S_t`$ only. They are the analogue to the kinetic term in variational quantum mechanical calculations, while the analogue of the contribution from a potential term (explicitly proportional to the strength of the coupling) resides in $`V`$. Similarly to Ref. , it is useful to express the retardation functions in terms of the variational “profile functions” $`A_i(E)`$ and the “pseudotimes” $`\mu _i^2(\sigma ),i=B,F`$ defined by $`A_i(E)`$ $`=`$ $`1+i\kappa _0{\displaystyle _0^{\mathrm{}}}𝑑\sigma g_i(\sigma )\mathrm{cos}(E\sigma )`$ (13) $`\mu _i^2(\sigma )`$ $`=`$ $`{\displaystyle \frac{4}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑E{\displaystyle \frac{1}{E^2A_i(E)}}\mathrm{sin}^2\left({\displaystyle \frac{E\sigma }{2}}\right),`$ (14) respectively It is easy to show from Eq. (14) that, as long as $`A_i(0)`$ and $`A_i(\mathrm{})`$ exist, for both asymptotically large and infinitesimally small times $`\sigma `$ the functions $`\mu _i^2(\sigma )`$ become proportional to $`\sigma `$; hence the label “pseudotime”. In the free case one has $`A_i(E)=1,\mu _i^2(\sigma )=\sigma ,\lambda =1,\mathrm{\Omega }_i=0`$. The other variational function $`A_{SO}(E)`$ and parameter $`\lambda _2`$ are linked to the spin structure of the propagator and therefore do not show up in Mano’s equation.. Furthermore, it is convenient to define $`\lambda =\lambda _1/A_B(0)`$. Indeed it turns out that the averages in Eq. (7) can be directly expressed in terms of these quantities. The kinetic terms in Eqs. (11) and (12) become $$\mathrm{\Omega }[A_i]=\frac{d\kappa _0}{2i\pi }_0^{\mathrm{}}𝑑E\left(\mathrm{log}A_i(E)+\frac{1}{A_i(E)}\mathrm{\hspace{0.25em}1}\right),$$ (15) where $`d`$ is the spacetime dimension $`d=42ϵ`$. This is identical to the result in Ref. if $`d=4`$ and $`\kappa _0=i`$ (i.e. the Euclidean formulation) are taken. The specific properties of QED are encoded in the “interaction” term $`V`$ which, with $`V=V_1+V_2`$, reads $`V_1[\mu _B^2,\mu _F^2]`$ $`=`$ $`(d1)\pi \alpha {\displaystyle \frac{\nu ^{2ϵ}}{\kappa _0}}{\displaystyle _0^{\mathrm{}}}𝑑\sigma {\displaystyle \frac{d^dk}{(2\pi )^d}\left\{\left[\dot{\mu }_F^2(\sigma )\right]^2\left[\dot{\mu }_B^2(\sigma )\right]^2\right\}E(k,\sigma )}`$ (16) $`V_2[\mu _B^2]`$ $`=`$ $`{\displaystyle \frac{4\pi \alpha \nu ^{2ϵ}\lambda ^2}{\kappa _0}}{\displaystyle _0^{\mathrm{}}}𝑑\sigma {\displaystyle \frac{d^dk}{(2\pi )^d}\frac{1}{k^2}\left[M^2+(d2)\frac{(kp)^2}{k^2}\right]E(k,\sigma )}.`$ (17) Note that by $`\dot{\mu }^2(\sigma )`$ we mean $`\frac{d}{d\sigma }\mu ^2(\sigma )`$ (and not $`\left[\frac{d}{d\sigma }\mu (\sigma )\right]^2`$), the function $`E(k,\sigma )`$ is defined to be $`E(k,\sigma )=\mathrm{exp}\left\{i\left[k^2\mu _B^2(\sigma )2\lambda kp\sigma \right]/(2\kappa _0)\right\}`$ and of course $`p^2=M^2`$. The fermionic contributions, both in the ‘kinetic term’ $`\mathrm{\Omega }_F`$ as well as in $`V_1`$, appear with an opposite sign to the bosonic contributions. The reason for the separation of $`V`$ into two pieces will become apparent below. By construction Mano’s equation is stationary under variation of the parameters. It is important to note that we have not demanded the various retardation functions $`g_{B,F}`$ (as well as $`g_{SO}`$, which only plays a role for the residue) to be identical (before variation). Had we done so, the resulting profile functions $`A_B`$ and $`A_F`$ would have also been identical, the pseudotimes $`\mu _{B,F}^2`$ would have been one and the same and hence $`\mathrm{\Omega }[A_B]\mathrm{\Omega }[A_F]`$ as well as $`V_1`$ would have vanished. The absence of a ‘kinetic’ contribution would have been fatal to the variational principle as this contribution provides the restoring ‘force’ to the potential $`V`$. On the other hand, closer examination of $`V_1`$ reveals that $`\dot{\mu }_B^2\dot{\mu }_F^2`$ is also dangerous: The contribution of each of these terms is quadratically (UV) divergent if the dimensional regularization is replaced by a momentum cutoff. This may be checked by either directly substituting the small $`\sigma `$ limit of $`\mu _i^2(\sigma )`$ into $`V_1`$ or by noting that for scalar QED, where the Grassmannian path integrals are absent, the remaining contribution from $`\dot{\mu }_B^2`$ gives rise to the quadratically divergent one-loop diagram of that theory. It is the combination $`(\dot{\mu }_B^2)^2(\dot{\mu }_F^2)^2`$ which displays the usual logarithmic UV divergence of QED. Although at leading order in the coupling we are guaranteed to reproduce the correct perturbative result \[see Eq. (7)\], at higher orders the cancellation of these quadratic divergences is ensured by the supersymmetry. To summarize, on the one hand the trial action cannot be restricted to contain only supersymmetric terms but on the other hand allowing non-supersymmetric terms may destroy the renormalizability of the theory. The way out of this predicament is provided by the variational principle itself: although it is unavoidable that the trial action breaks supersymmetry, the actual solutions to the variational equations may in fact be nearly supersymmetric. That this indeed turns out to be the case may be seen by recognizing that $`V_1`$ is the most singular part of the interaction whereas $`V_2`$, which involves only bosonic contributions and is the only source of supersymmetry breaking, is similar in structure to the scalar super-renormalizable model studied before . Divergent contributions in the limit $`ϵ0`$ to the variational equations are solely determined by $`V_1`$. Therefore, the divergent contributions to $`A_B(E)`$ and $`A_F(E)`$, and hence to $`\dot{\mu }_B^2`$ and $`\dot{\mu }_F^2`$, are identical. In this paper we confine ourselves to studying this divergent structure and so it is sufficient to set $`A_B(E)=A_F(E)A(E)`$. The corresponding variational equation becomes, after performing the $`k`$-integration in Eq. (16), $$A(E)=\mathrm{\hspace{0.25em}1}+(1ϵ)c_ϵ\nu ^{2ϵ}_0^{\mathrm{}}𝑑\sigma \frac{\mathrm{sin}E\sigma }{E}\frac{\dot{\mu }^2(\sigma )}{\left[\mu ^2(\sigma )\right]^{2ϵ}}\mathrm{exp}\left[i\frac{\lambda ^2M^2\sigma ^2}{2\kappa _0\mu ^2(\sigma )}\right].$$ (18) Note that here we have now also dropped the subscript on the pseudotime as it is no longer relevant and we have defined $$c_ϵ=\frac{\alpha }{\pi }\left(\frac{2\pi i}{\kappa _0}\right)^ϵ\frac{32ϵ}{(1ϵ)(2ϵ)}\stackrel{ϵ0}{}\frac{3\alpha }{2\pi }.$$ (19) Since $`\mu ^2(\sigma )\sigma `$ for small $`\sigma `$ one sees that the $`\sigma `$-integral in Eq. (18) would diverge for $`ϵ=0`$; this just reflects the $`1/\sigma `$ behaviour of the retardation function in Eq. (13) as was discussed before. The crucial difference between super-renormalizable and renormalizable theories therefore is that for the latter ones the variational equations themselves are UV-divergent. In this way the divergent structure of higher-order diagrams is effectively summed up. We may now simplify $`V`$ by making use of the above “asymptotic” supersymmetry. The only remaining contribution is that of $`V_2\lambda ^2M^2W_2`$ which becomes, after carrying out the integration over the momentum $`k`$, $$W_2=\frac{(2ϵ)(1ϵ)}{2}c_ϵ\nu ^{2ϵ}_0^{\mathrm{}}\frac{d\sigma }{\left[\mu ^2(\sigma )\right]^{1ϵ}}_0^1\frac{du}{u^ϵ}[ϵ+(1ϵ)u]\mathrm{exp}\left(\frac{i}{2\kappa _0}\frac{\lambda ^2M^2\sigma ^2}{\mu ^2(\sigma )}u\right)$$ (20) where the $`u`$ integration arises from an exponentiation of the photon propagator in Eq. (17) in a similar way as in Ref. . With this, the variational equation for $`\lambda `$ in this asymptotic limit becomes $$\lambda =\mathrm{\hspace{0.25em}1}\frac{}{\lambda }(\lambda ^2W_2).$$ (21) ## III Mass Renormalization Renormalizability of (quenched) QED means that all divergences can be collected in the mass and wave function renormalization constants. In the present investigation we concentrate on the mass renormalization constant in the MS scheme, $`Z_M^{\mathrm{MS}}`$, defined via $`M_0=Z_M^{\mathrm{MS}}M_\nu `$ where $`M_\nu `$ is an intermediate mass scale. In this scheme it has the perturbative expansion $$Z_M^{\mathrm{MS}}=\mathrm{\hspace{0.25em}1}+\frac{b_{11}}{ϵ}\frac{\alpha }{\pi }+\left[\frac{b_{22}}{ϵ^2}+\frac{b_{12}}{ϵ}\right]\left(\frac{\alpha }{\pi }\right)^2+\mathrm{},$$ (22) where it is known from perturbation theory that the expansion coefficients $`b_{ij}`$ are pure, i.e. mass independent, numbers. Furthermore, the renormalization group provides relations between many of these coefficients; at order $`n`$ in perturbation theory only the coefficient $`b_{1n}`$ contains new information. This is encapsulated in the solution of the renormalization group equation for $`Z_M^{\mathrm{MS}}`$, namely $$Z_M^{\mathrm{MS}}=\mathrm{exp}\left[\frac{1}{2ϵ}_0^\alpha 𝑑x\frac{\gamma _m(x)}{x}\right]=\mathrm{exp}\left[\frac{1}{2ϵ}\underset{n=1}{\overset{\mathrm{}}{}}\frac{\gamma _{n1}}{n}\left(\frac{\alpha }{\pi }\right)^n\right],$$ (23) where $`\gamma _m(\alpha )`$ is the anomalous mass dimension of the electron . In perturbation theory, $`\gamma _m(\alpha )`$ can be extracted from perturbative QCD calculations, which have been performed up to 4-loop order. One obtains $`\gamma _0=3/2,\gamma _1=3/16`$ , $`\gamma _2=\frac{129}{64}=2.0156`$ and $`\gamma _3=\frac{1}{128}\left[\frac{1261}{8}+336\zeta (3)\right]=4.3868`$ . As the variational calculation is applicable for arbitrary values of the coupling, comparison to perturbation theory provides a useful guide to its utility. As mentioned before, to first order in the coupling the calculation is guaranteed to be exact as long as one has used a trial action which can reduce to the free action in the limit $`\alpha 0`$. A genuine test of the variational scheme is only obtained by comparing the coefficients in higher order. It should be noted that this test is much more demanding than in the polaron case where one can only compare the numerical value of the second-order coefficient for the energy: here, in addition, one tests the $`ϵ`$-dependence of this coefficient and also whether it is mass-independent as it should be in the exact theory. In order to know $`\mathrm{\Omega }`$ and $`V`$ at second order in $`\alpha `$ one requires the variational parameters up to first order in $`\alpha `$. These may be obtained by inserting the zeroth order results $`\mu ^2(\sigma )=\sigma `$ and $`\lambda =1`$ into the variational equations for the profile function (18) and $`\lambda `$ (21). The solutions then need to be substituted back into $`\mathrm{\Omega }`$ \[Eq. (15)\] and $`V_2`$ \[Eq. (20)\]. Having done this, $`Z_M^{\mathrm{MS}}`$ may then be extracted from Mano’s equation, yielding $`b_{22}^{\mathrm{var}}=9/32`$, which is correct, and $`b_{12}^{\mathrm{var}}=0`$, which should be compared to the exact value of $`b_{12}=3/64`$. As in the Wick-Cutkosky model, the $`\lambda `$-variation is of crucial importance: for example, fixing $`\lambda =1`$ would give a wrong result for $`b_{22}`$ and a logarithmic mass-dependence for $`b_{12}`$. It is possible to develop the perturbative expansion of the variational result further, with the result that no mass dependence in the coefficients appears even at higher order. Indeed, it turns out that it is in fact possible to obtain the full analytic expression for the anomalous mass dimension in the worldline variational approximation. We shall sketch the derivation below, leaving the technical details for the Appendix to this paper. To begin with, we first drop the mass term in the variational equation (18) since it only affects long-distance physics and not the ultra-violet behaviour contained in $`Z_M^{\mathrm{MS}}`$ . Then we change variables from $`\sigma `$ to $`y=c_ϵ(\nu ^2\sigma )^ϵ`$, and equivalently for $`E`$. This has the effect of making the system of integral equations (14,18,20) independent of the coupling. We write these explicitly in the Appendix, where it is shown that if, for small $`ϵ`$, the pseudotime has the form $$\frac{\mu ^2(\sigma )}{\sigma }=\mathrm{exp}\left[\frac{\omega _0(y)}{ϵ}+O(ϵ^0)\right]$$ (24) then the anomalous mass dimension may be written in terms of this function $`\omega _0(y)`$, i.e. $$\gamma _m^{\mathrm{var}}=\frac{v(y_0)}{1v(y_0)},$$ (25) where $`v(y)=y\omega _0^{}(y)`$ and $`y_0`$ is determined by the implicit equation $$y_0=\frac{3\alpha }{2\pi }e^{\omega _0(y_0)}.$$ (26) On the other hand, it can also be shown (see the Appendix) that the variational equation for the pseudotime translates into an equation for $`\omega _0(y)`$, i.e. $$\frac{e^{\omega _0(y)}}{y}=\frac{\pi }{2}[\mathrm{\hspace{0.17em}1}v(y)]\mathrm{cot}\left[\frac{\pi }{2}v(y)\right].$$ (27) This equation must in general be solved numerically. However, we note that the calculation of the anomalous mass dimension in Eq. (25) only requires knowledge of the function $`v(y)`$ at $`y=y_0`$. Furthermore, it is remarkable that, at this value of $`y`$, the combination $`e^{\omega _0(y)}/y`$ is precisely the combination that is fixed in terms of the coupling constant \[see Eq. (26)\]. Hence, at $`y=y_0`$, the the L.H.S of Eq. (27) may be written in terms of $`\alpha `$ while on the R.H.S. we can eliminate $`v(y_0)`$ completely in terms of $`\gamma _m^{\mathrm{var}}`$ by making use of Eq. (25). One is left with a simple implicit algebraic equation for the anomalous dimension $$\frac{3}{4}\alpha =\left(\mathrm{\hspace{0.17em}1}+\gamma _m^{\mathrm{var}}\right)\mathrm{tan}\left(\frac{\pi /2\gamma _m^{\mathrm{var}}}{1+\gamma _m^{\mathrm{var}}}\right),$$ (28) without ever having actually solved the variational equations themselves. Eq. (28) is the main result of this paper. ## IV Discussion When expanded in powers of $`\alpha `$, Eq. (28) immediately yields $$\gamma _m^{\mathrm{var}}(\alpha )=\frac{3}{2}\frac{\alpha }{\pi }\frac{9}{32}\pi ^2\left(\frac{\alpha }{\pi }\right)^3+\frac{27}{32}\pi ^2\left(\frac{\alpha }{\pi }\right)^4\frac{243}{128}\pi ^2\left(1\frac{\pi ^2}{20}\right)\left(\frac{\alpha }{\pi }\right)^5+𝒪(\alpha ^6),$$ (29) which may be compared to perturbation theory. Numerically the values of the coefficients are different but of the same order of magnitude as the exact perturbative results. Note, however, that this comparison is not particularly meaningful: the variational result is an approximation which is valid at all $`\alpha `$. It need not have the same, or even approximately the same, perturbative expansion in $`\alpha `$ as the exact result. It should, however, be numerically similar. In Fig. 1 we plot the variational result as a function of the coupling and compare it to perturbation theory up to 4-loop order. For $`\alpha `$ $`\stackrel{>}{}`$ $`1`$ the 3- and 4-loop anomalous dimensions start to deviate so much from each other that one cannot trust either of them. Also shown is the result up to 5 loops, where the 5-loop coefficient has been estimated from Padé approximations to the perturbation theory (see Eq. (2.12) of Ref. , which needs to be adapted to QED with $`n_f=0`$ flavours; one finds $`\gamma _4^{\mathrm{Pade}}=\mathrm{\hspace{0.25em}3.848}`$). Clearly this does not significantly extend the numerical validity of the perturbative result. In short, the variational estimate for $`\gamma _m`$ is roughly in agreement with (albeit apparently a little below) the perturbative result in the region where the perturbative result can be trusted. Also shown in Fig. 1 is the only other easily available non-perturbative result for $`\gamma _m^{\mathrm{MS}}`$ based on the use of dimensionally regularized Dyson-Schwinger (DS) equations in “rainbow approximation” within the Landau gauge. This may be obtained by adapting the discussion in Ref. to finite $`M_0`$, with the result that $`\gamma _m^{\mathrm{DS}}=1\sqrt{13\alpha /\pi }`$ (which is the same as derived by Miransky using a hard momentum cutoff). We see that this result deviates from perturbation theory in a region where, at least numerically, perturbation theory still appears to converge. Above $`\alpha =\pi /3=1.047`$ the DS result becomes complex, this value of the coupling constant coinciding with the coupling $`\alpha _{cr}`$ at which the onset of chiral symmetry breaking takes place in those calculations. This is in contrast to the variational result which remains real for all values of the coupling and in fact has the strong coupling limit $$\gamma _m^{\mathrm{var}}(\alpha )\stackrel{\alpha \mathrm{}}{}\frac{1}{4}\sqrt{6\pi \alpha }\frac{1}{2}+𝒪\left(\frac{1}{\sqrt{\alpha }}\right).$$ (30) Further investigations are necessary to clarify the absence of any obvious sign of chiral symmetry breaking in the variational result for $`\gamma _m^{\mathrm{MS}}(\alpha )`$ The reader should note that the issue of dynamical chiral symmetry breaking in a dimensionally regulated theory is a notoriously subtle problem; see Ref. . In particular, it was shown there that if four dimensional quenched QED breaks chiral symmetry above a critical coupling than the dimensionally regularized theory will break it for all couplings at finite $`ϵ`$.. Indeed, in order to investigate the issue of dynamical chiral symmetry breaking, it would seem to be more straightforward, at least conceptually, to set $`M_0`$ on the right hand side of Mano’s equation (12) to zero and to see if the variational equations can be satisfied in this case (for a finite physical mass $`M`$). This, however, goes considerably beyond the scope of this paper: we have merely calculated $`Z_M^{\mathrm{MS}}=M_0/M_\nu `$ (or, more precisely, $`\gamma _m^{\mathrm{MS}}`$), which meant that we could simplify the calculation by i) restricting ourselves to considering the supersymmetric massless limit of the variational equations in Section II and ii) only taking into account the most divergent contributions (as $`ϵ0`$) to the variational equations, as well as to $`W_2`$, in Section III and the Appendix of this paper. In a full calculation of the R.H.S. of Mano’s equation (i.e. the additional calculation of the finite renormalization $`M_\nu /M`$) these two simplifications should not be made <sup>§</sup><sup>§</sup>§This situation is analogous to what is the case in perturbative calculations, where anomalous dimensions of operators are far easier to calculate than finite contributions.. In other words, even if $`Z_M^{\mathrm{MS}}0`$, dynamical chiral symmetry breaking can still occur if $`M_\nu /M`$ vanishes for finite $`M`$. It is interesting to note, however, that there are also some strong similarities in the analytic structure of the variational and DS result. A perturbative inversion of Eq. (28) , i.e. the expansion $`\gamma _m^{\mathrm{var}}(\alpha )=_{n=1}^{\mathrm{}}c_n\alpha ^n`$, has a finite radius of convergence due to a branch cut in the complex $`\alpha `$ plane. The position of this cut, and hence the radius of convergence, can be determined most easily by searching for the value of $`\alpha `$ at which Eq. (28) has two solutions for $`\gamma _m^{\mathrm{var}}`$ which are infinitesimally close to each other. This amounts to demanding that Eq. (28) is satisfied and at the same time the derivative of its R.H.S. vanishes, i.e. $$0=\mathrm{cot}\left(\frac{\pi /2}{1+\gamma _m^{\mathrm{var}}}\right)+\frac{\pi /2}{1+\gamma _m^{\mathrm{var}}}/\mathrm{sin}^2\left(\frac{\pi /2}{1+\gamma _m^{\mathrm{var}}}\right).$$ (31) One finds that $`\alpha _{con}=0.7934`$, which is not too different from the radius of convergence of the DS result . It is not clear whether this similarity between $`\alpha _{cr}`$ and $`\alpha _{con}`$ is accidental or not. In connection with this, it is interesting to note that for large $`n`$ the behaviour of the expansion coefficients $`c_n`$ in both the variational result as well as the DS result are rather similar: $$c_n\alpha _{con}^n\frac{e^\beta }{n^{3/2}}\mathrm{sin}\left[\left(a+\frac{5\pi }{7}\right)n\frac{3\pi }{7}+b\right],$$ (32) where numerically $`\beta 1.38`$, $`a2.3\times 10^3`$ and $`b8.27\times 10^2`$. For the DS result one obtains $`\beta =\mathrm{log}(2\sqrt{\pi })=1.27`$ and the sine function is absent. It is the sine function in the variational result which is responsible for placing the branchpoint (which, for the DS result, is on the positive real axis) into the complex plane. Furthermore, it is remarkable that the large-$`\alpha `$ limit of $`\|\gamma _m(\alpha )\|`$ obtained in Eq. (30) is almost the same as for the DS result: $`\|\gamma _m^{\mathrm{var}}(\alpha )\|1.09\sqrt{\alpha }`$ vs. $`\|\gamma _m^{\mathrm{DS}}(\alpha )\|0.98\sqrt{\alpha }`$. It should be pointed out that a finite radius of convergence of the perturbation expansion is not what one generally expects from calculations of large orders of perturbation theory using the methods of Lipatov and others . Rather, the factorial growth of the number of diagrams at n<sup>th</sup> order in perturbation theory tends to lead to a vanishing radius of convergence. As has been observed elsewhere , it can be shown that the variational calculation contains (pieces of) all possible Feynman diagrams at any order in perturbation theory. One concludes, therefore, that at n<sup>th</sup> order in perturbation theory there are either strong cancellations between diagrams in the variational calculation or that $`𝒪(n!)`$ of them give a vanishing contribution. ## V Summary and Outlook We have applied polaron variational techniques to quenched QED in $`3+1`$ dimensions and obtained, within the MS scheme, a remarkably simple expression for the anomalous mass dimension valid for arbitrary couplings. The approach has considerable advantages over other techniques in that it automatically maintains gauge invariance, as well as the requirements of the renormalization group, and corrections can be systematically calculated (as has been done in the polaron case ). Furthermore, we have shown that the numerical results for $`\gamma _m`$ are rather reasonable at small coupling and that at large couplings the perturbative expansion of this quantity fails in a way similar to rainbow DS results. It would be interesting to compare to DS calculations which go beyond the ladder approximation, thus decreasing the strong gauge dependence inherent in that approximation. Furthermore, variational calculations with more general trial actions could give an indication whether this analytic structure is robust, thus indicating possible large cancellations between diagrams at high order in the perturbation theory of quenched QED, or whether this structure is just an artifact of the particular trial action used in this paper. Finally, we note that the calculation of physical observables or application to bound state problems also seem feasible within the variational worldline approach developed here. ###### Acknowledgements. We would like to thank Reinhard Alkofer for helpful discussions. One of us (AWS) is supported by the Australian Research Council through an Australian Research Fellowship. C.A. would like to thank PSI for its hospitality on several visits during which parts of this work were done. ## Appendix In this Appendix we provide some of the technical details which enter into the derivation of the variational approximation to the anomalous dimension. To begin with, we shall scale the trivial $`\sigma `$ dependence out of $`\mu ^2(\sigma )`$ and define the reduced pseudotime $`s(\sigma )`$ as $$\mu ^2(\sigma )=\sigma s(\sigma ).$$ (A.1) As argued in the main text, mass terms can be dropped for the calculation of the mass anomalous dimension. A perturbative evaluation of the variational equations (18) and (14) for $`M=0`$ then shows that profile function and reduced pseudotime have an expansion in powers of $`E^ϵ`$ and $`\sigma ^ϵ`$, respectively: $$A(E)=\mathrm{\hspace{0.25em}1}+\underset{n=1}{}A_n\left(\frac{\nu ^2}{E}\right)^{nϵ},s(\sigma )=\mathrm{\hspace{0.25em}1}+\underset{n=1}{}s_n\left(\nu ^2\sigma \right)^{nϵ}.$$ (A.2) One finds $`A_1`$ $`=`$ $`c_ϵ\mathrm{\Gamma }(ϵ)\mathrm{cos}\left({\displaystyle \frac{ϵ\pi }{2}}\right)\stackrel{ϵ0}{}{\displaystyle \frac{3\alpha }{2\pi }}{\displaystyle \frac{1}{ϵ}},s_1={\displaystyle \frac{c_ϵ}{ϵ(1+ϵ)}}\stackrel{ϵ0}{}{\displaystyle \frac{3\alpha }{2\pi }}{\displaystyle \frac{1}{ϵ}}`$ (A.3) $`A_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{c_ϵ}{ϵ}}\right)^2{\displaystyle \frac{1ϵ}{1+ϵ}}\mathrm{\Gamma }(1+2ϵ)\mathrm{cos}(ϵ\pi )\stackrel{ϵ0}{}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{3\alpha }{2\pi }}{\displaystyle \frac{1}{ϵ}}\right)^2`$ (A.4) $`s_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{c_ϵ}{ϵ}}\right)^2{\displaystyle \frac{1}{1+2ϵ}}\left[\left(1+{\displaystyle \frac{1}{\mathrm{cos}(ϵ\pi )}}\right){\displaystyle \frac{\mathrm{\Gamma }^2(1+ϵ)}{\mathrm{\Gamma }(1+2ϵ)}}{\displaystyle \frac{1ϵ}{1+ϵ}}\right]\stackrel{ϵ0}{}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{3\alpha }{2\pi }}{\displaystyle \frac{1}{ϵ}}\right)^2.`$ (A.5) This suggests that perhaps the leading $`ϵ`$-behaviour of the coefficients is $$A_n\stackrel{ϵ0}{}\frac{1}{n!}\left(\frac{3\alpha }{2\pi }\frac{1}{ϵ}\right)^n,s_n\stackrel{ϵ0}{}\frac{(1)^n}{n!}\left(\frac{3\alpha }{2\pi }\frac{1}{ϵ}\right)^n.$$ (A.6) In view of these results we rewrite all equations in terms of the dimensionless quantities $$y=c_ϵ\left(\nu ^2\sigma \right)^ϵ,z=c_ϵ\left(\frac{\nu ^2}{E}\right)^ϵ$$ (A.7) where we have also rescaled by $`c_ϵ`$ (which is linear in the coupling) because the mass scale $`\nu `$ always appears in the combination $`\alpha \nu ^{2ϵ}`$ in dimensional regularization. In a similar way, the variational parameter $`\lambda `$ almost always appears in the combination $`\lambda M`$, so it is convenient to define the dimensionless combination $$a_ϵ=c_ϵ\left(\frac{2\kappa _0\nu ^2}{i\lambda ^2M^2}\right)^ϵ.$$ (A.8) With these definitions, the massless variational equation (18) for $`A(E)`$ may be brought into the form $$A(z)=\mathrm{\hspace{0.25em}1}+\frac{1}{ϵ}_0^{\mathrm{}}𝑑y\frac{\mathrm{cos}(y/z)^{1/ϵ}}{[s(y)]^{1ϵ}}.$$ (A.9) The reduced pseudotime \[ i.e. the rescaled version of Eq. (14) \] is now given by $$s(y)=\frac{2}{\pi }\frac{1}{ϵ}_0^{\mathrm{}}𝑑z\frac{1}{z}\left(\frac{z}{y}\right)^{1/ϵ}\frac{1\mathrm{cos}(y/z)^{1/ϵ}}{A(z)},$$ (A.10) while the rescaled potential $`W_2`$ becomes a function of $`a_ϵ`$ alone: $$W_2(a_ϵ)=\frac{(2ϵ)(1ϵ)}{2ϵ}_0^{\mathrm{}}\frac{dy}{[s(y)]^{1ϵ}}_0^1\frac{du}{u^ϵ}[ϵ+(1ϵ)u]\mathrm{exp}\left[(y/a_ϵ)^{1/ϵ}\frac{u}{s(y)}\right]$$ (A.11) and hence the variational equation (21) for $`\lambda `$ becomes $$\frac{1}{\lambda }=\mathrm{\hspace{0.25em}1}+\mathrm{\hspace{0.25em}2}W_2(a_ϵ)\mathrm{\hspace{0.25em}2}ϵa_ϵW_2^{}(a_ϵ).$$ (A.12) The anomalous mass dimension in the MS scheme may be defined \[ see Eq. (23) \] through $$\gamma _m=\underset{ϵ0}{lim}ϵ\frac{\alpha }{Z_M^2}\frac{}{\alpha }Z_M^2.$$ (A.13) Note that this equation is correct independently of whether $`Z_M`$ has been calculated in the MS scheme or whether it is defined through Mano’s equation by $`Z_MM_0/M`$. It is because of this fact that we can derive a nonperturbative expression for $`\gamma _m`$ in the MS scheme, even though this scheme is usually only used within the context of perturbation theory. As the $`\alpha `$ dependence of $`Z_M`$ now only enters through the variable $`a_ϵ`$ (and of course implicitly through the variational parameters), it is not surprising that we may use the variational equation (A.12) for $`\lambda `$ to simplify $`\gamma _m`$. Indeed, differentiating Mano’s equation with respect to the coupling gives $$\frac{}{\alpha }Z_M^2=\frac{\lambda }{\alpha }\frac{}{\lambda }Z_M^2\mathrm{\hspace{0.25em}2}\frac{a_ϵ}{\alpha }\lambda ^2W_2^{}(a_ϵ).$$ (A.14) The first term is zero because of the variational equation for $`\lambda `$, $`a_ϵ/\alpha `$ is just $`a_ϵ/\alpha `$ and by substituting the variational equation for $`\lambda `$ into Mano’s equation we find $$Z_M^2=\lambda [\mathrm{\hspace{0.17em}1}2ϵ\lambda a_ϵW_2^{}(a_ϵ)].$$ (A.15) Hence the anomalous mass dimension is just given by $$\gamma _m=\underset{ϵ0}{lim}\frac{2ϵ\lambda a_ϵW_2^{}(a_ϵ)}{1\mathrm{\hspace{0.25em}2}ϵ\lambda a_ϵW_2^{}(a_ϵ)}.$$ (A.16) In order to proceed further, we need to evaluate $`W_2(a_ϵ)`$. In general one would need to do this numerically, however fortunately in Eq. (A.16) only the small-$`ϵ`$ limit is required. Let us assume that the reduced pseudotime may be written as $$s(y)=\mathrm{exp}\left[\frac{\omega (y,ϵ)}{ϵ}\right],$$ (A.17) where $`\omega _0(y)lim_{ϵ0}\omega (y,ϵ)`$ is finite. This is supported by the perturbative results given in Eqs. (A.2) and (A.6) and we shall show that this holds in general when we solve the variational equations below. In this case the exponential in Eq. (A.11) has the argument $$\left(\frac{y}{a_ϵ}e^{\omega (y,ϵ)}\right)^{1/ϵ}u.$$ (A.18) If the term in brackets is larger than one, this argument will become arbitrarily large (and negative) in the limit $`ϵ0`$, hence it will lead to a vanishing contribution to the integral. If the term in brackets is smaller than one, however, the argument goes to zero, the exponential factor in Eq. (A.11) may be replaced by unity and the integral over $`u`$ may be performed, yielding $$W_2(a_ϵ)\stackrel{ϵ0}{}\frac{1}{2ϵ}_0^{y_0}𝑑y\mathrm{exp}\left[\frac{\omega (y,ϵ)}{ϵ}(1ϵ)\right],$$ (A.19) where $`y_0`$ is given by the equation $$y_0e^{\omega _0(y_0)}=\underset{ϵ0}{lim}a_ϵ.$$ (A.20) We have assumed here that $`ye^{\omega _0(y)}`$ is an increasing function of $`y`$, which will turn out to be the case. The leading term in Eq. (A.19) may be obtained by integration by parts, with the result $$W_2(a_ϵ)\stackrel{ϵ0}{}\frac{1}{2\omega ^{}(y_0,ϵ)}\mathrm{exp}\left[\frac{\omega (y_0,ϵ)}{ϵ}(1ϵ)\right].$$ (A.21) We also require the derivative of this function, which is most easily obtained by direct differentiation of Eq. (A.19): $$W_2^{}(a_ϵ)\stackrel{ϵ0}{}\frac{1}{2ϵ}\mathrm{exp}\left[\frac{\omega (y_0,ϵ)}{ϵ}(1ϵ)\right]\frac{e^{\omega (y_0,ϵ)}}{1+y_0\omega ^{}(y_0,ϵ)}.$$ (A.22) Substitution into the variational equation for $`\lambda `$ yields $$\lambda \stackrel{ϵ0}{}\omega ^{}(y_0,ϵ)\frac{1+y_0\omega ^{}(y_0,ϵ)}{\mathrm{exp}\left[\frac{\omega (y_0,ϵ)}{ϵ}(1ϵ)\right]}$$ (A.23) and hence $$2ϵ\lambda a_ϵW_2^{}(a_ϵ)\stackrel{ϵ0}{}y_0\omega ^{}(y_0,ϵ)$$ (A.24) so that the anomalous dimension becomes $$\gamma _m=\frac{y_0\omega _0^{}(y_0)}{1y_0\omega _0^{}(y_0)}.$$ (A.25) We stress that only the last line is exact while the previous ones have correction terms for finite $`ϵ`$. In particular, the calculation of $`Z_M=M_0/M`$ (as opposed to $`Z_M^{\mathrm{MS}}=M_0/M_\nu `$) would require these additional terms and hence the result, unfortunately, does not shed light on whether $`Z_M`$ could in fact be zero for finite $`M`$, which would signal chiral symmetry breaking. Note that the limit $`ϵ0`$ in Eq. (A.20) needs some care: naively, one would conclude from the definition (A.8) that the R.H.S. equals $`lim_{ϵ0}c_ϵ=3\alpha /(2\pi )`$ but Eq. (A.23) shows that the variational parameter $`\lambda `$ vanishes like $`\mathrm{exp}(\omega _0(y_0)/ϵ)`$ and therefore also gives a contribution: $$y_0e^{\omega _0(y_0)}=\frac{3\alpha }{2\pi }e^{2\omega _0(y_0)},$$ (A.26) this being the result (26) quoted in the main text. It now remains to calculate the function $`\omega _0(y)`$. The arguments used to derive the approximate expression for $`W_2(a_ϵ)`$ in Eq. (A.19) are more difficult to apply to the variational equation (A.9) for $`A(z)`$ and the definition (A.10) of $`s(y)`$ because of the rapidly oscillating trigonometric functions appearing in their integrands. We shall therefore adopt a more systematic approach at this stage and note that it is possible to write these equations in a differential form. Consider, for example, an integral of the type $$I_ϵ(z)=_0^{\mathrm{}}dyf(y)\mathrm{cos}\left(\frac{y}{z}\right)^{1/ϵ}.$$ (A.27) By changing integration variable to $`y^{1/ϵ}`$ and Taylor expanding the function $`f(y)`$ we may carry out the integration term by term by making use of the integral $$_0^{\mathrm{}}𝑑yy^{q1}\mathrm{cos}y=\mathrm{\Gamma }(q)\mathrm{cos}\left(q\frac{\pi }{2}\right).$$ (A.28) Hence we obtain $$I_ϵ(z)=\underset{n=0}{\overset{\mathrm{}}{}}f^{(n)}(0)\frac{z^{n+1}}{(n+1)!}\mathrm{\Gamma }\left[1+(n+1)ϵ\right]\mathrm{cos}\left[(n+1)ϵ\frac{\pi }{2}\right].$$ (A.29) This expression may be resummed, by defining the dilatation operator $`D_zz\frac{d}{dz}`$, into the compact form $$I_ϵ(z)=\mathrm{\Gamma }(1+ϵD_z)\mathrm{cos}\left(\frac{\pi }{2}ϵD_z\right)_0^zdyf(y)=:\gamma _c\left(ϵD_z\right)_0^zdyf(y).$$ (A.30) Hence the variational equation for $`A(z)`$ becomes $$A(z)=\mathrm{\hspace{0.25em}1}+\frac{1}{ϵ}\gamma _c(ϵD_z)_0^z𝑑y\frac{1}{[s(y)]^{1ϵ}}$$ (A.31) and in a similar way we can rewrite Eq. (A.10) as $$s(y)=\frac{1}{1+ϵD_y}\frac{1}{\gamma _c(ϵD_y)}\frac{1}{A(y)}.$$ (A.32) Inverting Eq. (A.32) and substituting into Eq. (A.31) eliminates the profile function $`A(z)`$: $$\frac{1}{\left(1+ϵD_y\right)\gamma _c(ϵD_y)s(y)}=\mathrm{\hspace{0.25em}1}+\frac{1}{ϵ}\gamma _c(ϵD_y)_0^y𝑑x\frac{1}{[s(x)]^{1ϵ}}.$$ (A.33) Finally, one can eliminate the integral by operating with $`D_y`$ on both sides of this equation, so that $$D_y\frac{1}{(1+ϵD_y)\gamma _c(ϵD_y)s(y)}=\frac{1}{ϵ}\gamma _c(ϵD_y)\frac{y}{[s(y)]^{1ϵ}}.$$ (A.34) This equation may be solved systematically by defining $`s(y)`$ in terms of the function $`\omega (y,ϵ)`$ \[ see Eq. (A.17) \] and by making the Ansatz that $`\omega (y,ϵ)`$ has a power expansion in $`ϵ`$ $$\omega (y,ϵ)=\omega _0(y)+ϵ\omega _1(y)+\mathrm{}.$$ (A.35) The crucial observation is that repeated application of the dilatation operator on an exponential of the form of Eq. (A.17) results in $`\left(ϵD_y\right)^n\mathrm{exp}\left(\omega /ϵ\right)=\left[(y\omega ^{})^n+𝒪(ϵ)\right]\mathrm{exp}\left(\omega /ϵ\right)`$ so that at leading order in $`ϵ`$, for any function $`F(ϵD_y)`$ acting on $`\mathrm{exp}(\pm \omega /ϵ)`$, we have $$F(ϵD_y)\mathrm{exp}(\pm \omega /ϵ)\stackrel{ϵ0}{}F(\pm y\omega _0^{})\mathrm{exp}(\pm \omega /ϵ).$$ (A.36) Applying this relation to Eq. (A.34) provides the following equation for $`\omega _0(y)`$ $$\frac{\omega _0^{}}{(1y\omega _0^{})\gamma _c(y\omega _0^{})}=\gamma _c(y\omega _0^{})e^{\omega _0}.$$ (A.37) which is $`ϵ`$-independent, justifying the Ansatz (A.35) a postiori. This equation may be simplified considerably by making use of the reflection formula $`\mathrm{\Gamma }(z)\mathrm{\Gamma }(1z)=\pi /\mathrm{sin}\pi z`$ for $`\mathrm{\Gamma }`$-functions. By defining $`v(y)y\omega _0^{}(y)`$, we then find $$\frac{e^{\omega _0(y)}}{y}=\frac{\pi }{2}[1v(y)]\mathrm{cot}\left[\frac{\pi }{2}v(y)\right],$$ (A.38) which is Eq. (27) in the main text. Together with the boundary condition $`\omega _0(0)=0`$ \[ i.e. $`\mu ^2(\sigma )\sigma `$ for $`\sigma 0`$, as discussed in the footnote below Eq. (14) \] the first-order nonlinear differential equation (A.38) determines the function $`\omega _0(y)`$. Remarkably, as shown in the main text, it is not actually necessary to solve it in order to obtain the anomalous mass dimension $`\gamma _m`$. It is also interesting to note that due to the reflection formula all $`\mathrm{\Gamma }`$-functions have disappeared, which has the consequence that in a perturbative expansion of $`\gamma _m^{\mathrm{var}}`$ no Riemann $`\zeta `$-functions, but only powers of $`\pi `$, occur.
warning/0003/gr-qc0003006.html	ar5iv	text	# Some Remarks on the Bel–Robinson Tensorfootnote footnote This paper is an amended version of an Essay which received a ”honorable mention” in the 1999 Essay Competition of the Gravity Research Foundation. ## I Introduction In analogy to the symmetric energy–momentum tensor of the electromagnetic field, the Bel–Robinson superenergy tensor is defined as follows (see eg. ) $`T^{iklm}`$ $`=`$ $`R^{iabl}R_{ab}^{k}{}_{}{}^{m}+R^{iabl}R_{ab}^{k}{}_{}{}^{m}`$ (1) $`=`$ $`R^{iabl}R_{ab}^{k}{}_{}{}^{m}+R^{iabm}R_{ab}^{k}{}_{}{}^{l}{\displaystyle \frac{1}{2}}g^{ik}R^{abcl}R_{abc}^m,`$ (2) where $`R_{iabl}`$ is the Riemann tensor and $``$ indicates the usual dual operation $$R_{iabl}:=\frac{1}{2}\eta _{iadc}R_{bl}^{dc}.$$ (3) The unusual properties (see eg. ) of the Bel–Robinson tensor intrigued physicists and a large number of papers were devoted to the understanding of its physical sense within the framework of the General Relativity (GR). Recently, in the Ref. (see also ) the authors try to connect this tensor with square of the energy–momentum and derive from it (by taking a suitable defined "square root") an energy–momentum tensor of the gravitational field. In some special cases such square root really exists. However, it seems that this is an incorrect physical idea. In order to understand the physical content of the Bel–Robinson tensor correctly in GR one should take into account the following, fundamental facts: 1. The Bel–Robinson tensor can be obtained as a consequence of the Bianchi identities independently of the Einstein equations. 2. In the framework of the standard General Relativity (GR) The standard General Relativity has a very good experimental confirmation; especially its main postulate — Einstein Equivalence Principle (see eg. ).the gravitational field has non–tensorial strenghts $`\mathrm{\Gamma }_{kl}^i`$ and has no (and cannot have any) energy–momentum tensor but only the so–called pseudotensors. It is a consequence of the Einstein Equivalence Principle (EEP). 3. The Bel–Robinson tensor appears explicitly in the expansion of the differences of the gravitational energy–momentum calculated in normal coordinates NC(P) by use of the canonical energy–momentum pseudotensor $`{}_{E}{}^{}t_{i}^{k}`$. Namely, it is a part of the differences <sup>§</sup><sup>§</sup>§The analogous expansions were obtained by using other energy–momentum pseudotensors of the gravitational field and also contain Bel–Robinson tensor. However, they are much more complicated; see e.g. . $`{}_{E}{}^{}t_{i}^{k}(y)_Et_i^k(P)`$. Here $`y`$ means normal coordinates NC(P) which have point P as origin. In this paper, in Sec. II and in Sec. III, we will consider the above three facts more intensively and show that they uniquely indicate the link between Bel–Robinson tensor and differences of the canonical gravitational energy–momentum calculated in NC(P). In Sec. IV we give conclusions and some remarks. ## II The Bel–Robinson tensor and Bianchi identities The Bel–Robinson tensor is a special case of the so–called Maxwellian tensors. The Maxwellian tensors generalize the symmetric energy–momentum tensor of the electromagnetic field onto antisymmetric tensor fields. The general method of construction of such tensors was developed in . In the following we apply this method to the Riemann tensor. Let us consider Bianchi identities for the Riemann tensor $`R_{iklm}=R_{lmik}=R_{kilm}=R_{ikml}`$ $$_{[a}R_{bc]de}_aR_{bcde}+_bR_{cade}+_cR_{abde}0$$ (4) and their non–vanishing contractions $$_aR_{cd}^{ab}2_{[c}R_{d]}^b.$$ (5) We realize that the identities (3)–(4) possess Maxwellian structure in the indices $`(a,b,c)`$. Let us multiply (3) by $`R_f^{bcd}`$. Then, after simple calculations we get the new identities $$R_f^{bcd}_bR_{acde}\frac{1}{2}R_f^{bcd}_aR_{bcde}0.$$ (6) Then let us transpose the indices $`f`$ and $`e`$ in (5) $$R_e^{bcd}_bR_{acdf}\frac{1}{2}R_e^{bcd}_aR_{bcdf}0.$$ (7) The sum of (5) and (6) gives $`R_f^{bcd}_bR_{acde}`$ $`+`$ $`R_e^{bcd}_bR_{acdf}`$ (8) $``$ $`{\displaystyle \frac{1}{2}}\left(R_f^{bcd}_aR_{bcde}+R_e^{bcd}_aR_{bcdf}\right)0.`$ (9) Due to the identities (4) one can rewrite the identities (7) in the following form $`_b(R_f^{bcd}R_{acde}`$ $`+`$ $`R_e^{bcd}R_{acdf}{\displaystyle \frac{1}{2}}\delta _a^bR_e^{bcd}R_{bcdf})`$ (10) $``$ $`2R_{ac}^{d}{}_{e}{}^{}_{[d}R_{f]}^c+2R_{ac}^{d}{}_{f}{}^{}_{[d}R_{e]}^c.`$ (11) The Bel–Robinson tensor $`T_{aef}^b`$ is easily indicated inside parenthesis on the left hand side of the identities (8) which determine the divergence of this tensor. Thus, we see that the Bel–Robinson tensor and its divergence arise as a consequence of the Bianchi identities (3) and their contractions (4) only and they are neither connected with Einstein equations nor with the canonical formalism of the energy–momentum in GR. However, by using of the Einstein equations $$R_k^i=\beta (T_k^i\frac{1}{2}\delta _k^iT)=:\beta E_k^i,$$ (12) where $`\beta =8\pi G/c^4`$, one can rewrite the identities (8) in the form $$_bT_{aef}^b=2\beta R_{ac}^{d}{}_{e}{}^{}_{[d}E_{f]}^c+2\beta R_{ac}^{d}{}_{f}{}^{}_{[d}E_{e]}^c.$$ (13) The equations (10) give the link between the divergence of the Bel–Robinson tensor and GR. It follows from (10) that, in vacuum, $$_bT_{aef}^b=0.$$ (14) The dimensions of the components of the Bel–Robinson tensor are $`(length)^{()4}`$; but it is a trivial fact that $`1/\beta ^2T_{afe}^b`$ have dimensions of the energy–momentum square. This trivial fact was used in Ref. 3 with the aim of connecting the Bel–Robinson tensor with square of an energy–momentum tensor. ## III The relation between the Bel–Robinson tensor and the canonical energy–momentum in General Relativity The problem of the energy–momentum in General Relativity (GR) was intensively studied by many authors (see e.g. \[10—17\]). The main results of these investigations are the following: We do not consider here the so–called quasilocal quantities \[18—22\] because already for the Kerr spacetime differently defined quasilocal quantities give different results (see eg. ). Moreover, the term ”quasilocal” is very obscure. We omit also Lorentz hypothesis (see eg. ) which is unsatisfactory from the physical point of view since it gives for gravitational field an energy–momentum tensor which vanishes in vacuum. But we think that Lorentz hypothesis is the best of the all trials to attribute an energy–momentum tensor to gravitational field. 1. Owing to the non–tensorial character of the gravitational strengths $`\mathrm{\Gamma }_{kl}^i=\{_{kl}^i\}`$ the gravitational field in standard GR has no (and cannot have) any energy–momentum tensor. Any attempt to introduce such a tensor leads us beyound standard GR. Moreover, it is speculative and contradicts EEP. From that it follows the non–localizability of the gravitational energy–momentum. 2. The best solution of the energy–momentum problem in standard GR seems to be given by the so–called canonical energy–momentum pseudotensor $`{}_{E}{}^{}t_{i}^{k}`$ proposed for gravitational field by Einstein and related to that pseudotensor, the canonical, double index, energy–momentum complex $${}_{E}{}^{}K_{i}^{k}:=\sqrt{\|g\|}\left(T_i^k+_Et_i^k\right),$$ (15) for matter and gravitation which satisfies $$\sqrt{\|g\|}\left(T_i^k+_Et_i^k\right)=_FU_{i}^{kl}{}_{,l}{}^{}.$$ (16) Here $${}_{F}{}^{}U_{i}^{kl}=()_FU_i^{lk}=\alpha \frac{g_{ia}}{\sqrt{\|g\|}}\left[(g)\left(g^{ka}g^{lb}g^{la}g^{kb}\right)\right]_{,b}$$ (17) are von Freud superpotentials, $`T_i^k`$ are the components of a symmetric energy–momentum tensor of matter (the sources in the Einstein equations) and $`g`$ is the determinant of the metric tensor; $`,i`$ or $`_i`$ denotes partial derivative. $`\alpha =1/2\beta =c^4/16\pi G`$. We have \[25—27\] $`{}_{E}{}^{}t_{i}^{k}`$ $`=`$ $`\alpha \{\delta _i^kg^{ms}(\mathrm{\Gamma }_{mr}^l\mathrm{\Gamma }_{sl}^r\mathrm{\Gamma }_{ms}^r\mathrm{\Gamma }_{rl}^l)`$ (18) $`+`$ $`g_{,i}^{ms}[\mathrm{\Gamma }_{ms}^k{\displaystyle \frac{1}{2}}(\mathrm{\Gamma }_{tp}^kg^{tp}\mathrm{\Gamma }_{tl}^lg^{kt})g_{ms}`$ (19) $``$ $`{\displaystyle \frac{1}{2}}(\delta _s^k\mathrm{\Gamma }_{ml}^l+\delta _m^k\mathrm{\Gamma }_{sl}^l)]\}.`$ (20) The equations (13) can be obtained by rearranging of the Einstein equations having $`T_i^k`$ as sources. From (13) there follow the local or differential conservation laws $$\left[\sqrt{\|g\|}\left(T_i^k+_Et_i^k\right)\right]_{,k}=0,$$ (21) and, by using Stokes integral theorem, the integral conservation laws $$\underset{\mathrm{\Omega }}{}\sqrt{\|g\|}\left(T_i^k+_Et_i^k\right)𝑑\sigma _k=0.$$ (22) $`\mathrm{\Omega }`$ is the boundary of a four–dimensional, compact domain $`\mathrm{\Omega }`$, and $`d\sigma _k`$ is the three–dimensional integration element (see e.g. ). The components $`{}_{E}{}^{}t_{i}^{k}`$ of the Einstein canonical energy–momentum pseudotensor of the gravitational field neither form a tensor nor other geometric object. Any attempt of the physical interpretation of the Bel–Robinson tensor in the framework of GR should take into account the connection of the Bel–Robinson tensor and its divergence with Bianchi identities, the above two fundamental facts and the next, more important fact as follows: the Bel–Robinson tensor appears explicitly in the expansion of the differences $${}_{E}{}^{}t_{i}^{k}(y)_Et_i^k(P)$$ (23) of the canonical energy–momentum calculated in normal coordinates \[28—30\] NC(P)(And in analogic expansions obtained when using other pseudotensors too ).Because $`{}_{E}{}^{}t_{i}^{k}(P)={}_{E}{}^{}t_{i}^{k}{}_{,l}{}^{}(P)=0,`$ one can also speak about expansion of $`{}_{E}{}^{}t_{i}^{k}`$ alone. But the differences $`{}_{E}{}^{}t_{i}^{k}(y)_Et_i^k(P)`$ have deeper physical meaning; for example, they admit introduction of superenergy and supermomentum tensors \[31—37\]. This fact gives the most important connection between the Bel–Robinson tensor and GR as follows. In the NC(P) $`\{y^i\}`$ having the point P as their origin we have $`{}_{E}{}^{}t_{i}^{k}(y)_Et_i^k(P)`$ $`=`$ $`{}_{E}{}^{}t_{i}^{k}(y)={\displaystyle \frac{1}{2}}{}_{E}{}^{}t_{i}^{k}{}_{,lm}{}^{}(P)y^ly^m+R_3`$ (24) $`=`$ $`{\displaystyle \frac{\alpha }{9}}[T_{ilm}^k(P)+P_{ilm}^k(P){\displaystyle \frac{1}{2}}\delta _i^kR_l^{abc}(P)(R_{abcm}(P)+R_{acbm}(P))`$ (25) $`+`$ $`2\delta _i^kR_{(l\|g}(P)R_{\|m)}^g(P)3R_{i(l\|}(P)R_{\|m)}^k(P)+R_{gi(l\|}^k(P)R_{\|m)}^g(P)`$ (26) $`+`$ $`R_{ig(l\|}^k(P)R_{\|m)}^g(P)]y^ly^m+R_3.`$ (27) In the formula (19) $`R_3`$ is the remainder of the third order, while $$T_{ilm}^k:=R_l^{kab}R_{iabm}+R_m^{kab}R_{iabl}\frac{1}{2}\delta _i^kR_l^{abc}R_{abcm}$$ (28) are the Bel–Robinson tensor components, and $$P_{ilm}^k:=R_l^{kab}R_{ibam}+R_m^{kab}R_{ibal}\frac{1}{2}\delta _i^kR_l^{abc}R_{acbm}$$ (29) are components of the tensor which is closely related to the Bel–Robinson tensor. By using the Einstein equations in the form (9), one can rewrite (19) to the form $`{}_{E}{}^{}t_{i}^{k}(y)_Et_i^k(P)`$ $`=`$ $`{}_{E}{}^{}t_{i}^{k}(y)={\displaystyle \frac{\alpha }{9}}[T_{ilm}^k(P)+P_{ilm}^k(P)`$ (30) $``$ $`{\displaystyle \frac{1}{2}}\delta _i^kR_l^{abc}(P)\left(R_{abcm}(P)+R_{acbm}(P)\right)`$ (31) $`+`$ $`2\delta _i^k\beta ^2E_{(l\|g}(P)E_{\|m)}^g(P)3\beta ^2E_{i(l\|}(P)E_{\|m)}^k(P)+\beta R_{gi(l\|}^k(P)E_{\|m)}^g(P)`$ (32) $`+`$ $`\beta R_{ig(l\|}^k(P)E_{\|m)}^g(P)]y^ly^m+R_3.`$ (33) In vacuum we have from (22) $`{}_{E}{}^{}t_{i}^{k}(y)`$ $``$ $`{}_{E}{}^{}t_{i}^{k}(P)=_Et_i^k(y)={\displaystyle \frac{\alpha }{9}}[T_{ilm}^k(P)+P_{ilm}^k(P)`$ (34) $``$ $`{\displaystyle \frac{1}{2}}\delta _i^kR_l^{abc}(P)(R_{abcm}(P)+R_{acbm}(P))]y^ly^m+R_3`$ (35) $`=`$ $`{\displaystyle \frac{4\alpha }{9}}\left[R_{(l\|}^{k(ab)}(P)R_{iab\|m)}(P){\displaystyle \frac{1}{2}}\delta _i^kR_l^{a(bc)}(P)R_{abcm}(P)\right]y^ly^m+R_3.`$ (36) We see from the above formulas (19)–(23) that the Bel–Robinson tensor really appears in the expansion of the differences (18). Years ago, by using the expansion (23), we have shoved that the infinitesimal differences $`\mathrm{}P_a`$ of the free gravitational energy–momentum calculated in NC(P) are proportional to the components $`T_{a00}^0`$ of the Bel–Robinson tensor multiplied by $`\alpha =c^4/16\pi G`$. ## IV Conclusion The Bel–Robinson tensor follows from the Bianchi identities as the Maxwellian tensor for the Riemann tensor $`R_{iklm}`$ and, within the framework of the standard GR, it can be connected with the differences of the canonical gravitational energy–momentum calculated in NC(P). It is easily seen from Sec. II, from the formulas (19)–(23) and from the Ref. 38. These facts give the most natural (and correct) physical interpretation of this tensor in the framework of the GR. Although the gravitational field $`\mathrm{\Gamma }_{kl}^i`$ has no energy–momentum tensor in standard GR, one can easily introduce there the so–called canonical superenergy tensor \[31—37\] for this field. The method of construction of the canonical superenergy tensor for gravitational field uses the expansion (19) and some kind of averaging. The Bel–Robinson tensor multiplied by $`\alpha =c^4/16\pi G`$ is the main, "Maxwellian part" of such a tensor. The canonical superenergy tensor for gravitational field does not vanish in vacuum; so, it gives us, for example, a very useful tool for local analysis of the gravitational radiation . The idea of superenergy and its tensor is universal and easy to generalize onto matter field too; for example, one can easily introduce the canonical superenergy tensor for matter \[31—35\]. On the "superenergy level" one can easily introduce the canonical angular supermomentum tensors for matter and gravitation as well. The canonical superenergy tensors of gravitation and matter and the canonical angular supermomentum tensors of gravitation and matter have much better geometrical and physical properties than the canonical objects from which they were obtained. Also, the integral superenergetic quantities have better properties than corresponding integral energetic quantities; especially, the superenergetic integrals have better convergence in asymptotically flat spacetimes (at spatial or null infinity). Finally, the canonical superenergy and angular supermomentum tensors give a very good tool for local (and also to global) analysis of the gravitational and matter fields within the framework of the standard GR. Acknowledgments I would like to thank Dr M.P. Da̧browski for his help in preparation English version of this paper.
warning/0003/cond-mat0003016.html	ar5iv	text	# Landau Theory of the Finite Temperature Mott Transition \[ ## Abstract In the context of the dynamical mean-field theory of the Hubbard model, we identify microscopically an order parameter for the finite temperature Mott endpoint. We derive a Landau functional of the order parameter. We then use the order parameter theory to elucidate the singular behavior of various physical quantities which are experimentally accessible. \] When the strength of the electron-electron interaction $`U`$ is increased compared to the bare bandwidth $`2D`$, a metal-insulator transition (MIT) occurs . This phenomenon, known as the Mott transition, can take place in the absence of magnetic long-range order, and is still an outstanding problem in condensed-matter physics. From a theoretical point of view, a difficulty is the absence of an obvious order parameter to systematize the critical behavior of the observable quantities when the metal insulator transition is not accompanied by the onset of magnetic long range order. These issues are experimentally relevant to systems such as V<sub>2</sub>O<sub>3</sub> and Ni(Se,S)<sub>2</sub> and are the subject of intensive experimental study. In recent years, great progress has been made by using the dynamical mean-field theory (DMFT) . This framework describes both paramagnetic metallic and paramagnetic insulating phases. The $`U`$-$`T`$ phase diagram ($`T`$ is the temperature) of the frustrated Hubbard model in the limit of large lattice coordination is qualitatively similar to that of the V<sub>2</sub>O<sub>3</sub> and Ni(Se,S)<sub>2</sub> systems: A first-order phase-transition line ends in a second-order critical point, henceforth referred to as the Mott critical point, which is the main focus of this letter. We will use this framework to address the fundamental questions raised in the previous paragraph. There are two earlier qualitative ideas as to what should be the order parameter to describe the physics around the finite temperature Mott point. One idea is to connect the order parameter to the notions of ”metallicity” or coherence. It can be traced back to the early paper of Brinkman and Rice and is captured in a slave boson formalism where the metallic state has a non zero expectation value of a Bose field which describes the coherent propagation of one particle excitations. In a very different picture, Castellani et al. viewed the metal as a liquid rich in doubly occupied sites, and the insulator as a liquid with few doubly occupied sites. The metal to insulator transition is viewed as a condensation of doubly occupied sites, and the order parameter is related to the Blume-Emery-Griffith model . The Landau approach presented here provides a synthesis of these ideas. It bridges naturally between a picture based on one particle excitations and a picture based on local collective excitations (or double occupancies). In agreement with Castellani et al. we find that the Mott transition has indeed an Ising-like character. On the other hand, we obtain a complementary description in terms of the one particle spectral function reminiscent of the slave boson picture. A simple and clear description of the critical behavior near the critical point emerges. It allows us to systematically derive the critical behavior of any observable quantity and to relate its non analytic dependence on $`T`$ and $`U`$ to that of the order parameter. Our results should be also of help in resolving some controversies on the solution of the Hubbard model in infinite dimensions by providing a theoretical framework in which to analyze numerical results on the finite temperature Mott transition. It can also be used to analyze results of photoemission and optical conductivity experiments. For simplicity, we focus on the single-band Hubbard model at half-filling, $$\widehat{H}=\frac{t}{\sqrt{z}}\underset{ij\sigma }{}c_{i\sigma }^+c_{j\sigma }+U\underset{i}{}\widehat{n}_i\widehat{n}_i.$$ (1) The first term describes the hopping between nearest neighbors on a lattice with coordination number $`z`$. The corresponding half bandwidth is our unit of energy, $`D=2t=1`$. The second term is an on-site interaction suppressing double occupancies by imposing an energy cost $`U`$ on each one. In the limit of infinite dimensions, $`z\mathrm{}`$, this model can be mapped onto a single-impurity Anderson model (SIAM) supplemented by a self-consistency condition. We adopt a semicircular density of states, which is realized on the Bethe lattice. The dynamical mean-field equations can be obtained by differentiating the Landau functional $$F_{\text{LG}}[\mathrm{\Delta }]=T\underset{n}{}\frac{\mathrm{\Delta }(i\omega _n)^2}{t^2}+F_{\text{imp}}[\mathrm{\Delta }],$$ (2) with respect to the hybridization function $`\mathrm{\Delta }(i\omega _n)`$ of the SIAM, which has the meaning of a Weiss field. $`i\omega _n`$ are fermionic Matsubara frequencies, while $`F_{\text{imp}}[\mathrm{\Delta }]`$ is the free energy of the SIAM, given by the action $`S_{\text{imp}}=S_{\text{loc}}[\mathrm{\Delta }=0]+_{\sigma ,n}f_\sigma ^+(i\omega _n)\mathrm{\Delta }(i\omega _n)f_\sigma (i\omega _n)`$. Here, $`S_{\text{loc}}[\mathrm{\Delta }=0]`$ is the action of the local $`f`$ level with the hybridization set to zero. The first term in Eq. (2) is the cost of forming the Weiss field $`\mathrm{\Delta }(i\omega _n)`$ around a given site, while the second one is the free energy of an electron at this site in the presence of the Weiss field. Using the Green’s function of the SIAM, $`G(i\omega _n)=(1/2T)\delta F_{\text{imp}}/\delta \mathrm{\Delta }(i\omega _n)`$, the mean-field equation reads $$\frac{t^2}{2T}\frac{\delta F_{\text{LG}}[\mathrm{\Delta }]}{\delta \mathrm{\Delta }(i\omega _n)}=t^2G(i\omega _n)[\mathrm{\Delta },\alpha ]\mathrm{\Delta }(i\omega _n)=0.$$ (3) Here, $`\alpha =(U,T)`$ comprises the control parameters. This Landau approach was used to describe the energetics of the Mott transition at zero temperature . We will show that near the finite temperature Mott point, the Weiss field has a singular dependence which can be parametrized by a single number which assumes the role of an effective order parameter for this transition. As in Landau theory, we assume that a finite temperature transition exists, and derive a complete description of the critical behavior near the transition as follows: First, we expand the mean-field equation (3) around the critical point, $`\alpha _c=(U_c,T_c)`$, up to third order in the deviation of the hybridization function from its value at the critical point, $`\delta \mathrm{\Delta }=\mathrm{\Delta }(\alpha _c+\delta \alpha )\mathrm{\Delta }(\alpha _c)`$, and to first order in $`\delta \alpha =(UU_c,TT_c)`$. This expansion is well-behaved because the impurity model at finite temperatures depends smoothly on $`\alpha `$ and $`\delta \mathrm{\Delta }(i\omega _n)`$. In order to carry out this expansion it is convenient to define a fluctuation matrix $$M_{nm}=\frac{t^2}{2T}\frac{\delta ^2F_{\text{LG}}[\mathrm{\Delta }]}{\delta \mathrm{\Delta }(i\omega _n)\delta \mathrm{\Delta }(i\omega _m)}\|_{\text{critical point}}$$ (4) $`M_{nm}`$ has the form $`\delta _{nm}+K_{nm}`$, where $`K_{nm}`$ is the Fourier transform of a kernel $`K(\tau ,\tau ^{})`$ which is proportional to the connected correlation function of an operator $`O(\tau )=_{0}^{}{}_{}{}^{\beta }𝑑uf^+(u+\tau )f(u)`$, $`<O(\tau )O(\tau ^{})><O(\tau )><O(\tau ^{})>`$ where the average $`<>`$ is calculated with the action of an Anderson impurity model. It is well known that the correlation functions of the Anderson impurity model are bounded, and therefore the Kernel $`K`$ is square integrable $`_{0}^{}{}_{}{}^{\beta }_{0}^{}{}_{}{}^{\beta }𝑑\tau 𝑑\tau ^{}\|K(\tau ,\tau ^{})\|^2<\mathrm{}`$. Therefore it $`K_{nm}`$ is a Fredholm operator which and has a discrete spectrum of eigenvalues which we labeled by the index $`l`$. At half-filling, particle-hole symmetry guarantees that the order parameter $`\mathrm{\Delta }(i\omega )`$ is odd and wholly imaginary. Accordingly, the fluctuation matrix is real and symmetric and has real eigenvalues $`m_l`$ belonging to eigenvectors $`\varphi _l(i\omega _n)`$ which can be chosen to be purely imaginary and to form an orthonormal basis. The critical point, in this description of the problem, is signaled by the appearance of a single zero eigenvalue, $`m_0=0`$, which indicates the occurrence of a simple bifurcation. Next, we represent $`\delta \mathrm{\Delta }`$ in the eigenbasis of the matrix (4), $`\delta \mathrm{\Delta }(i\omega _n)=_l\eta _l\varphi _l(i\omega _n)`$, where all $`\eta _l`$ are real. By projecting the mean-field equation (3) onto the eigenbasis $`\varphi _l`$, we obtain an equation of the form $`m_l\eta _l+F_l^{(0)}[\{\eta _{j0}\}]+F_l^{(1)}[\{\eta _{j0}\}]\eta _0`$ (5) $`+F_l^{(2)}[\{\eta _{j0}\}]\eta _0^2+F_l^{(3)}\eta _0^3=0`$ , (6) which holds for all $`l`$. $`F_l^{(0)}`$ is of order $`\delta \alpha `$. $`F_l^{(1)}`$ and $`F_l^{(2)}`$ have Taylor expansions in the $`\eta _{j0}`$, where $`F_l^{(1)}`$ starts with the linear order. We solve Eq. (6) iteratively for all $`\eta _{l0}`$ to obtain $`\eta _{l0}=a_l+c_l\eta _0^2+d_l\eta _0^3`$. Here, $`a_l`$ is of first order in $`\delta \alpha `$, (which assures us that the leading singular dependence of the spectral function is proportional to $`\varphi _0`$) further corrections have the form $`b_l\eta _0`$ with $`b_l`$ also of order $`\delta \alpha `$. By inserting this expression into the $`l=0`$ case of Eq. (6), we derive an effective equation for the zero-mode amplitude $`\eta _0`$. We can think of $`\eta _0`$ as the soft mode near the transition and the $`\eta _{l0}`$ as massive modes. The elimination of the massive modes renormalizes the coefficients of the effective action for the soft mode. In the resulting cubic equation for $`\eta _0`$, we eliminate the quadratic term by shifting $`\eta _0`$ by an appropriately chosen linear function in $`\delta \alpha `$, $`\eta =\eta _0+\text{const}_1\times (TT_c)+\text{const}_2\times (UU_c)`$. Close to the critical point, $`\eta `$ and $`\eta _0`$ are dominated by non analytic terms and are therefore essentially equal. We thus obtain an equation of state without quadratic term in $`\eta `$: $$p\eta +c\eta ^3=h.$$ (7) Here, all quantities are real. As in Landau theory, a microscopic calculation of the Landau coefficients (p,c,h) is difficult. However we can extract exact information about the critical behavior from the knowledge that they are smooth functions of the control parameters, i.e. $`c`$ is finite at the critical point, whereas $`p`$ and $`h`$ are linear functions of $`\delta \alpha `$, $`h=h_1(UU_c)+h_2(TT_c)`$ and $`p=p_1(UU_c)+p_2(TT_c)`$. As a consequence, $`\eta `$ has a singular dependence on $`U`$ and $`T`$ near the critical point. At $`U=U_c`$, and for $`T`$ near $`T_c`$, $$\eta (U_c,T)\text{sign}(h_2/c)\text{sign}(TT_c)\|TT_c\|^{1/3}.$$ (8) The mean-field equation (7) describes the Mott transition close to the critical point in terms of the order parameter $`\eta `$. In this form, the analogy with the liquid gas transition is evident. The Mott transition takes place on the line in the $`U`$-$`T`$ plane where $`h`$ vanishes and the system has full Ising symmetry. The critical point, $`(U_c,T_c)`$, divides this line into two half-lines. On the half-line where $`T<T_c`$, there are two solutions, $`\eta =\pm \sqrt{\|p/c\|}`$. We will see later that $`\eta `$ parametrizes the strength of the quasiparticle resonance of the single-particle spectrum (see Fig. 2). A positive or negative ”field” $`h`$ increases or decreases this component of the spectral function, respectively. The field $`h`$ decreases when $`U`$ or $`T`$ is increased, because either increase eliminates the metallic coherence and thus reduces the value of $`\eta `$. We have used the sign convention whereby is positive. We now turn to various consequences of our construction. From Eq. (7), we can obtain the shape of the coexistence region near the critical point, where two solutions of the mean field equations coexist. It is centered symmetrically about the $`h=0`$ line, and its width along $`T=\text{const}`$ lines, $`\mathrm{\Delta }U`$, scales with $`(T_cT)^{3/2}`$. The constant of proportionality is given by $`(4/\sqrt{c}\|h_1\|)[(p_2p_1h_2/h_1)/3]^{3/2}`$. An important quantity which is measured in numerical simulations is the double occupancy. It is connected to our order parameter $`\eta `$ as follows: $`d=(T/U)_n\{[(i\omega _n+\mu )G(i\omega _n)1]e^{i\omega _n0^+}t^2G(i\omega _n)^2\}=d_c+c_1^{(d)}\eta +c_2^{(d)}\eta ^2`$. In this expansion about the critical point, we have only retained the leading and next to leading nonanalytic terms responsible for the critical behavior. The susceptibility $`\chi =d/U`$ diverges at the critical point. For example: $$\chi (U,T_c)(c_1^{(d)}/3)\text{sign}(h_1/c)\|h_1/c\|^{1/3}\|UU_c\|^{2/3}.$$ (9) The double occupancy is related to the magnetization by the identity $`(n_{}n_{})^2=12d`$. The magnetic response will therefore also exhibit nonanalytic dependences on the control parameters. There has been several numerical studies of the finite temperature Mott transition in this model. The Landau approach predicts the functional dependence of various quantities near the transition, and therefore the expressions derived in this paper, are useful for interpreting the numerical work. To illustrate how our approach sheds new light on previously obtained numerical data we compare in Fig. 1 the results for the double occupancy $`d`$ obtained within the IPT and QMC calculations with $`\mathrm{\Delta }\tau =0.5/D`$, after carrying out the shifts and the rescaling described in the figure caption. Within the statistical errors of the QMC calculation, the agreement is excellent. This surprising result is consistent with the Landau theory: different approximations for the solution of the impurity model reduce to the same Landau theory near the critical point, but with different values of the Landau coefficients. Therefore, with a suitable rescaling, the results near the critical point should agree with each other, and with a fit based on the Landau theory which is shown in the red line in figure 1. Small changes in the values of $`\mathrm{\Delta }\tau `$ result in shifts of $`U_c`$, $`T_c`$, and $`d`$ at criticality, but does not change the form of the critical behavior. We also note that the critical slowing down which has been observed in the iterative solutions of the mean field equations are a direct consequence of the presence of the soft mode $`\eta `$ described in the Landau approach. From our construction it is clear that $`\eta `$ provides the leading non analytic behavior of the Weiss field. In order to get a better feeling for its physical significance we have to understand how it can be probed experimentally. Since the order parameter is closely related to the amplitude of the quasiparticle peak, photoemission is an ideal tool to probe the temperature and pressure dependence of the order parameter near the critical point. This experimental technique, in the angle integrated mode, would also measure the convolution of the Fermi function with the analytically continued eigenfunction of the zero mode, $`\text{Im}\varphi _0(i\omega _n=\omega i\delta )`$. To visualize the shape of the spectral function near the critical point we must resort to calculations based on analytic methods such as IPT. The inset of Fig. 2 shows the spectral function very near the critical point, computed within the IPT. It illustrates how the compromise between metallic and insulating features is realized. A finite $`\eta `$, depending on its sign, adds or subtracts spectral weight to the coherent low energy feature immersed in a constant backround in between the Hubbard bands. The zero mode is seen to affect mainly the low-energy part of the spectrum, which determines whether the system is metallic or insulating. The strong temperature dependence has been noticed in previous theoretical and experimental studies. Its origin and connection to an order-parameter description of the Mott transition, however, had not been recognized until now. In the main panel of Fig. 2 we display the height of the quasiparticle peak $`A_0=i\mathrm{\Delta }(i0^+)/\pi t^2`$, for $`UU_c`$, as a function of temperature in the vicinity of $`T_c`$. The rapid variation seen in the figure is consistent with the form $`A_0=A_{0c}+c_1^{(A)}\eta +c_2^{(A)}\eta ^2`$ with coefficients $`c_i^{(A)}`$ independent of U and temperature. Optical techniques are probably the best tool available to test the predictions of our theory. For instance, one may consider the integral of the optical conductivity up to some cuttoff, $`N_{\text{eff}}(T)`$. Since the optical conductivity in infinite dimensions is directly expressed in terms of the single-particle Green’s function, $`N_{\text{eff}}(T)`$ must also exhibit the singular temperature dependence near the transition. We would therefore expect the temperature variation of this quantity to be most visible for a relatively small cuttoff, displaying a rapid variation with $`T`$ similarly as for $`A_0`$. Since the singular dependence arises from the order parameter $`\eta `$, it should be possible to fit the Drude weight by $`N_{\text{eff}}(T)=N_{\text{eff}}(T_c)+c_1^{(N)}\eta (T)+c_2^{(N)}\eta ^2(T)`$. $`N_{\text{eff}}(T)`$ has recently been measured in NiS<sub>2-x</sub>Se<sub>x</sub> , the observed strong temperature dependence of the effective number of carriers is consistent with our predictions. In summary, we derived an order parameter description of the Mott transition near its critical point in the $`U`$-$`T`$ plane. We showed that the critical behavior in proximity to this point is governed by an Ising-like Landau functional and is present in a large number of observable quantities. We predict that any physical quantity which is sensitive to the single-particle spectrum exhibits singular dependences on the control parameters close to the finite-temperature Mott point. The leading non analytic behavior of other physical quantities can be obtained along similar lines, i.e. by recognizing their coupling to the order parameter. This involves a few coefficients, (i.e. the $`c^{(A)}`$’s) which depend on the observable (and on the approximation method) and, as in Landau theory, should be taken as parameters. The dependence on temperature and on pressure is completely determined from the temperature or pressure dependence of the order parameter that follows from Eq. (7). ACKNOWLEDGMENT This work was supported by NSF 95-29138. E.L. was partially supported by the Deutsche Forschungsgemeinschaft. M.J.R. acknowledges support of Fundación Antorchas, CONICET (PID $`N^o4547/96`$), and ANPCYT (PMT-PICT1855). We thank R. Chitra for discussions and D. Vollhardt for useful comments on the mansucript.
warning/0003/quant-ph0003115.html	ar5iv	text	# Coherent States of Non-Linear Algebras: Applications to Quantum Optics. ## 1 Introduction Till recently, in quantum optics, only linear Lie algebras have been used to give mutiphoton coherent (including squeezed) states . It is well-known that if we have bilinear Hamiltonians for two mode radiation fields characterized by operators $`a`$, $`b`$, $`a^{}`$ and $`b^{}`$, then the simplest types of coherent states that can be constructed are the product states $`\|\alpha >\|\beta >`$ where $`\|\alpha >`$ and $`\|\beta >`$ are the eigenstates corresponding to $`a`$ and $`b`$ respectively. However, if the system has an added symmetry or conservation law, then, a set of coherent states restricted by the extra symmetry can be constructed, by a suitable projection from the ordinary product states. Examples of such coherent states include the coherent states of a radiation field with arbitrary polarization such that $`a^{}a+b^{}b`$ is conserved. Here, the symmetry algebra is $`SU\left(2\right)`$, and the corresponding states are the $`SU\left(2\right)`$ coherent states . In the frequency conversion of photons of a given frequency $`\mathrm{\Omega }`$ into two photons of frequencies $`\omega _a`$ and $`\omega _b`$, when the two photons are created or destroyed together such that the operator $`Q=a^{}ab^{}b`$ is conserved, the relevant states are the ‘pair coherent states’ or the SU(1,1) ‘Barut-Girardello (BG)’ states. . The symmetry algebra in this case is $`SU(1,1)`$, defined by $`K_+=a^{}b^{},K_{}=ab,K_0=\frac{1}{2}\left(a^{}a+b^{}b+1\right)`$ and the Casimir operator $`C=\frac{1}{4}\left(1\left(a^{}ab^{}b\right)^2\right)=\frac{1}{4}\left(1Q^2\right)`$. The coherent states are the simultaneous eigenstates of $`K_{}`$ and $`C`$. Other coherent states are the SU(1,1) Perelomov states, $`\lambda >=e^{\left(\lambda K_+\lambda ^{}K_{}\right)}0,0>`$ ,which are the ‘Caves-Schumaker’ states that represent two-mode squeezing . A third set of eigenstates $`A\psi >=\left(uK_0+\nu K_{}+cK_+\right)\psi >=\stackrel{~}{\lambda }\psi >`$ called algebraic coherent states have also been constructed by Agarwal, Shanta et.al , Trifonov and others . The BG states can be produced by a dissipative process governed by Hamiltonians of the form $`H=H_0+\omega _1ab+\omega _1^{}a^{}b^{}=H_0+H_{int}`$, as the steady states of the master equation $`\frac{d\rho }{dt}=[\rho ,H_{int}]`$. The $`SU(1,1)`$ Perelomov states represent the time evolution of the same Hamiltonian. Thus, the corresponding Lie algebraic structure has proved instrumental in studying the quantum optical properties of two mode radiation fields. We show, in this paper, that when one has Hamiltonians representing interactions of multimode radiation fields, i.e., three or more modes, then the dynamical symmetry algebra of the Hamiltonian becomes a polynomially deformed algebra. The deformation is quadratic for the three mode case and cubic for the four mode case. The quadratic algebra was discovered by Sklyanin , in the context of statistical physics and field theory . It was shown to be the symmetry algebra of a two-dimensional anisotropic harmonic oscillator and the isotropic harmonic oscillator in curved space . The well-known Higgs algebra, a cubic algebra, occurs in the study of the dynamical symmetries of the Coulomb problem in a space of constant curvature. These algebras have now found a place in quantum optics with the observation that quantum optical Hamiltonians describing multiphoton processes have symmetries which can be described by polynomially deformed SU(1,1) and SU(2) algebras . A polynomial deformation of a Lie algebra is defined in the following fashion in the Cartan-Weyl basis, $$[H,E_\pm ]=\pm E_\pm ,[E_+,E_{}]=f\left(H\right),$$ (1) where $`f\left(H\right)`$ is a polynomial function of $`H`$. The corresponding Casimir can be written in the form , $`C`$ $`=`$ $`E_{}E_++g\left(H\right),`$ (2) $`=`$ $`E_+E_{}+g\left(H1\right),`$ where, $$f\left(H\right)=g\left(H\right)g\left(H1\right).$$ (3) The form of $`g\left(H\right)`$ can be determined up to the addition of a constant. The eigenstates are characterized by the values of the Casimir operator and the Cartan subalgebra $`H`$. In particular, a polynomial deformation of $`SL(2,R)`$ is of the form $`N_0=J_0`$, $`N_+=F(J_0,J)J_+`$, $`N_{}=F(J_0,J)J_{}`$, where the $`J_i`$’s are the ordinary $`SL(2,R)`$ generators. The commutation relations are $`[N_0,N_\pm ]=\pm N_\pm `$ and $`[N_+,N_{}]=F\left(N_0\right)`$ . When $`F\left(N_0\right)`$ is quadratic in $`N_0`$ the algebra is called a quadratic algebra and if it is cubic in $`N_0`$ the ”Higgs” algebra results. As an example of occurrence of non-linear algebras in quantum optics, consider the triboson Hamiltonian $$H=\omega _aa^{}a+\omega _bb^{}b+\omega _cc^{}c+\kappa ab^{}c^{}+\kappa ^{}a^{}bc.$$ (4) Physically, for this Hamlitonian, $`a`$,$`b`$ and $`c`$ represent the pump, signal and idler modes for parametric amplification and the idler, pump and signal modes in frequency conversion. Raman and Brillouin scattering can be described by $`H`$, if $`a`$, $`b`$ and $`c`$ represent input, vibrational and Stokes modes for a Stokes process and anti-Stokes, input and vibrational modes for an anti-Stokes process. It also represents the interaction of N identical two-level atoms with a single mode radiation field. This has been considered by many authors by ordinary linear Lie algebraic methods leading to approximate results for specific cases. Infinite dimensional Lie algebraic techinques have also been attempted and the physics has been extracted by a truncation of these algebras, hence the results obtained have again been approximate, with a number of assumptions . In this paper, we show that this Hamiltonian and its generalizations have a non-linear algebra as its dynamical symmetry algebra and the construction of the coherent states is straightforward using the representation theory of these algebras. Furthermore, all three types of coherent states can be constructed on the basis of our method. We shall show that this Hamiltonian is formed by operators which obey a finite quadratic polynomial deformation of $`SL(2,R)`$ and the construction of CS for this Hamiltonian is a fairly straightforward process based on the undeformed algebra. For the Hamiltonian given in equation (4), let $`a`$ represent a pump system and $`b`$ and $`c`$ represent the signal and idler variables. The interaction Hamiltonian between the pump and signal-idler subsystem is given by $$H_{int}=\kappa ab^{}c^{}+\kappa ^{}a^{}bc.$$ (5) Energy conservation requires that $`\omega _a=\omega _b+\omega _c`$. If the signal and idler frequencies are equal then $`\omega _b=\frac{\omega _a}{2}`$ and $`\omega _c=\frac{\omega _b}{2}`$ with, $`H_0^{free}=\omega _a\left(a^{}a+\frac{b^{}b+c^{}c}{2}\right)`$. The generators of the polynomial quadratic algebra are defined by the operators $$J_0=\frac{1}{2}\left(a^{}aK_0\right)$$ (6) $$J_{}=ab^{}c^{}=aK_+$$ (7) $$J_+=a^{}bc=a^{}K_{}$$ (8) where $`K_0`$, $`K_{}`$ and $`K_+`$ form SU(1,1) generators. The algebra closes only if we define an additional conserved quantity $`H_0`$ given by : $$H_0=\frac{a^{}a+K_0}{2}.$$ (9) Since $`H_0`$ is related to $`H_0^{free}`$ through $`H_0^{free}=2H_0\frac{1}{2}`$, we see that physically this condition is satisfied. In fact, $`H_0`$ can also be related to the Manley-Rowe invariants of the system. The algebra is given by: $$[J_+,J_{}]=3J_0^2+\left(2H_01\right)J_0C_{bc}\left(K_0\right)+H_0\left(H_0+1\right)$$ (10) Where $`C_{bc}=\frac{1}{4}\frac{\left(b^{}bc^{}c\right)^2}{4}=\frac{1Q^2}{4}`$ is the Casimir operator for the idler-signal system , for which Q is a conserved quantity. For the case special $`b^{}bc^{}c=0`$ , $`C_{bc}=\frac{1}{4}.`$ The two commuting generators are then $`H_0`$ and $`J_0`$ and a general eigenstate of the system is labelled by the quantum numbers corresponding to their eigenvalues and is given by $`\|h_0,j_0>`$. Similarly, the symmetry algebra for four photon processes is a Higgs algebra. For general multiphoton Hamiltonians: $$H=H_0+\kappa \left(a_0\right)^m\left(a_1^{}\right)^n+c.c$$ (11) we can define $`N_0,N_{},N_+`$ in an analogous way $`N_+`$ $`=`$ $`a_0^m\left(a_1^{}\right)^n`$ $`N_{}`$ $`=`$ $`a_1^n\left(a_0^{}\right)^m`$ $`N_0`$ $`=`$ $`{\displaystyle \frac{1}{m+n}}\left(a_1^{}a_1a_0^{}a_0\right)`$ (12) and show that we get n-dimensional polynomial algebras as the symmetry algebra if $`H_0=\frac{1}{m+n}\left(a_0^{}a_0+a_1^{}a_1\right)`$ is conserved. Similarly n-photon Dicke Models with Hamiltonians of the form: $$H=H_0+\kappa \underset{i}{}\sigma _{}\left(i\right)\left(a_1^{}\right)^n+\kappa ^{}\underset{i}{}\sigma _+\left(i\right)\left(a_1\right)^n$$ (13) with $`N_0`$ $`=`$ $`{\displaystyle \underset{i}{\overset{n}{}}}\sigma _0\left(i\right)a_1^{}a_1`$ $`N_{}`$ $`=`$ $`{\displaystyle \underset{i}{\overset{n}{}}}\sigma _{}\left(i\right)\left(a_1^{}\right)^n`$ $`N_+`$ $`=`$ $`{\displaystyle \underset{i}{}}\sigma _+\left(i\right)\left(a_1\right)^n`$ (14) satisfying a polynomial Lie Algebra of order n if $`H_0=ϵ_i^n\sigma _0\left(i\right)+w_1a_1^{}a_1`$ is conserved. We present a unified approach for finding the coherent states (CS) of these algebras . Apart from its application to quantum optics, the method of construction presented here is quite general and will greatly facilitate the physical applications of these algebras to many quantum mechanical problems. This method is a generalisation to non-linear algebras of the procedure for constructing multiphoton states outlined in reference . For ordinary Lie algebras, the construction of the CS for the non-compact cases, was shown to be a two step procedure . First, the canonical conjugate of the lowering operator were found and the CS of these algebras were obtained by the action of the exponential of the respective conjugate operators on the vacuum . This method was in complete parallel to the one used for constructing the coherent states for harmonic oscillator algebras. Another CS, dual to the first one, naturally follows from the above construction. Here, we generalise the above construction to non-linear algebras and provide a mapping between the deformed algebras and their undeformed counterparts. This mapping is utilized to find the CS in the Perelomov sense. Apart from obtaining the known CS of the $`SU(1,1)`$ algebra, we construct the CS for the quadratic and cubic polynomial algebras. Other coherent states in the literature which are essentially special cases of this construction are the ‘f-oscillator states’ and the non-linear states $`f\left(N_0\right)a\|\lambda >=\lambda \|\lambda >`$, which have been shown to be useful for the trapped ion problem. While these are non-linear harmonic oscillator coherent states, the CS that we construct may be called non linear SU(1,1) (or SU(2)) coherent states. These states would give a multi-mode generalization of the type $`f(n_a,n_b)a^nb^m\|\lambda >=\lambda \|\lambda >`$, as one of the possible coherent states. Thus our construction encompasses existing non-linear states and allows for the construction of new physical states. One such state, for example, is the case n=1 and m=1, which is a two-mode realization of the non-linear coherent states. Our method is quite general and encompasses q-deformations of linear Lie algebras. In this work, we concentrate on finite, polynomial SU(2) and SU(1,1) algebras, in view of their importance in quantum optics. ## 2 Construction of Coherent States of Non-Linear Algebras:Formalism. Having seen that polynomially deformed algebras occur in a large class of systems, we now give the formalism for the construction of coherent states of these algebras. For the purpose of clarity, we start with Lie Algebras and then extend the method to the deformed algebras in a straightforward way. In the next section, we shall show how this formalism can be used to explicitly construct the coherent states for application to mutiphoton processes. We introduce our method by first considering SU(1,1) for which the generators satisfy the commutation relations $$[K_+,K_{}]=2K_0,[K_0,K_\pm ]=\pm K_\pm .$$ (15) Thus for this case one finds, $`f\left(K_0\right)=2K_0`$ and $`g\left(K_0\right)=K_0\left(K_0+1\right)`$. The quadratic Casimir operator is given by $`C=K_{}K_++g\left(K_0\right)=K_{}K_+K_0\left(K_0+1\right)`$. $`\stackrel{~}{K_+}`$, the canonical conjugate of $`K_{}`$, satisfying $$[K_{},\stackrel{~}{K_+}]=1,$$ (16) can be written in the form, $$\stackrel{~}{K_+}=K_+F(C,K_0).$$ (17) Eq.(17) then yields, $$F(C,K_0)K_{}K_+F(C,K_01)K_+K_{}=1;$$ (18) making use of the Casimir operator relation given earlier, one can solve for $`F(C,K_0)`$ in the form, $$F(C,K_0)=\frac{K_0+\alpha }{C+K_0\left(K_0+1\right)}.$$ (19) The constant, arbitrary, parameter $`\alpha `$ in F can be determined by demanding that Eq.$`\left(17\right)`$ be valid in the entire Hilbert space. For the purpose of illustration, we demonstrate our method, by using the one oscillator realization of the $`SU(1,1)`$ generators $`K_{}=a^2`$,$`K_+=a^^2`$, $`K_0=\frac{1}{4}\left(aa^{}+a^{}a\right)`$. Since the coherent states of this realization have been studied extensively in the literature , this provides a good testing ground for our method. The ground states defined by $`K_{}0>=\frac{1}{2}a^20>=0`$, are, $`0>`$ and $`1>=a^{}0>`$, in terms of the oscillator Fock space. $$K_00>=\frac{1}{4}\left(2a^{}a+1\right)0>=\frac{1}{4}0>,$$ (20) and $$C0>=\frac{3}{16}0>.$$ (21) Thus $`[K_{},\stackrel{~}{K_+}]0>=K_{}\stackrel{~}{K_+}0>`$ yields $`\alpha =\frac{3}{4}`$. Similarly, for the other ground state $`\|1>`$, $$[K_{},\stackrel{~}{K_+}]1>=1>,$$ (22) leads to $`\alpha =\frac{1}{4}`$. Hence, there are two disjoint sectors characterized by the $`\alpha `$ values $`\frac{3}{4}`$ and $`\frac{1}{4}`$, respectively. These results match identically with the earlier known ones , once we rewrite $`F`$ as, $$F(C,K_0)=\frac{K_0+\alpha }{C+K_0\left(K_0+1\right)},$$ (23) $$=\frac{K_0+\alpha }{K_{}K_+}.$$ (24) The unnormalized coherent state $`\beta >`$, which is the annihilation operator eigenstate, i.e, $`K_{}\beta >=\beta \beta >`$, is given in the vacuum sector $`\|0>`$ by $$\beta >=e^{\beta \stackrel{~}{K^+}}0>.$$ (25) For the vacuum sector $`\|1>`$, where $`\alpha =\frac{1}{4}`$, a similar construction holds. These states, which provide a realization of the Cat states, play a prominent role in quantum measurement theory. The canonical conjugate $`\stackrel{~}{K_+}`$ such that: $$[\stackrel{~}{K_+^{}},K_+]=1.$$ (26) can be constructed, as in ref. from this, one can find the eigenstate of $`\stackrel{~}{K_+^{}}`$ operator, in the form, $$\gamma >=e^{\gamma K_+}0>.$$ (27) This CS, after proper normalization, is the well-known Yuen (squeezed) state: $`e^{\mu a^^2\mu ^{}a^2}`$, with $`\gamma =\frac{\mu }{\left\|\mu \right\|}tanh\left(\left\|\mu \right\|\right)`$ . Our construction can be easily generalized to various other realizations of the $`SU(1,1)`$ algebra, such as the two mode realization, where the corresponding states are the Pair coherent and Perelomov(Cave-Schumaker) states. We now extend the above procedure to the quadratic algebra, which is the relevant algebra in considering the coherent states of trilinear boson Hamiltonians . The algebra is given by: $$[J_0,J_\pm ]=\pm J_\pm ,[J_+,J_{}]=+\left(2H_01\right)J_03J_0^2\frac{1q^2}{4}+H_0\left(H_0+1\right).$$ (28) where the positive or negative sign of ($`2H_01`$) determines whether the algebra is a quadratic deformation of $`SU\left(2\right)`$ or $`SU(1,1)`$ respectively. In this case, $`f_1\left(J_0\right)`$ $`=`$ $`\left(2H_01\right)J_03J_0^2{\displaystyle \frac{1Q^2}{4}}+H_0\left(H_0+1\right)`$ (29) $`=`$ $`g_1\left(J_0\right)g_1\left(J_01\right),`$ where $$g_1\left(J_0\right)=J_0\left[H_0\left(H_0+1\right)\frac{1Q^2}{4}+\left(H_0\frac{1}{2}\right)\left(J_0+1\right)\right]J_0\left(J_0+1\right)\left(J_0+\frac{1}{2}\right).$$ (30) In this case we have three different vacua, $`h_0,h_0+\frac{q}{2}>`$, $`h_0,h_0\frac{q}{2}>`$ and $`h_0,h_0>`$ where $`h_0`$ is the eigenvalue of the operator $`H_0`$ and q is the eigenvalue of $`Q=b^{}bc^{}c`$. This is a special case of the general quadratic algebra: $$[N_0,N_\pm ]=\pm N_\pm ,[N_+,N_{}]=\pm 2bN_0+aN_0^2+c.$$ (31) In this case, $`f_1\left(N_0\right)=\pm 2bN_0+aN_0^2+c=g_1\left(N_0\right)g_1\left(N_01\right)`$. with $`g_1\left(N_0\right)=\frac{a}{3}N_0\left(N_0+1\right)\left(N_0+\frac{1}{2}\right)+N_0\left(c\pm b\left(N_0+1\right)\right)`$. The representation theory of the quadratic algebra has been studied in the literature,. It shows a rich structure depending on the values of ‘a’. In the non-compact case, i.e, for polynomial deformations of $`SU(1,1)`$, the unitary irreducible representations (UIREP) are either bounded below or above, we can construct the canonical conjugate $`\stackrel{~}{N_+}`$ of $`N_{}`$ such that $`[N_{},\stackrel{~}{N_+}]=1`$. It is given by $`\stackrel{~}{N_+}=N_+F_1(C,N_0)`$, with $$F_1(C,N_0)=\frac{N_0+\delta }{C\left(N_0\right)\frac{a}{3}N_0\left(N_01\right)\left(N_0+\frac{1}{2}\right)N_0\left(c\pm b\left(N_0+1\right)\right)}.$$ (32) As can be seen easily, in the case of the finite dimensional UIREP, $`\stackrel{~}{N_+}`$ is not well defined since $`F_1(C,N_0)`$ diverges on the highest state. The values of $`\delta `$ can be fixed by demanding that the relation, $`[N_{},\stackrel{~}{N_+}]=1`$, holds in the vacuum sector $`v_i>`$, where , $`v_i>`$’s are annihilated by $`N_{}`$. This gives $`N_{}\stackrel{~}{N_+}v_i>=v_i>`$, which leads to $`\left(N_0+\delta \right)v_i>=v_i>`$. The value of the Casimir operator, $`C=N_{}N_++g_1\left(N_0\right)`$, can then be calculated. Hence, the unnormalized coherent state $`\alpha >`$, such that $`N_{}\alpha >=\alpha \alpha >`$ is given by $`e^{\alpha \stackrel{~}{N_+}}v_i>`$. We can define the canonical conjugate of $`N_+`$ by $`[\stackrel{~}{N_+^{}},N_+]=1`$. The other coherent state is $`\gamma >=e^{\gamma N_+}\stackrel{~}{v}_i>`$ , where $`\stackrel{~}{N_+^{}}\|\stackrel{~}{v}_i>=0`$. Depending on whether the UIREP is infinite or finite dimensional, this quadratic algebra can also be mapped onto the $`SU(1,1)`$ and $`SU\left(2\right)`$ algebras, respectively. Leaving aside the commutators not affected by this mapping, one gets, $$[N_+,\overline{N_{}}]=2bN_0;$$ (33) where $`b=1`$ corresponds to the $`SU(1,1)`$ and $`b=1`$ gives the $`SU\left(2\right)`$ algebra. Explicitly, $$\overline{N_{}}=N_{}G_1(C,N_0),$$ (34) and $$G_1(C,N_0)=\frac{\left(N_0^2N_0\right)b+ϵ}{Cg_1\left(N_01\right)},$$ (35) $`ϵ`$ being an arbitrary constant. One can immediately construct CS in the Perelomov sense (see page 73-74 in ref) as $`\|\xi >=Uv_i>`$, where $`U=e^{\eta N_+\eta ^{}\overline{N_{}}}`$, with $`\xi =\frac{\eta }{\left\|\eta \right\|}tanh\left(\left\|\eta \right\|\right)`$. For the compact case, the CS are analogous to the spin and atomic coherent states. The cubic algebra, which is also popularly known as the Higgs algebra in the literature, appears in the study of the Coulomb problem in a curved space and in quantum optics for quadrilinear boson Hamiltonians. The generators satisfy, $$[M_0,M_\pm ]=\pm M_\pm ,[M_+,M_{}]=2cM_0+\mathrm{\hspace{0.17em}4}hM_0^3,$$ (36) where, $`f_2\left(M_0\right)=2cM_0+\mathrm{\hspace{0.17em}4}hM_0^3=g_2\left(M_0\right)g_2\left(M_01\right)`$, and $$g_2\left(M_0\right)=cM_0\left(M_0+1\right)+hM_0^2\left(M_0+1\right)^2.$$ (37) Analysis of its representation theory yields a variety of UIREP’s, both finite and infinite dimensional, depending on the values of the parameters $`c`$ and $`h`$ . In the non-compact case the canonical conjugate is given by, $$\stackrel{~}{M_+}=M_+F_2(C,M_0),$$ (38) where, $$F_2(C,M_0)=\frac{M_0+\zeta }{CcM_0\left(M_0+1\right)hM_0^2\left(M_0+1\right)^2}.$$ (39) As before, the annihilation operator eigenstate is given by $$\rho >_i=e^{\rho \stackrel{~}{M_+}}p_i>,$$ (40) where, $`p_i>`$ are the states annihilated by $`M_{}`$. Like the previous cases, the dual algebra yields another coherent state. This algebra can also be mapped in to $`SU(1,1)`$ and $`SU\left(2\right)`$ algebras, as has been done for the quadratic case: $$[M_+,\overline{M_{}}]=2dM_0,$$ (41) where, $`d=1`$ and $`d=1`$ correspond to the $`SU(1,1)`$ and $`SU\left(2\right)`$ algebras respectively. Here, $$\overline{M_{}}=M_{}G_2(C,M_0),$$ (42) where, $$G_2(C,M_0)=\frac{\left(M_0^2M_0\right)d+\sigma }{Cg_2\left(M_01\right)},$$ (43) $`\sigma `$ being a constant. The coherent state in the Perelomov sense is then $`\zeta >=Up_i>`$, where, $`U=e^{\zeta M_+\zeta ^{}\overline{M_{}}}`$. In earlier works on non-linear algebras, the generators of the deformed algebra have been written in terms of the undeformed ones. However, in our approach the undeformed $`SU(1,1)`$ and $`SU\left(2\right)`$ generators are constructed from the deformed generators. ## 3 Explicit Construction of the Coherent States for Physical Application. We now give an outline of the method of explicit construction of coherent states for general multiphoton processes for which the generators satisfy the algebra : $`[N_0,N_\pm ]=\pm N_\pm `$ and $`\left[N_+N_{}\right]=g\left(N_0\right)g\left(N_01\right)`$ The action on eigenstates of $`N_0`$ is given by $$N_0j,m>=\left(j+m\right)j,m>$$ (44) $$N_+j,m>=\sqrt{C\left(j\right)g\left(j+m\right)}j,m+1>$$ (45) $$N_{}j,m>=\sqrt{C\left(j\right)g\left(j+m1\right)}j,m1>$$ (46) where $`C\left(j\right)=g\left(j1\right)`$. Depending on the order of the polynomial algebra n, there will be n+1 degenerate states annihilated by $`N_{}`$. We denote these as $`\|j,0>_i`$. For each, the vaule of $`\delta =\delta _i`$ is appropriately chosen as shown earlier. The coherent state is given by $`\alpha >`$ $`=`$ $`Ae^{\alpha \stackrel{~}{N_+}}j,0>_i`$ (47) $`=`$ $`A{\displaystyle \underset{n}{}}{\displaystyle \frac{\alpha ^n}{n!}}\left(N_+\right)^n{\displaystyle \frac{N_0+\delta _i}{g\left(j1\right)g\left(N_0\right)}}\mathrm{}{\displaystyle \frac{N_0+n1+\delta _i}{g\left(j1\right)g\left(N_0+n1\right)}}j,0>`$ $`=`$ $`A{\displaystyle \underset{n}{}}\alpha ^n{\displaystyle \frac{1}{\sqrt{\left(g\left(j1\right)g\left(j\right)\right)\mathrm{}\left(g\left(j1\right)g\left(j+n1\right)\right)}}}j,n>`$ ‘A’ being the normalization constant. A discussion of coherent states is incomplete without showing that these states do give a resolution of the identity and that they are overcomplete. From the resolution of the identity we have: $$𝑑\sigma (\alpha ^{},\alpha )\|\alpha \alpha \|=\mathrm{𝟏}$$ (48) Within the polar decomposition ansatz $$d\sigma (\alpha ^{},\alpha )=\sigma \left(r\right)d\theta rdr$$ (49) with $`r=\left\|\alpha \right\|`$ and an yet unknown positive density $`\sigma `$ which provides the measure. For the general case we have: $$2\pi _0^{\mathrm{}}dr\sigma \left(r\right)r^{2n+1}=A(g(j1)g\left(j\right))\mathrm{}..(g(j1)g(j+n1))$$ (50) For the various cases the substitution of the explicit value of g(j) then reduces the expression on the R.H.S to a rational function of Gamma Functions and the measure $`\sigma `$ can be found by an inverse Mellin transform. For the general case the measure is a Meijer’s G-function. The fact that these states are overcomplete (i.e; $`<\beta \|\alpha >0`$) can be shown for explicit examples. This we shall later show in the case of a quadratic algebra. We now construct the state explicitly for purposes of application. First we show that, this method indeed, gives us the well known SU(1,1) Barut-Girardello (pair coherent) states for SU(1,1) in the familair form . The action on Hilbert Space of the generators is given in the original BG representation by: $$K_0\varphi >=\left(\varphi +m\right)\varphi ,m>,$$ (51) $$K_+\varphi ,m>=\frac{1}{\sqrt{2}}\sqrt{\left(m+1\right)\left(2\varphi +m\right)}\varphi ,m+1>,$$ (52) $$K_{}\varphi ,m>=\frac{1}{\sqrt{2}}\sqrt{m\left(2\varphi +m1\right)}\varphi ,m1>.$$ (53) There are two vacuua annihilated by $`K_{}`$, they correspond to $`\varphi ,0>`$ and $`\varphi ,2\varphi +1>`$. The coherent state $`\alpha >`$ constructed on the vacuum $`\varphi ,0>`$ gives us, $`\delta =\varphi +1`$, so that $`\left(N_0+\delta \right)\varphi ,0>=\varphi ,0>`$ and the resulting coherent state is: $$\alpha >=e^{\alpha \stackrel{~}{K_+}}\varphi ,0>$$ (54) where, $`[K_{},\stackrel{~}{K_+}]=1`$ and $`\stackrel{~}{K_+}=K_+F(C,K_0)`$ with $$F(C,K_0)=\frac{K_0+\delta }{Cg\left(K_0\right)}=\frac{K_0+\delta }{C+\frac{1}{2}K_0\left(K_0+1\right)}.$$ (55) Hence $$\alpha >=\underset{n}{}\frac{\alpha ^n}{n!}\left(K_+F(C,K_0)\right)^n\varphi ,0>=\underset{n}{}\frac{\alpha ^n}{n!}\left(K_+\right)^nF(C,K_0)\mathrm{}F(C,K_0+n1)\varphi ,0>$$ (56) substituting the values of $`F`$ we get: $`\|\alpha >`$ $`=`$ $`A{\displaystyle \underset{n}{}}{\displaystyle \frac{\alpha ^n}{n!}}{\displaystyle \frac{\left(K_0+\delta \right)\left(K_0+\delta +1\right)\mathrm{}.\left(K_0+\delta +n1\right)}{\left(\varphi +\frac{1}{2}\right)\left(\varphi +\frac{3}{2}\right)\mathrm{}\mathrm{}.\left(\varphi +\frac{n1}{2}\right)}}\left(K_+\right)^n\varphi ,0>`$ (57) $`=`$ $`A{\displaystyle \underset{n}{}}\left(2\alpha \right)^n{\displaystyle \frac{\mathrm{\Gamma }\left(2\varphi \right)}{\mathrm{\Gamma }\left(n+1\right)\mathrm{\Gamma }\left(2\varphi +n\right)}}{\displaystyle \frac{\sqrt{n!\left(2\varphi +n1\right)!}}{\left(\sqrt{2}\right)^n\sqrt{\mathrm{\Gamma }\left(2\varphi \right)}}}\varphi ,n>`$ $`=`$ $`A\sqrt{\mathrm{\Gamma }\left(2\varphi \right)}{\displaystyle \underset{n}{}}\left(\sqrt{2}\alpha \right)^n{\displaystyle \frac{1}{\left(\mathrm{\Gamma }\left(n+1\right)\mathrm{\Gamma }\left(2\varphi +n\right)\right)^{\frac{1}{2}}}}\varphi ,n>,`$ which is precisely the well-known state of Barut and Girardello upto the normalization coefficient A. For example for the SU(1,1), $`g\left(j\right)=\frac{1}{2}\left(j\right)\left(j+1\right)`$ then the right hand side becomes in the BG representation $$2\pi _0^{\mathrm{}}𝑑r\sigma \left(r\right)r^{2n+1}=A\frac{\mathrm{\Gamma }\left(n+1\right)\mathrm{\Gamma }\left(2\varphi +n\right)}{\mathrm{\Gamma }\left(2\varphi \right)},$$ (58) where A is a numerical constant and from the inverse Mellin transform, we get $`\sigma \left(r\right)=Ar^{2\varphi +1}K_{\frac{1}{2}+\varphi }\left(2r\right)`$ The second state annihilated by $`K_{}`$ is the state $`\varphi ,2\varphi +1>`$ and this corresponds to $`\delta =\varphi `$. The coherent state is: $$\|\alpha >=A^{}\sqrt{\mathrm{\Gamma }\left(2\varphi \right)}\underset{n}{}\left(\sqrt{2}\alpha \right)^n\frac{1}{\sqrt{\mathrm{\Gamma }\left(n+1\right)\mathrm{\Gamma }\left(2\varphi +n\right)}}\varphi ,2\varphi +1+n>$$ (59) The third state given by Eq. (29) is : $$\|\gamma >=B\left(\gamma \right)\underset{n}{}\left(\frac{\gamma }{\sqrt{2}}\right)^n\sqrt{\frac{\mathrm{\Gamma }\left(n2\varphi \right)}{\mathrm{\Gamma }\left(n+1\right)\mathrm{\Gamma }\left(2\varphi \right)}}\varphi ,n>.$$ (60) This is the state constructed by Perelomov , upto a normalisation constant $`B\left(\gamma \right)`$. For the quadratic case, we take an illustrative algebra relevant to the trilinear boson cases described in the introduction. For convenience we rewrite the three boson algebra as: $`[N_0,N_\pm ]=N_\pm `$ and $`\left[N_+N_{}\right]=3N_0^2+4ϵN_0ϵ^2`$ with $`ϵ=2H_01`$. We define $`n=\left(h_0+j_0\right)`$, where $`h_0`$ and $`j_0`$ are the quantum numbers associated with $`H_0`$ and $`N_0`$ respectively. The state $`\|n>`$ corresponds to the state $`\|h_0,h_0+j_0>=\|ϵ,n>`$ and the three states annihilated by $`N_{}`$ are given by $`\|ϵ,0>,\|ϵ,ϵ\frac{1}{2}>,\|ϵ,ϵ+\frac{1}{2}>`$. The action of the operators on eigenfunctions of $`N_0`$ is given by: $$N_0ϵ,n>=\left(n\right)ϵ,n>,$$ (61) $$N_+ϵ,n>=\sqrt{\left(n+\frac{3}{2}ϵ\right)\left(n+1\right)\left(n+\frac{1}{2}ϵ\right)}ϵ,n+1>,$$ (62) $$N_{}ϵ,n>=\sqrt{\left(n\frac{1}{2}ϵ\right)n\left(n+\frac{1}{2}ϵ\right)}ϵ,n1>.$$ (63) We give the explicit construction of the coherent state for the case $`\|v_i>=\|ϵ,0>`$ , for which $`\delta =1`$. Suitable choices of $`\delta `$ will give the other two coherent states. Here $`g\left(N_0\right)=\left(N_0+\frac{3}{2}ϵ\right)\left(N_0\right)\left(N_0+\frac{1}{2}ϵ\right)`$ and $`g\left(N_0\right)g\left(N_01\right)=3N_0^24ϵN_0+ϵ^2`$. From our construction the CS is: $$\alpha >=e^{\alpha \stackrel{~}{N_+}}ϵ,0>=\underset{n}{}\frac{\alpha ^n}{n!}\left(\stackrel{~}{N_+}\right)^nϵ,0>$$ (64) Thus: $$\|\alpha >=A\underset{n}{}\frac{\alpha ^n}{n!}\left(N_+F(N_0,C)\right)^n0>$$ (65) Constructing the $`F^{}s`$ from $`g\left(N_0\right)`$ we get: $`\|\alpha >`$ $`=`$ $`A{\displaystyle \underset{n}{}}{\displaystyle \frac{\alpha ^n}{n!}}\left(N_+\right)^nF\left(N_0\right)F\left(N_0+1\right)\mathrm{}F\left(N_0+n1\right)0>`$ (66) $`=`$ $`A{\displaystyle \underset{n}{}}{\displaystyle \frac{\alpha ^n}{n!}}\left(N_+\right)^n{\displaystyle \frac{N_0+\delta }{\left(N_0ϵ\right)\left(N_0\right)\left(N_0+1ϵ\right)}}\mathrm{}{\displaystyle \frac{N_0+n1+\delta }{\left(N_0+n1ϵ\right)\left(N_0+n\right)\left(N_0+nϵ\right)}}ϵ,0>`$ $`=`$ $`A{\displaystyle \underset{n}{}}\alpha ^n{\displaystyle \frac{\left(\frac{1}{2}ϵ\right)!\left(\frac{1}{2}ϵ\right)!}{\left(n\frac{1}{2}ϵ\right)!n!\left(n+\frac{1}{2}ϵ\right)!}}\left(N_+\right)^nϵ,0>`$ $`=`$ $`A\sqrt{\mathrm{\Gamma }\left({\displaystyle \frac{1}{2}}ϵ\right)\mathrm{\Gamma }\left({\displaystyle \frac{3}{2}}ϵ\right)}{\displaystyle \underset{n}{}}{\displaystyle \frac{\alpha ^n}{\sqrt{\mathrm{\Gamma }\left(n+\frac{1}{2}ϵ\right)\mathrm{\Gamma }\left(n+1\right)\mathrm{\Gamma }\left(n+\frac{3}{2}ϵ\right)}}}ϵ,n>,`$ $`A`$ is the normalization coefficient, which can be easily determined to be $`\frac{1}{(^0F_2(\frac{1}{2}ϵ,\frac{3}{2}ϵ,\|\alpha \|^2))^{\frac{1}{2}}}`$. These set of states can be shown to be overcomplete: $$\left\|<\beta \|\alpha >\right\|^2=\frac{{}_{}{}^{0}F_{2}^{}(\frac{1}{2}ϵ,\frac{3}{2}ϵ,\alpha \beta ^{})}{(^0F_2(\frac{1}{2}ϵ,\frac{3}{2}ϵ,\|\alpha \|^2))^{\frac{1}{2}}\left)(^0F_2(\frac{1}{2}ϵ,\frac{3}{2}ϵ,\|\beta \|^2)\right)^{\frac{1}{2}}}.$$ (67) The completeness relation is given by $$_0^{\mathrm{}}𝑑r\sigma \left(r\right)r^{2n+1}=\mathrm{\Gamma }\left(n+1\right)\frac{\mathrm{\Gamma }\left(\frac{1}{2}ϵ+n\right)\mathrm{\Gamma }\left(\frac{3}{2}ϵ+n\right)}{\mathrm{\Gamma }\left(\frac{1}{2}ϵ\right)\mathrm{\Gamma }\left(\frac{3}{2}ϵ\right)}$$ (68) and $`\sigma \left(r\right)`$ can be determined to be a confluent hypergeometric function from the inverse Mellin transformation formula. The resolution of the identity can thus be obtained. $$\sigma \left(r\right)=\frac{1}{\mathrm{\Gamma }(\frac{1}{2}ϵ)(\mathrm{\Gamma }(\frac{3}{2}ϵ)}G_{\mathrm{0\hspace{0.17em}3}}^{\mathrm{3\hspace{0.17em}0}}\left(r_{0,\frac{1}{2}ϵ,\frac{1}{2}ϵ}^0\right),$$ (69) Where $`G_{\mathrm{0\hspace{0.17em}3}}^{\mathrm{3\hspace{0.17em}0}}\left(x\right)`$ is a Meijer’s G-function. The other two coherent states based on the two vacuua, $`\|ϵ,ϵ+\frac{1}{2}>`$ and $`\|ϵ,ϵ\frac{1}{2}>`$ can similarly be constructed by chosing $`\delta =\frac{3}{2}ϵ`$ and $`\delta =\frac{1}{2}ϵ`$. The state corresponding to the Perelomov state is : $$\|\gamma >=B^{}\left(\gamma \right)\underset{n}{}\left(\gamma \right)^n\sqrt{\frac{\mathrm{\Gamma }\left(n+\frac{3}{2}ϵ\right)\mathrm{\Gamma }\left(n+\frac{1}{2}ϵ\right)}{\mathrm{\Gamma }\left(n+1\right)\mathrm{\Gamma }\left(\frac{1}{2}ϵ\right)\mathrm{\Gamma }\left(\frac{3}{2}ϵ\right)}}\varphi ,n>.$$ (70) The normalisation constant can be calculated easily and using a method similar to the one used for Eq.(68), the overcompleteness of these states and the resolution of the identity can also be easily obtained. ## 4 Conclusion To conclude, we have found a general method for constructing the coherent states for various polynomially deformed algebras for quantum optical systems, whose dynamics are governed by multilinear boson Hamiltonians. Since our method is algebraic and relies on the group structure of well-known algebras, the precise nature of the non-classical behaviour of these CS can be easily inferred from our construction. It will be of particular interest to see the time development of the system and the role of the deformation parameters in the physical system described in the text. For a system initially in a coherent state, it is fairly straightforward to calculate the time evolution of the system exactly using the methods of reference . Since many of these algebras are related to quantum mechanical problems with non-quadratic, non-linear Hamiltonians, a detailed study of the properties of the CS associated with non-linear and deformed algebras is of physical relevance . This is the subject of our current and future work . The authors take the pleasure to thank Prof. S. Chaturvedi and Prof. C. Mukku for stimulating conversations. VSK acknowledges useful discussions with Mr. N. Gurappa.
warning/0003/hep-ph0003226.html	ar5iv	text	# Northwestern University: N.U.H.E.P. Report No. 708 March, 2000 Breaking the Barriers—Uniting Accelerator and Cosmic Ray p-p Cross Sections ## 1 Introduction This communication is divided into three sections. First, we show that the data on the total cross section, the slope parameter $`B`$ of the elastic differential cross section, and the ratio of the real to imaginary part of the forward scattering amplitude $`\rho `$ for $`pp`$ and $`\overline{p}p`$ interactions can be nicely described by a model where high energy cross sections grow with energy as a consequence of the increasing number of soft partons populating the colliding particles ,. The differential cross sections for the Tevatron and LHC are predicted. Next, we verify the model by showing that the known experimental data on $`\gamma p`$ and $`\gamma \gamma `$ interactions can be derived from our $`pp`$ and $`\overline{p}p`$ forward scattering amplitudes using vector meson dominance (VMD) and the additive quark model. Finally, we use the high energy predictions of our QCD-inspired parameterization of accelerator data on forward proton-proton and antiproton-proton scattering amplitudes, along with Glauber theory, to predict proton–air cross sections at energies near $`\sqrt{s}`$ 30 TeV. All cross sections will be computed in an eikonal formalism guaranteeing unitarity throughout: $`\sigma _{tot}(s)`$ $`=`$ $`2{\displaystyle \left\{1e^{\chi __I(b,s)}\mathrm{cos}[\chi __R(b,s)]\right\}d^2\stackrel{}{b}}.`$ (1) Here, $`\chi `$ is the complex eikonal ($`\chi =\chi __R+i\chi __I`$), and $`b`$ is the impact parameter. The even eikonal profile function $`\chi ^{even}`$ receives contributions from quark-quark, quark-gluon and gluon-gluon interactions, and therefore $`\chi ^{even}(s,b)`$ $`=`$ $`\chi _{qq}(s,b)+\chi _{qg}(s,b)+\chi _{gg}(s,b)`$ (2) $`=`$ $`i\left[\sigma _{qq}(s)W(b;\mu _{qq})+\sigma _{qg}(s)W(b;\sqrt{\mu _{qq}\mu _{gg}})+\sigma _{gg}(s)W(b;\mu _{gg})\right],`$ where $`\sigma _{ij}`$ are the cross sections of the colliding partons, and $`W(b;\mu )`$ their overlap function in impact parameter space, parameterized as the Fourier transform of a dipole form factor. The impact parameter space distribution function $$W(b;\mu )=\frac{\mu ^2}{96\pi }(\mu b)^3K_3(\mu b)$$ (3) is normalized so that $`W(b;\mu )d^2\stackrel{}{b}=1`$. As a consequence of both factorization and the normalization chosen for the $`W(b;\mu )`$, it should be noted that $$\chi ^{even}(s,b)d^2\stackrel{}{b}=i\left[\sigma _{gg}(s)+\sigma _{qg}(s)+\sigma _{qq}(s)\right],$$ (4) so that $`\sigma _{tot}^{even}(s)=2\mathrm{Im}\left\{i\left[\sigma _{gg}(s)+\sigma _{qg}(s)+\sigma _{qq}(s)\right]\right\}`$, for small $`\chi `$. This formalism is identical to the one used in “mini-jet” models , as well as in simulation programs for minimum-bias hadronic interactions such as PYTHIA and SIBYLL. In this model hadrons asymptotically evolve into black disks of partons. The rising cross section, asymptotically associated with gluon-gluon interactions, is simply parameterized by a normalization, an energy scale, and two parameters: $`\mu _{gg}`$ which describes the “area” occupied by gluons in the colliding hadrons, and $`J(=1+ϵ)`$. Here, $`J`$ is defined via the gluonic structure function of the proton, which is assumed to behave as $`1/x^J`$ for small x. It therefore controls the soft gluon content of the proton. The introduction of the quark-quark and quark-gluon terms allows us to adequately parameterize the data at all energies, since the “size” of quarks and gluons in the proton can be different. In the present context, this model represents a convenient parameterization of the $`pp`$ and $`p\overline{p}`$ forward scattering amplitude. The photoproduction cross sections are calculated from this parameterization assuming vector meson dominance and the additive quark model. For the probability that the photon interacts as a hadron ($`P_{had}`$), we use the value $`P_{had}=1/240`$ which can be derived from vector meson dominance. Our results show that its value is indeed independent of energy. It is, however, uncertain by 20% because it depends on whether we relate photoproduction to $`\pi `$-nucleon or nucleon-nucleon data (In other words, $`\pi N`$ and $`NN`$ total cross sections only satisfy the additive quark model to this accuracy). Subsequently, following reference , we obtain $`\gamma p`$ cross sections from the assumption that, in the spirit of VMD, the photon is a 2 quark state in contrast with the proton which is a 3 quark state. The $`\gamma p`$ total cross section is obtained from the even eikonal for $`pp`$ and $`\overline{p}p`$ by the substitutions $`\sigma _{ij}\frac{2}{3}\sigma _{ij}`$ and $`\mu _i\sqrt{\frac{3}{2}}\mu _i`$. We will thus produce a parameter-free description of the total photoproduction cross section, the phase of the forward scattering amplitude and the forward slope for $`\gamma pVp`$, where $`V=\rho ,\omega ,\varphi `$. Interestingly, our results on the phase of $`VpVp`$ are in complete agreement with the values derived from Compton scattering results ($`\gamma +p\gamma +p`$) using dispersion relations. We also calculate the total elastic and differential cross sections for $`\gamma pVp`$. This wealth of data is accommodated without discrepancy. The $`\gamma \gamma `$ cross sections are derived following the same procedure. We now substitute $`\sigma _{ij}\frac{4}{9}\sigma _{ij}`$ and $`\mu _i\frac{3}{2}\mu _i`$ into the nucleon-nucleon even eikonal, and predict the total cross section and the differential cross sections for all reactions $`\gamma \gamma V_iV_j`$ at a variety of energies, where $`V=\rho ,\omega ,\varphi `$. The high energy $`\gamma \gamma `$ total cross section have been measured by two experiments at LEP. While these measurements yield new information on its high energy behavior at center-of-mass energies in excess of $`\sqrt{s}=15`$ GeV, they may represent the last opportunity to measure the $`\gamma \gamma `$ cross section, and the two data sets appear to disagree. However, it has been argued that the original data are consistent within the errors and that the observed disagreements are due to two different Monte Carlo’s used to extract the quoted values. We here point out that our analysis nicely accommodates the L3 result . Our model approximately satisfies the factorization theorem, $`\sigma _{pp}/\sigma _{\gamma p}=\sigma _{\gamma p}/\sigma _{\gamma \gamma }`$, because of its small eikonal. The OPAL data do not satisfy it. In fact, no model incorporating the additive quark model and factorization can accommodate the OPAL data. VMD and factorization are sufficient to prevent one from adjusting $`P_{had}`$, or any other parameters, to change this conclusion. ## 2 High energy proton-proton and proton-antiproton scattering In this section we discuss our QCD-inspired parameterization of the forward amplitudes. To determine its parameters, we fit all high energy forward $`\overline{p}p`$ and $`pp`$ scattering data above 15 GeV, for the total cross section ($`\sigma _{tot}`$), the ratio of the real to the imaginary part of the forward scattering amplitude ($`\rho `$), and the logarithmic slope of the differential elastic scattering cross section in the forward direction ($`B`$). Then, we compare the experimental data for the elastic scattering cross section and for the differential elastic scattering with our results. Finally, a prediction is made for the differential elastic scattering at the LHC. Our QCD-inspired parameterization satisfies crossing symmetry, i.e., it is either even or odd under the transformation $`EE`$, where $`E`$ is the laboratory energy. This allows us to simultaneously describe $`\overline{p}p`$ and $`pp`$ scattering. It also satisfies analyticity, and unitarity because of the eikonal formalism. Since the total cross section asymptotically rises as $`\mathrm{log}^2s`$, our QCD-inspired parameterization complies with the Froissart bound. The eikonal formalism for calculating $`\sigma _{tot}`$, $`\rho `$ and $`B`$, along with details on the analyticity, the Froissart bound, and the QCD-inspired eikonal are given in ref. . In all 11 parameters are used. The low energy region, where the differences between $`\overline{p}p`$ and $`pp`$ scattering are substantial, largely determines the 7 parameters necessary to fit the odd eikonal and the quark-quark and quark-gluon contribution to the even eikonal. Thus, they largely decouple from the high energy behavior. Hence, for $`\sqrt{s}25`$ GeV, where the difference between $`\overline{p}p`$ and $`pp`$ scattering becomes small, only 4 parameters describe all data. We fit all the highest energy cross section data (E710 , CDF and the unpublished Tevatron value ), which anchor the upper end of our cross section curves. The results of the fit are shown in Fig. 1. Data for $`\rho `$ values and $`B`$ are confronted with our model in Figs. 2 and 3. It can be seen from those figures that we obtain a satisfactory description of all 3 quantities, for both $`\overline{p}p`$ and $`pp`$ scattering. The $`\chi ^2`$ of the fit is reasonably good (considering the large spread in some of the experimental data, as well as the discrepancies in the highest energy cross sections), giving a $`\chi ^2/d.f.=1.66`$, for 75 degrees of freedom. The model splits the difference between the measurements of the total cross section at $`\sqrt{s}=1800`$ GeV (see Fig. 1). From Fig. 2, we note that the fit to $`\rho `$ is anchored at $`\sqrt{s}=550`$ GeV by the very accurate measurement of UA4/2 and passes through the E710 point . The statistical uncertainty of the fitted parameters is such that at 25 GeV the cross section predictions are statistically uncertain to $`1.3`$%, at 500 GeV are uncertain to $`1.6`$%, and at 2000 GeV are uncertain to $`2.5`$%. In Fig. 4 we show the prediction for the elastic cross section along with the data for both $`\overline{p}p`$ and $`pp`$. The agreement is excellent. We note that $`\sigma _{elastic}`$ is rising more sharply with energy than the total cross section $`\sigma _{tot}`$. Comparing Fig. 1 with Fig. 4, we see that the ratio of the elastic to total cross section is rising with energy. The ratio is, of course, bounded by the value for the black disk , i.e., 0.5, as the energy goes to infinity. Having fixed all parameters specifying our eikonal, we calculate $`d\sigma /dt`$, for various values of $`\sqrt{s}`$. The differential cross section at the Tevatron ($`\sqrt{s}=1800`$ GeV) is shown in Fig. 5 along with E710 data. The agreement over 4 decades is striking. Our prediction for the differential cross section at $`\sqrt{s}=14`$ TeV, the energy of the LHC, is plotted in Fig. 6. In particular, at small $`\|t\|`$, we predict that the curvature parameter $`C`$ ($`d\sigma /dte^{Bt+Ct^2}`$ for small $`t`$; see ref. for details) is negative. For energies much lower than 1800 GeV, the observed curvature has been measured as positive. For 1800 GeV, we see from Fig. 5 that the curvature parameter $`C`$ is compatible with being zero. Block and Cahn have pointed out that the curvature is predicted to go through zero near the Tevatron energy and that it should become negative thereafter. Asymptotically the proton approaches a black disk. Its curvature is always negative , $`C=R^4/192`$, where $`R`$ is the radius of the disk. Thus, the curvature has to pass through zero as the energy increases. ‘Asymptopia’ is the energy region (energies much larger than the Tevatron) where the scattering approaches that of a sharp disk. With the parameters we obtained from our fit, the total cross section at the LHC (14 TeV) is predicted to be $`\sigma _{tot}=108.0\pm 3.4`$ mb, where the error is due to the statistical errors of the fitting parameters. ## 3 Photon-proton reactions We assume that the photon behaves like a two quark system when it interacts strongly. We therefore obtain $`\gamma p`$ scattering amplitudes by performing the substitutions $`\sigma _{ij}\frac{2}{3}\sigma _{ij}`$ and $`\mu _i\sqrt{\frac{3}{2}}\mu _i`$ in the even eikonal for nucleon–nucleon scattering, so that $$\chi ^{\gamma p}(s,b)=i\left[\frac{2}{3}\sigma _{qq}(s)W(b;\sqrt{\frac{3}{2}}\mu _{qq})+\frac{2}{3}\sigma _{qg}(s)W(b;\sqrt{\frac{3}{2}\mu _{qq}\mu _{gg}})+\frac{2}{3}\sigma _{gg}(s)W(b;\sqrt{\frac{3}{2}}\mu _{gg})\right].$$ (5) Using vector dominance, the photon-proton total cross section is then written as $$\sigma _{tot}^{\gamma p}(s)=2P_{had}\left\{1e^{\chi _I^{\gamma p}(b,s)}\mathrm{cos}[\chi _R^{\gamma p}(b,s)]\right\}d^2\stackrel{}{b},$$ (6) where $`P_{had}`$ is the probability that a photon interacts as a hadron. We use the value $`P_{had}=1/240`$. This value is found by normalizing the total $`\gamma p`$ cross section to the low energy data, and is very close to that derived from vector dominance, 1/249. Using $`f_\rho ^2/4\pi =2.2`$, $`f_\omega ^2/4\pi =23.6`$ and $`f_\varphi ^2/4\pi =18.4`$, we find $`\mathrm{\Sigma }_V(4\pi \alpha /f_V^2)=1/249`$, where $`V=\rho ,\omega ,\varphi `$ (see Table XXXV, pag. 393 of Ref. ). With all eikonal parameters fixed by the nucleon-nucleon data, we can now calculate $`\sigma _{tot}^{\gamma p}(s)`$. The result is shown in Fig. 7. It reproduces the rising cross section for $`\gamma p`$, using the parameters fixed by nucleon-nucleon scattering. This prediction only uses the 9 parameters of the even eikonal, of which but 4 are important in the upper energy region. The accuracy of our predictions are $`1.5\%`$, from the statistical uncertainty in our eikonal parameters. We next consider the ‘elastic’ scatterings $`\gamma +p`$ $``$ $`\rho _{virtual}+p\rho +p,`$ $`\gamma +p`$ $``$ $`\omega _{virtual}+p\omega +p,`$ $`\gamma +p`$ $``$ $`\varphi _{virtual}+p\varphi +p.`$ (7) Here the photon virtually transforms into a vector meson which elastically scatters off of the proton. The strengths of these reactions is $`𝒪(\alpha )`$ times a strong interaction cross section. The true elastic cross section is given by Compton scattering on the proton, $`\gamma +p\gamma +p`$, which we can visualize as $`\gamma +p`$ $``$ $`\rho _{virtual}+p\rho +p\gamma +p,`$ $`\gamma +p`$ $``$ $`\omega _{virtual}+p\omega +p\gamma +p,`$ $`\gamma +p`$ $``$ $`\varphi _{virtual}+p\varphi +p\gamma +p.`$ (8) It is clearly $`𝒪(\alpha ^2)`$ times a strong interaction cross section, and hence is much smaller than ‘elastic’ scattering of Eq. (7). Thus, we justify the use of Eq. (6) to calculate the total cross section, since only reactions with a photon in the final state are neglected. We evaluate $`\rho `$ and the slope $`B`$ for the ‘elastic’ scattering expressed in Eq. (7), with $`\rho `$ and $`B`$ being the same for all 3 reactions. The dependence of $`\rho `$ with the energy is shown in Fig. 8. Damashek and Gilman have calculated the $`\rho `$ value for Compton scattering on the proton using dispersion relations, i.e., the true elastic scattering reaction for photon-proton scattering. We compare this calculation, the dotted line in Fig. 8, with our prediction of $`\rho `$ (the solid line). The agreement is so close that we had to move the two curves apart so that they may be viewed more clearly. In Fig. 9 we show our results for the slope $`B`$ as a function of the energy. The available experimental data for ‘elastic’ $`\rho p`$ and $`\omega p`$ final states are also plotted. Again, the agreement of theory and experiment is very good. To calculate the elastic cross sections $`\sigma _{elastic}^{Vp}`$ and differential cross sections $`d\sigma ^{Vp}/dt`$ as a function of energy, we use $$\sigma _{elastic}^{Vp}(s)=P_{had}^{Vp}\left\|1e^{i\chi ^{\gamma p}(b,s)}\right\|^2d^2\stackrel{}{b},$$ (9) where $`P_{had}^{Vp}`$ is the appropriate probability for a photon to turn into $`V`$, with $`V=\rho ,\omega `$ or $`\varphi `$. The differential scattering cross section is given by $$\frac{d\sigma ^{Vp}}{dt}(s,t)=\frac{P_{had}^{Vp}}{4\pi }\left\|J_0(qb)(1e^{i\chi ^{\gamma p}(b,s)})d^2\stackrel{}{b}\right\|^2,$$ (10) where $`t=q^2`$. Since we normalize our data to the cross section found with $`\chi ^{\gamma p}`$, and not to $`(\sigma _{tot}^{\pi ^+}+\sigma _{tot}^\pi ^{})/2`$, we must multiply all $`f_V^2/4\pi `$ by 1.65. Hence, our effective couplings are $`f_{\rho \mathrm{eff}}^2/4\pi =3.6`$, $`f_{\omega \mathrm{eff}}^2/4\pi =38.9`$, and $`f_{\varphi \mathrm{eff}}^2/4\pi =30.4`$. Our evaluation of the ‘elastic’ cross section for the reactions $`\gamma +p\rho ^0+p`$ and $`\gamma +p\omega ^0+p`$ are shown in Figs. 10 and 11, respectively. The differential cross section, $`d\sigma /dt`$, for the ‘elastic’ reactions $`\gamma +p\rho ^0+p`$, $`\gamma +p\omega ^0+p`$ and $`\gamma +p\varphi ^0+p`$ are plotted in Figs. 12, 13, and 14, respectively. The agreement, in absolute normalization and shape, of our results for all three light vector mesons with the experimental data for all available energies reinforces our confidence in the model. ## 4 Photon-Photon Interactions In this section, we consider $`\gamma \gamma `$ interactions. As it was done for $`\gamma p`$ interactions, we will start from the eikonal $`\chi ^{\gamma p}(s,b)`$ and multiply every cross section by 2/3 and multiply each $`\mu `$ by $`\sqrt{3/2}`$. Therefore, $$\chi ^{\gamma \gamma }(s,b)=i\left[\frac{4}{9}\sigma _{qq}(s)W(b;\frac{3}{2}\mu _{qq})+\frac{4}{9}\sigma _{qg}(s)W(b;\frac{3}{2}\sqrt{\mu _{qq}\mu _{gg}})+\frac{4}{9}\sigma _{gg}(s)W(b;\frac{3}{2}\mu _{gg})\right].$$ (11) Using vector dominance we obtain, $$\sigma _{tot}^{\gamma \gamma }(s)=2P_{had}^2\left\{1e^{\chi _I^{\gamma \gamma }(b,s)}\mathrm{cos}[\chi _R^{\gamma p}(b,s)]\right\}d^2\stackrel{}{b},$$ (12) where $`P_{had}=1/240`$ is the probability that a photon will interact as a hadron. In Fig. 15 we plot our results for $`\sigma _{tot}^{\gamma \gamma }(s)`$ as a function of the energy, and compare it to the various sets of experimental data. We note that our prediction fits the L3 data, but doesn’t fit the OPAL results. ## 5 Proton-air cross sections Cosmic ray experiments measure the penetration in the atmosphere of particles with energies in excess of those accelerated by existing machines—interestingly, their energy range covers the energy of the Large Hadron Collider (LHC) and extends beyond it. However, extracting proton–proton cross sections from cosmic ray observations is far from straightforward . By a variety of experimental techniques, cosmic ray experiments map the atmospheric depth at which cosmic ray initiated showers develop. The measured shower attenuation length ($`\mathrm{\Lambda }_m`$) is not only sensitive to the interaction length of the protons in the atmosphere ($`\lambda _{p\mathrm{air}}`$), with $$\mathrm{\Lambda }_m=k\lambda _{p\mathrm{air}}=k\frac{13.5m_p}{\sigma _{p\mathrm{air}}^{\mathrm{inel}}},$$ (13) but also depends on the rate at which the energy of the primary proton is dissipated into electromagnetic shower energy observed in the experiment. The latter effect is parameterized in Eq. (13) by the parameter $`k`$; $`m_p`$ is the proton mass and $`\sigma _{p\mathrm{air}}^{\mathrm{inel}}`$ the inelastic proton-air cross section. The value of $`k`$ depends on the inclusive particle production cross section in nucleon and meson interactions on the light nuclear target of the atmosphere and its energy dependence. We here ignored the fact that particles in the cosmic ray ”beam” may be nuclei, not just protons. Experiments allow for this by omitting from their analysis showers which dissipate their energy high in the atmosphere, a signature that the initial energy is distributed over the constituents of a nucleus. The extraction of the pp cross section from the cosmic ray data is a two step process. First, one calculates the $`p`$-air total cross section from the measured inelastic cross section $$\sigma _{p\mathrm{air}}^{\mathrm{inel}}=\sigma _{p\mathrm{air}}\sigma _{p\mathrm{air}}^{\mathrm{el}}\sigma _{p\mathrm{air}}^{q\mathrm{el}}.$$ (14) Next, the Glauber method is used to transform the measured value of $`\sigma _{p\mathrm{air}}^{\mathrm{inel}}`$ into a proton–proton total cross section $`\sigma _{pp}`$; all the necessary steps are calculable in the theory. In Eq. (14) the measured cross section for particle production is supplemented with the elastic and quasi-elastic cross section, as calculated by the Glauber theory, to obtain the total cross section $`\sigma _{p\mathrm{air}}`$. The subsequent relation between $`\sigma _{p\mathrm{air}}^{\mathrm{inel}}`$ and $`\sigma _{pp}`$ involves the slope of the forward scattering amplitude for elastic $`pp`$ scattering, $$B=\left[\frac{d}{dt}\left(\mathrm{ln}\frac{d\sigma _{pp}^{\mathrm{el}}}{dt}\right)\right]_{t=0},$$ (15) and is shown in Fig. 16, which plots $`B`$ against $`\sigma _{pp}`$, for 5 curves of different values of $`\sigma _{p\mathrm{air}}^{\mathrm{inel}}`$. This summarizes the reduction procedure from $`\sigma _{p\mathrm{air}}^{\mathrm{inel}}`$ to $`\sigma _{pp}`$ . Also plotted in Fig. 16 is a curve of $`B`$ vs. $`\sigma _{pp}`$ which will be discussed later. A significant drawback of the method is that one needs a model of proton–air interactions to complete the loop between the measured attenuation length $`\mathrm{\Lambda }_m`$ and the cross section $`\sigma _{p\mathrm{air}}^{\mathrm{inel}}`$, i.e., the value of $`k`$ in Eq. (13). We minimize the impact of theory by using our QCD-inspired parameterization of the forward proton–proton and proton–antiproton scattering amplitudes which is analytic, unitary and simultaneously fits all data of $`\sigma _{\mathrm{tot}}`$, $`B`$ and $`\rho `$. Using vector meson dominance and the additive quark models, we have shown that it accommodates a wealth of data on photon-proton and photon-photon interactions without the introduction of new parameters. Because the model is both unitary and analytic, it has high energy predictions that are essentially theory–independent. In particular, it also simultaneously fits $`\sigma _{pp}`$ and $`B`$, forcing a relationship between the two. Specifically, the $`B`$ vs. $`\sigma _{pp}`$ prediction of the model is shown as the dashed curve in Fig. 16. The dot corresponds to our prediction of $`\sigma _{pp}`$ and $`B`$ at $`\sqrt{s}`$ = 30 TeV. It is seen to be slightly below the curve for 490 mb, the lower limit of the Fly’s Eye measurement, which was made at $`\sqrt{s}`$ 30 TeV. In Fig. 17, we have plotted the values of $`\sigma _{pp}`$ vs. $`\sigma _{p\mathrm{air}}^{\mathrm{inel}}`$ that are deduced from the intersections of the $`B`$-$`\sigma _{pp}`$ curve with the $`\sigma _{p\mathrm{air}}^{\mathrm{inel}}`$ curves of Fig. 16. Figure 17 allows the conversion of the measured $`\sigma _{p\mathrm{air}}^{\mathrm{inel}}`$ to $`\sigma _{pp}`$ . Our prediction for the total cross section $`\sigma _{pp}`$ as a function of energy is confronted with all of the accelerator and cosmic ray measurements in Fig. 18. For inclusion in Fig. 18, we have calculated the cosmic ray values of $`\sigma _{pp}`$ from the published experimental values of $`\sigma _{p\mathrm{air}}^{\mathrm{inel}}`$, using the results of Fig. 17. We note the predicted curve is systematically lower than the cosmic ray points, roughly about the level of one standard deviation. It is at this point important to recall Eq. (13) and consider the fact that the extraction of $`\sigma _{p\mathrm{air}}^{\mathrm{inel}}`$ from the measurement of $`\mathrm{\Lambda }_m`$ requires a determination of the parameter $`k`$. The measured depth $`X_{\mathrm{max}}`$ at which a shower reaches maximum development in the atmosphere, which is the basis of the cross section measurement in Ref. , is a combined measure of the depth of the first interaction, which is determined by the inelastic cross section, and of the subsequent shower development, which has to be corrected for. The position of $`X_{\mathrm{max}}`$ also directly affects the rate of shower attenuation with atmospheric depth which is the alternative procedure for extracting $`\sigma _{p\mathrm{air}}^{\mathrm{inel}}`$. The model dependent rate of shower development and its fluctuations are the origin of the deviation of $`k`$ from unity in Eq. (13). Its values range from 1.5 for a model where the inclusive cross section exhibits Feynman scaling, to 1.1 for models with large scaling violations. The comparison between data and experiment in Fig. 18 is further confused by the fact that the AGASA and Fly’s Eye experiments used different values of $`k`$ in the analysis of their data, i.e., AGASA used $`k=1.5`$ and Fly’s Eye used $`k=1.6`$. We therefore decided to match the data to our prediction and extracted a common value for $`k=1.33\pm 0.04`$. This neglects the possibility that $`k`$ may show a weak energy dependence over the range measured. In Fig. 19 we have replotted the high energy cosmic ray data for our prediction of $`\sigma _{p\mathrm{air}}^{\mathrm{inel}}`$ vs. $`\sqrt{s}`$, with the common value of $`1.33`$ obtained from a $`\chi ^2`$ fit. Clearly, we have an excellent fit, with good agreement between AGASA and Fly’s Eye. The analysis gives $`\chi ^2=1.75`$ for 6 degrees of freedom (the low $`\chi ^2`$ is probably due to overestimates of experimental errors). This result for $`k`$ is interesting—it is close to the value of $`1.2`$ obtained using the SIBYLL simulation for inclusive particle production. This represents a consistency check in the sense that our model for forward scattering amplitudes and SIBYLL share the same underlying physics. The increase of the total cross section with energy to a black disk of soft partons is the shadow of increased particle production which is modeled by the production of (mini)-jets in QCD. The difference between the $`k`$ values of 1.20 and 1.33 could be understood because the experimental measurement integrates showers in a relatively wide energy range, which tends to increase the value of $`k`$. In the near term, we look forward to the possibility of repeating this analysis with the higher statistics of the HiRes cosmic ray experiment that is currently in progress and the Auger Observatory. In conclusion, we have successfully united the high energy cross section results ($`\sqrt{s}30`$ TeV) of the cosmic ray measurements with the accelerator cross section measurements, under a common rubric, the QCD-inspired analysis.
warning/0003/astro-ph0003125.html	ar5iv	text	# 1 Specklegram of a binary star, HR4689 obtained at VBT, Kavalur, India. Speckle Imaging: a boon for astronomical observations S. K. Saha Indian Institute of Astrophysics Bangalore - 560 034 e-mail: sks@iiap.ernet.in Abstract The speckle imaging is a photographic technique that resolves objects viewed through severely distorted media. The results are insensitive to the errors caused by apparent size of the isoplanatic patch and the telescope aberrations. In this article, a short descriptions of the atmospheric turbulence and its effect on the flat wavefront from a stellar source is presented; the shortcomings of the conventional long-exposure images in the presence of Earth’s atmosphere are discussed. The advantages of the speckle interferometric technique over conventional imaging are enumerated. The technical details of the method, basic Fourier optics, data analysis procedures are also described. 1. Introduction The atmosphere of earth restricts the resolution of conventional ground based astronomical images to about 1 arc second. This is due to the refractive index variations of the atmosphere through which the light rays reaching the telescope. When a star image is observed through a telescope with high magnification, the observed image structure is usually far from the theoretical pattern. The appearance of image depends strongly on the size of aperture of the telescope. Large telescope helps in gathering more optical energy, as well as in obtaining better angular resolution. The resolution increases with the diameter of the aperture. For example, the diffraction-limit of a 2.34 meter telescope is about 0.05 arc second. But in reality, the image is degraded by factor of 20. With small aperture a random motion of a image is often affect the main effect, whereas, with large aperture spreading and blurring of the image occur. Owing to the diffraction phenomenon, the image of the point source (unresolved stars) cannot be smaller than a limit at the focal plane of the telescope. This phenomenon can be observed in ocean, when regular waves pass through an aperture. It is present in the sound waves, as well as in the electro-magnetic spectrum too starting from gamma rays to radio waves. Speckle interferometric technique (Labeyrie, 1970) yields the diffraction-limited autocorrelation of the object. The diffraction-limited resolution of celestial objects viewed through the Earth’s turbulent atmosphere could be achieved with the largest optical telescope, by post detection processing of a large data set of short-exposure images using Fourier-domain methods. Certain specialized moments of the Fourier transform of a short-exposure image contain diffraction-limited information about the object of interest. In this article, I shall describe a few basic theorems of Fourier optics, which are essential to understand this technique followed by the theory of speckle interferometry. The data processing method to analyze specklegrams of close binary stars obtained with 2.34 meter Vainu Bappu telescope (VBT), Kavalur, India is also discussed. 2. Preamble In order to obtain high angular resolution of an stellar object, Fizeau (1868) had suggested to install a screen with two holes on top of the telescope that produce Young’s fringes at its focal plane as the fringes remain visible in presence of seeing, thus allowing measurements of stellar diameters. Stefan attempted with 1 meter telescope at Observatoire de Marseille but could not notice any significant drop of fringe visibility and opined none of the observed stars approached 0.1 arc-seconds in angular size. About half a century later, Michelson could measure the diameter of the satellites of Jupiter with Fizeau interferometer on top of the Yerkes refractor. With the 100 inch telescope at Mt. Wilson (Anderson, 1920), the angular separation of spectroscopic binary star Capella was measured. To overcome the restrictions of the baseline, Michelson (1920) constructed the stellar interferometer equipped with 4 flat mirrors to fold the beams by installing a 7 meter steel beam on top of the afore-mentioned 100 inch telescope; the supergiant star $`\alpha `$ Orionis were resolved (Michelson and Pease, 1921). Due to the various difficulties, viz., (i) effect of atmospheric turbulence, (ii) variations of refractive index above small sub-apertures of the interferometer, and (iii) mechanical instability, the project was abandoned. The field of optical interferometry lay dormant until it was revitalized by the development of intensity interferometry (Brown and Twiss, 1958). Success in completing the intensity interferometer at radio wavelengths (Brown et al., 1952), in which the signals at the antennae are detected separately and the angular diameter of the source is obtained by measuring correlation of the intensity fluctuations of the signals as a function of antenna separation, Brown and Twiss (1958) demonstrated its potential at optical wavelengths by measuring the angular diameter of Sirius. Subsequent development of this interferometer with a pair of 6.5 meter light collector on a circular railway track spanning 188 meter (Brown et al., 1967), depicted the measurements of 32 southern binary stars with angular resolution limit of 0.5 milliarcseonds (Brown, 1974). The project was abandoned due to lack of photons beyond 2.5 magnitude stars. Meanwhile, Labeyrie (1975) had developed a long baseline interferometer $``$ Interf$`\stackrel{`}{e}`$rom$`\overline{e}`$tre $`\stackrel{`}{a}`$ deux telescope (I2T) $``$ using a pair of 25 cm telescopes at Observatoire de Calern, France. His design combines features of the Michelson and of the radio interferometers. The use of independent telescopes increases the resolving capabilities. In this case, coude beams from both the telescopes arrive at central station and recombines them. This interferometer obtained the first measurements for a number of giant stars (Labeyrie, 1985). Following the success of its operation, he undertook a project of building large interferometer known as grand interf$`\stackrel{`}{e}`$rom$`\overline{e}`$tre $`\stackrel{`}{a}`$ deux telescope (GI2T) at the same observatory. This interferometer comprises of two 1.5 meter spherical telescopes on a North-South baseline, which are movable on a railway track (Labeyrie et al., 1986). Mourard et al. (1989) had resolved the rotating envelope of hot star $`\gamma `$ Cassiopeiae using this interferometer. The technical details of this kind of interferometers can be found in the recent article by Saha (1999a). There are several long baseline interferometers that are in operation; some are at various stages of development (Saha, 1999a). These interferometers are based on the principle of merging speckles from both the telescopes. In other words, the fringed speckle can be visualized when a speckle from one telescope is merged with the speckle from the other telescope. Therefore, it is necessary to get acquainted with relatively new topics, speckle interferometry. In what follows, the formation of speckles and the way to decode the atmospherically degraded informations are discussed in brief. 3. Convolution and its applications The convolution of two functions is a mathematical procedure (Goodman, 1968) which simulates phenomena such as a blurring of a photograph. That may be caused by poor focus, by the motion of a photographer during the exposure, by dirt on the lens etc. In such blurred picture each point of object is replaced by a spread function. The spread function is disk shaped in the case of poor focus, line shaped if the photograph has moved, halo shaped if there is a dust on lens. In other words, we know that delta function has value at a single point otherwise it is zero. But generally the measurement does not produce this. Let us consider an input curve that can be represented by the curve $`f(x)`$ in terms of lot of close delta functions which are spread. Here, the shape of the response of the system including unwanted spread, is same for all values of $`x`$ (invariant for each considered delta function). Now, the value of the function $`f(x)`$ at $`x_1`$ is $`f(x^{})g(x_1x^{})`$. This is similar for each considered point on the curve. So for the whole curve, we define mathematically $$h(x)=_{\mathrm{}}^+\mathrm{}f(x^{})g(xx^{})𝑑x^{},$$ (1) where, $`h(x)`$ is the output value at particular point $`x`$, $``$ stands for convolution. This integral is defined as convolution of $`f(x)`$ and $`g(x)`$. $$h(x)=f(x)g(x),$$ (2) where, $`g(x)`$ is referred to as a blurring function or line spread function (LSF) or in two dimensions, the point spread function (PSF). The Fourier transform of a convolution of two functions is the product of the Fourier transform of the two functions. Therefore, in the Fourier plane the effect becomes a multiplication, point by point, of the transform of $`\widehat{F}(u)`$ with the transfer function $`\widehat{G}(u)`$. 4. Atmospheric turbulence and speckle formation Owing to the turbulent phenomena associated with heat flow and winds in the atmosphere, the density of air fluctuates in space and time. The inhomogeneities of the refractive index of the air can have devastating effect on the resolution achieved by any large telescope. The disturbance takes the form of distortion of the shape of the wavefront and variations of the intensity across the wavefront. Due to the motion and temperature fluctuations in the air above the telescope aperture, inhomogeneities in the refractive index develop. These inhomogeneities have the effect of breaking the aperture into cells with different values of refractive index that are moved by the wind across the telescope aperture. Kolmogrov law represents the distribution of turbule sizes, from millimeters to meters, with lifetimes varying from milliseconds to seconds. Changes in the refractive index in different portions of the aperture result to the phase changes in the value of the aperture function. The time evolution of the aperture function implies that the point spread function is time dependent. When the Reynolds number exceeds some value in a pipe (depending on the geometry of the pipe), the transition from laminar flows to turbulent flows occur. The dimensionless quantity Reynolds number is defined as $`QL/\nu `$, where $`Q`$ is the mean flow speed, $`L`$ is the transverse size of the pipe and $`\nu `$ is kinematic viscosity of the fluid. If $`L`$ is taken as some characteristic size of the flow of atmosphere the result holds good for atmospheric case. The power spectral density of refractive index fluctuations caused by the atmospheric turbulence follows a power law with large eddies having greater power (Tatarski, 1961). A plane wave propagating through the atmosphere of earth is distorted by refractive index variation in the atmosphere (troposphere); it suffers phase fluctuations and reaches the entrance pupil of a telescope with patches of random excursions in phase (Fried, 1966). Therefore, the image of the star in the focal plane of a large telescope is larger then the Airy disk of the telescope. The size is equivalent to the atmospheric point spread function (point spread function is a modulus square of the Fourier transform of the aperture function). The resolution at the image plane of the telescope is determined by the width of the PSF which is of the order of $`(1.22\lambda /r_{})`$, where, $`\lambda `$ is a wavelength of light and $`r_{}`$ is the average size of the turbulence cell, which is of the order of 10 cm. Therefore, resolution ($``$ 1 arc second) is unfortunately much larger then the theoretical size $`1.22\lambda /D`$ of the Airy disk of a large telescope (Rayleigh limit or diffraction limit), where, $`D`$ is the diameter of the telescope. The variance of phase difference fluctuations between any two points in the wave-front increases as the 5/3 power of their separation. When this variance exceeds $`\pi ^2`$ rad for some separation $`r_{}`$, then all details in the smaller than $`\lambda /r_{}`$ will be obliterated in the long exposure images. If the exposure time is shorter than the evolution time of the phase inhomogeneities, then each patch of the wave-front with diameter $`r_{}`$ $``$ Fried parameter $``$ would act independently of the rest of the wavefront resulting in multiple images of the source. These sub-images or ‘speckles’, as they are called and spread over the area defined by the long exposure image, can occur randomly along any direction within an angular patch of diameter $`\lambda /r_{}`$. The average size of the speckle is of the same order of magnitude as the Airy disk of the telescope in the absence of atmospheric turbulence and the lifetime of individual speckle is of the order of 0.1 to 0.01 seconds. Figure 1 depicts the speckles of the star HR4689; observations were carried out at 2.34 meter VBT, Kavalur, India with the speckle interferometer (Saha et al., 1997, 1999a). A snap shot taken later will show a different pattern but with similar probability of the angular distribution. A sum of similar exposures is the conventional image. It is easy to visualize that the sum of several statistically uncorrelated speckle patterns from a point source can result in an uniform patch of light a few arc-seconds wide (Saha, 1999b). figure 2 shows the result of summing 128 specklegrams demonstrating the destructions of finer details of the image by the atmospheric turbulence. Venkatakrishnan et al., (1989) had also generated the intensity distribution for $`r_{}`$ in the plane of a large telescope; the smallest contours have the size of the Airy disk of the telescope. The result of summing 100 such distributions showed similar concentric circle of equal intensity, destroying the finer details of an image. The method of such computer simulations runs as follows: The intensity distribution in the focal point of the telescope for the atmospheric cell ($`r_{}`$) of 10cm size and $`L_{}`$ = 200 cm for an entrance pupil of 200cm diameter were taken as samples. A power spectral density of the form $$\rho (k)\frac{L_{}^{\frac{11}{3}}}{(1+k^2L_{}^2)^{\frac{11}{6}}},$$ (3) was multiplied with a random phase factor $`e^{i\varphi }`$, one for each value of $`(k_x,k_y)`$, with $`\varphi `$ uniformly distributed between $`\pi `$ and $`\pi `$. The resulting 2-D pattern in $`k_x,k_y`$ space was Fourier transformed to obtain one realization of the wavefront $`W(x,y)`$. The Fraunhofer diffraction pattern of a piece of this wavefront with the diameter of the entrance pupil gives angular distribution of amplitudes, while the squared modulus of this field gives the intensity distribution in the focal plane of the telescope. The long-exposure PSF is defined by the ensemble average, $`<S(x,y)>`$, independent of any direction. The average illumination, $`I(x,y)`$ of a resolved object, $`O(x,y)`$ obeys convolution relationship, $$<I(x,y)>=O(x,y)<S(x,y)>,$$ (4) where, $`(x,y)`$ is a 2-dimensional space vector. Using 2-dimensional Fourier transform, this equation can be read as, $$<\widehat{I}(u,v)>=\widehat{O}(u,v)<\widehat{S}(u,v)>,$$ (5) where, $`\widehat{O}(u,v)`$ is the object spectrum, $`<\widehat{S}(u,v)>`$ is the transfer function for long-exposure images and is the product of the transfer function of the atmosphere $`\widehat{B}(u,v)`$, as well as the transfer function of the telescope, $`\widehat{T}(u,v)`$. $`(u,v)`$ is the spatial frequency vector. The transfer function for long-exposure image can be expressed as, $$<\widehat{S}(u,v)>=\widehat{B}(u,v)\widehat{T}(u,v).$$ (6) The benefit of the short-exposure images over long-exposure can be visualized by the following explanation. Let us consider two seeing cells separated by a vector in the telescope pupil, $`\lambda (u,v)`$, where $`\lambda `$ is the mean wavelength. If a point source is imaged through the telescope by using pupil function consisting of two apertures, corresponding to the two seeing cells, then a fringe pattern is produced with narrow spatial frequency bandwidth. If the major component $`\widehat{I}(u,v)`$ at the frequency $`(u,v)`$ is produced by contributions from all pairs of points with separations $`\lambda (u,v)`$, with one point in each aperture and is averaged over many frames, then the result for frequencies greater than $`r_o/\lambda `$ tends to zero. The Fourier component performs a random walk in the complex plane and average to zero, $`<\widehat{I}(u,v)>`$ = 0, when $`u>r_o/\lambda `$. For a large telescope, the aperture, P, can be sub-divided into a set of sub-apertures, $`p_i`$. According to the diffraction theory (Born and Wolf, 1984) the image at the focal plane of the telescope is obtained by adding all such fringe patterns produced by all possible pairs of sub-apertures. With increasing distance of the baseline between two sub-apertures, the fringes move with an increasingly larger amplitude. On a long-exposure images, no such shift is observed, which implies the loss of high frequency components of the image. While, in the short-exposure images ($`<`$20 msec), the interference fringes are preserved. 4.1 Seeing The common query for an observer is about ‘seeing’. This important parameter changes very fast, every now and then. It is the total effect of distortion in the path of the light via different contributing layers of the atmosphere to the detector placed at the focus of the telescope (detail discussions can be found in recent article by Saha, 1999a). The major sources of image degradation predominantly comes from the surface layer, as well as from the aero-dynamical disturbances in the atmosphere surrounding the telescope and its enclosure, namely, (i) thermal distortion of primary and secondary mirrors when they get heated up, (ii) dissipation of heat by the latter mirror, (iii) rise in temperature at the primary cell, (iv) at the focal point causing temperature gradient close to the detector etc. Saha and Chinnappan (1999) have found that the seeing at VBT improved gradually in the latter part of the night. The resolution $`\theta `$ of a large telescope, limited by the atmospheric turbulence, as defined by the Strehl criterion is $$\theta =\frac{4}{\varphi }\frac{\lambda }{r_o},$$ (7) Then the question arises how to measure seeing? The qualitative method is from the short exposure images using speckle interferometric technique, where, the area of the telescope aperture divided by the estimated number of speckles gives the wavefront coherence area $`\sigma `$, from which $`r_o`$ can be found by using relation, $$\sigma =0.342\left(\frac{r_o}{\lambda }\right)^2.$$ (8) If the autocorrelations of the short exposure images are summed, it contains autocorrelation of the seeing disk together with the autocorrelation of the mean speckle cell. It is width of the speckle component of the autocorrelation that provides information on the size of the object being observed (Saha and Chinnappan, 1999, Saha et al., 1999a). Systematic studies of this parameter would enable to understand the various causes of the local seeing, for example, thermal inhomogeneities associated with the building, aberrations in the design, manufacture and alignment of the optical train etc. Doom seeing plays a vital role in deteriorating image quality. It is necessary to take precautionary measure to avoid hot air entrapment. The important ones are (i) bring down the control room from the observing floor, (ii) improve the cross ventilation at each of the floors, (iii) remove unused machines or any other equipments, excess cemented portions available around the building, as well as all glass enclosures from the building. Mirror seeing is another important source of image spread and has the longest time-constant. The spread amounts to 0.5” for a 1 different in temperature. Therefore, It is essential to make arrangement to cool the primary mirror and try to maintain uniform temperature in and around the primary mirror cell (Saha and Chinnappan, 1999). 5. Speckle The term ’Speckle’ refers to a grainy structure observed when an uneven surface of an object is illuminated by a fairly coherent source. A good example of speckle phenomena may be observed at the river port when many boats are approaching towards the former at a particular time or in the swimming pool when many swimmers are present. Each boat or swimmer emits wave trains and interference between these random trains causes a speckled wave field on the water surface. Depending on the randomness of the source, spatial or temporal, speckles tend to appear. Spatial speckles may be observed when all parts of the source vibrate at same constant frequency but with different amplitude and phase, while temporal speckles are produced if all parts of it have uniform amplitude and phase. With a non-monochromatic vibration spectrum, in the case of random sources of light, spatio-temporal speckles are produced. The ground illumination produced by any star has fluctuating speckles, known as star speckles. It is too fast and faint, therefore, cannot be seen directly. Atmospheric speckles can be observed easily in a star image at the focus of a large telescope using a strong eyepiece. The star image looks like a pan of boiling water. If a short exposure image is taken, speckles can be recorded. The number of correlation cells is determined by the equation $`N=D/r_o`$. As the seeing improves, the number decreases. It is clear that speckles are caused by interference effect between wave element having random phases, rather then by ray bending effect. Its structure in astronomical images is the result of constructive and destructive bi-dimensional interferences between rays coming from different zones of incident wave. The statistical properties of speckle pattern depend both on the coherence of the incident light and the properties of random medium. Mathematically, speckles are simply the result of adding many sine functions having different, random characteristics. Since the positive and negative values cannot cancel out everywhere, adding an infinite number of such sine functions would result in a function with 100$`\%`$ constructed oscillations. 5.1 Imaging in the Presence of Atmosphere In the ideal condition, the resolution can be achieved in an imaging experiment, is limited only on imperfections in the optical system. If a collimated beam passes through the atmosphere and is collected by a telescope, the quality of the image formed is influenced by the atmospherically produced disturbances. Let the modulation transfer function (MTF) of an optical system be composed of the atmosphere and a telescope. The random wavefront tilt displaces the image but does not reduce its sharpness. If a short exposure is recorded, the image sharpness and MTF are insensitive to the tilt. While, in the case of long exposure ($`>`$ 20 msec) image, the image is spread during the exposure by its random variations of the tilt. Therefore, the image sharpness and the MTF are affected by wavefront tilt, as well as by the more complex shapes. In the short exposure case, a random factor associated with the tilt is extracted from the MTF before being taken the average, where in the long exposure case, no such factor is removed. Let us consider an imaging system consists of a simple lens based telescope in which the point spread function (PSF) is invariant to spatial shifts. An object (point source) at a point $`(x^{},y^{})`$ anywhere in the field of view will, therefore, produce a pattern $`S[(x,y)(x^{},y^{})]`$ across the image. If the object can emit incoherently, the image $`I(x,y)`$ of a resolved object $`O(x,y)`$ obeys a convolution relationships. The mathematical description of the convolution of two functions is of the form: $$I(x,y)=O(x^{},y^{})S[(x,y)(x^{},y^{})]d(x^{},y^{}),$$ (9) 5.2. Outline of the theory of speckle interferometry By integrating autocorrelation function of the successive short-exposure records rather than adding the images themselves, the diffraction-limited information can be obtained. Indeed, autocorrelation function of a speckle images preserves some of the information in the way which is not degraded by the co-adding procedure. For each of the short-exposure instantaneous record, the quasi-monochromatic incoherent imaging equation applies, $$I(x,y)=O(x,y)S(x,y),$$ (10) where, $`I(x,y)`$ is the instantaneous image intensity, $`O(x,y)`$ the object intensity, $`S(x,y)`$ the instantaneous PSF. The analysis of data may be carried out in two equivalent ways. In the spatial domain the ensemble average space autocorrelation is found giving the resultant imaging equation. In the Fourier plane the effect becomes a multiplication, point by point, of the transform of the object $`\widehat{O}(u,v)`$ with the transfer function $`\widehat{S}(u,v)`$, and therefore, equation 10 leads to, $$\widehat{I}(u,v)=\widehat{O}(u,v)\widehat{S}(u,v).$$ (11) The ensemble average of the power spectrum is given by, $$<\|\widehat{I}(u,v)\|^2>=\|\widehat{O}(u,v)\|^2<\|\widehat{S}(u,v)\|^2>.$$ (12) Hence, if $`<\|\widehat{S}(u,v)\|^2>`$ is known, $`\|\widehat{O}(u,v)\|`$<sup>2</sup> can be estimated. To find this, one has to observe a point source close to the object. The Fourier transform of a point source (delta function) is a constant, $`C_n`$. Then, equation 12, for a point source is $$<\|\widehat{I}_s(u,v)\|^2>=C_n^2<\|\widehat{S}(u,v)\|^2>.$$ (13) To find $`C_n`$<sup>2</sup>, one has to find the boundary condition. At the origin of the Fourier plane, $`S(u=0,v=0)`$ is unity. This is true for an incoherent source. Hence, $`C_n`$<sup>2</sup> is given by, $$C_n^2=<\|\widehat{I}_s(0,0)\|^2>/<\|\widehat{S}(0,0)\|^2>,$$ (14) i.e., $$C_n^2=<\|\widehat{I}_s(0,0)\|^2>.$$ (15) Using equation 15 in equation 13 gives, $$<\|\widehat{S}(u,v)\|^2>=<\|\widehat{I}_s(u,v)\|^2>/<\|\widehat{I}_s(0,0)\|^2>.$$ (16) From equations 16 and 11, the power spectrum of the object is given by $$\|\widehat{O}(u,v)\|^2=<\|\widehat{I}(u,v)\|^2>/[<\|\widehat{I}_s(u,v)\|^2>/<\|\widehat{I}_s(0,0)\|^2],$$ (17) i.e., the power spectrum of the object is the ratio of the average power spectrum of the image to the normalized average power spectrum of the point source. By Wiener-Kinchin theorem, the inverse Fourier transform of equation 17 gives the autocorrelation of the object. $$A[O(x,y)]=FT^{}[\|\widehat{O}(u,v)\|^2],$$ (18) where, A stands for autocorrelation. 6. Observational technique at VBT The programme of observing close binary systems (separation $`<`$ 1$``$) has been going on since 1996 using speckle interferometer at the Cassegrain focus of the 2.34 VBT, Kavalur, India (Saha, 1999a). The details of this interferometer can be found in the articles (Saha et al., 1997, 1999a), that samples the image scale at the Cassegrain focus of the said telescope to 0.015” per pixel of the intensified CCD. The wave-front falls on the focal plane and passes on to a microscope objective through a circular aperture of $``$350 $`\mu `$m of an optical flat kept at an angle of 15. This aperture was developed on a low expansion optical glass by devising a fine grinding precision mechanism at the laboratory (A. P. Jayarajan and S. K. Saha). The aperture was ground into the glass held at the angle of 75 with respect to the grinding axis. Interested readers may try it out to repeat the same. The enlarged beam is recorded after passing through a narrow band filter by a Peltier-cooled ICCD (386$`\times `$576) camera as detector which offers various option of exposure time, viz., 1 msecs, 5 msecs, 10 msecs, 20 msecs etc. It can operate in full frame, frame transfer and kinetic modes. Since CCD is cooled to $``$40, the dark noise is considerably low. Unlike the uncooled ICCD where data is stored in 8 bits, in this system, data is stored to 12 bits and can be archived to a Pentium PC. In full frame, as well as in frame transfer modes, the region of interests can be acquired at a faster speed. While in the kinetic mode, the image area can be kept small to satisfy requirements, therefore, the rest of the area is usable for the data storage. The surrounding star field of diameter $`\varphi `$ 10 mm gets reflected from the optical flat on to a plane mirror and is re-imaged on to an uncooled ICCD (Chinnappan et al., 1991) for guiding the object. 7. Summary Among others the most important observations made by means of speckle interferometry is the discovery of compact cluster, R136a (HD38268), of Doradus nebula in the Large Magallanic Clouds (Weigelt and Baier, 1985). Recent observations with adaptive optics system (Brandl et al., 1996) have revealed over 500 stars within the field of view 12.8”$`\times `$12.8” covering a magnitude range 11.2. Baba et al., (1994) have observed a binary star, $`\varphi And`$ (separation 0.53”) using imaging speckle spectroscopic method and found that the primary star (Be star) has an H$`\alpha `$ emission line while the companion has an H$`\alpha `$ absorption line. Developments of high angular resolution imaging have been going on in our Institute over a decade. Several experiments have been conducted at Vainu Bappu Observatory (Saha et al., 1987). The programme of speckle imaging at VBT has been a successful one. Now we are in a position to obtain informations of Fourier phase of the objects too (Saha et al., 1999b). Mapping of the certain interesting objects, viz., active galactic nuclei, proto-planetary nebulae will be undertaken in near future. References Anderson J. A., 1920, Astrophys. J., 51, 263. Baba N., Kuwamura S., Miura N., Norimoto Y., 1994, Ap. J., 431, L111. Born M., Wolf E., 1984, Principles of Optics, Pergamon Press. Brandl B., Sams B. J., Bertoldi F., Eckart A., Genzel R., Drapatz S., Hofmann R., Lowe M., Quirrenbach A., 1996, Ap. J., 466, 254. Brown R. H., 1974, ‘The Intensity Interferometry, its Applications to Astronomy’, Taylor & Francis, London. Brown R. H. and Twiss R. Q., 1958, Proc. Roy. Soc. A., 248, 222. Brown R. H., Davis J. and Allen L. R., 1967, MNRAS, 137, 375. Brown R. H., Jennison R. C. and Das Gupta M. K., 1952, Nature, 170, 1061. Chinnappan V., Saha S. K., Faseehana, 1991, Kod. Obs. Bull. 11, 87. Fizeau H., 1868, C. R. Acad. Sci. Paris, 66, 934. Fried D. C., 1966, J. Opt. Soc. Am., 56, 1972. Goodman J. W., 1968, Introduction to Fourier optics, McGraw Hill Book Co. NY. Labeyrie A., 1970, Astron. & Astrophys., 6, 85. Labeyrie A., 1975, Astrophys. J., 196, L71. Labeyrie A., 1985, 15th. Advanced Course, Swiss Society of Astrophys. and Astron. ed.. A. Benz, M. Huber and M. Mayor, 170. Labeyrie A., Schumacher G., Dugu$`\stackrel{´}{e}`$ M., Thom C., Bourlon P., Foy F., Bonneau D. and Foy R., 1986, Astron. & Astrophys., 162, 359. Michelson A. A., 1920, Astrophys. J., 51, 257. Michelson A. A., and Pease F. G., 1921, Astrophys. J., 53, 249. Mourard D., Bosc I., Labeyrie A., Koechlin A. and Saha S., 1989, Nature, 342, 520. Saha S. K., 1999a, Bull. Astron. Soc. Ind., 27, 443. Saha S. K., 1999b, Ind. J. Phys., 73b, 553. Saha S. K., Chinnappan V., 1999, Bull. Astron. Soc. Ind., 27, 327. Saha S. K., Jayarajan A. P., Sudheendra G., Umesh Chandra A., 1997, Bull. Astron. Soc. Ind., 25, 379. Saha S. K., Sridharan R., Sankarasubramanian K., 1999b, submitted to Bull. Astron. Soc. Ind. Saha S. K., Sudheendra G., Umesh Chandra A., Chinnappan V., 1999a, Exp. Astron, 9, 39. Saha S. K., Venkatakrishnan P., Jayarajan A. P., Jayavel N., 1987, Curr. Sci., 56, 985. Tatarski V. I., 1961, ‘Wave Propagation in a Turbulent Medium’ McGraw Hill. Venkatakrishnan P., Saha S. K., Shevgaonkar R. K., 1989, Proc. ‘Image Processing in Astronomy’, ed., T. Velusamy, 57. Weigelt G., Baier G., 1985, A & A., 150, L18. Appendix I Theorems of Fourier Transform 1. Linearity theorem $$F(\alpha g+\beta h)=\alpha F(g)+\beta F(h)$$ (19) i.e., the transform of sum of two functions is simply the sum of their individual transforms. 2. Similarity theorem A stretching of the co-ordinates in the space domain $`(x,y)`$ results in the construction of the co-ordinates in the frequency domain $`(f_x,f_y)`$ plus a change in the overall amplitude of the spectrum. i.e., if $$F(g(x,y))=G(f_x,f_y)$$ (20) then, $$F(g(ax+by))=G(f_x/a,f_y/b)/\|ab\|$$ (21) 3. Shift theorem Translation of a function in a space domain introduces a linear phase shift in the frequency domain. i.e., if $$F(g(x,y))=G(f_x,f_y)$$ (22) then, $$F(g(xa,yb))=G(f_x,f_y)\mathrm{exp}(2\pi j(f_xa+f_yb))$$ (23) 4. Perseval’s theorem This theorem is generally interpretable as a statement of conservation of energy. It says that the total energy in the real domain is equal to the total energy in the Fourier domain. i.e., $$F(g(x,y))=G(f_x,f_y)$$ (24) then $$_{\mathrm{}}^+\mathrm{}_{\mathrm{}}^+\mathrm{}\|g(x,y)\|^2𝑑x𝑑y=_{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}\|G(f_x,f_y)\|^2𝑑f_x𝑑f_y$$ (25) 5. Convolution Theorem The convolution of two functions in the space domain (an operation that will be found to arise frequently in the theory of linear system) is entirely equivalent of the more simple operation of multiplying their individual transform. i.e., $$F(g(x,y))=G(f_x,f_y)$$ (26) and $$F(h(x,y))=H(f_x,f_y)$$ (27) then $$F(_{\mathrm{}}^+\mathrm{}_{\mathrm{}}^+\mathrm{}g(\xi ,\eta )h(x\xi ,y\eta )𝑑\xi 𝑑\eta )=G(f_x,f_y)H(f_x,f_y)$$ (28) 6. Autocorrelation theorem This theorem may be regarded as special case of convolution theorem. The Fourier transform of autocorrelation of a function is the squared modulus of the Fourier transform. i.e. if $$F(g(x,y))=G(f_x,f_y)$$ (29) then, $$F(_{\mathrm{}}^+\mathrm{}_{\mathrm{}}^+\mathrm{}g(\xi ,\eta )g^{}(\xi x,\eta y)𝑑\xi 𝑑\eta )=\|G(f_x,f_y)\|^2$$ (30)
warning/0003/math0003056.html	ar5iv	text	# Untitled Document SUPER-BROWNIAN MOTION WITH REFLECTING HISTORICAL PATHS Research partially supported by NSF grant DMS-9700721. Krzysztof Burdzy Jean-François Le Gall Abstract. We consider super-Brownian motion whose historical paths reflect from each other, unlike those of the usual historical super-Brownian motion. We prove tightness for the family of distributions corresponding to a sequence of discrete approximations but we leave the problem of uniqueness of the limit open. We prove a few results about path behavior for processes under any limit distribution. In particular, we show that for any $`\gamma >0`$, a “typical” increment of a reflecting historical path over a small time interval $`\mathrm{\Delta }t`$ is not greater than $`(\mathrm{\Delta }t)^{3/4\gamma }`$. 1. Introduction. The present article has been inspired by two probabilistic models—superprocesses with interactions and reflected particle systems. The first person to study a reflecting system of particles was Harris \[H\] who considered an infinite system of Brownian particles on the line. He proved that if the initial positions of the particles are points of a Poisson point process, then for a large time $`t`$ the distribution of a single particle is normal with the standard deviation $`(2t/\pi )^{1/4}`$. Spitzer \[S\] analyzed a similar model with particles moving along straight lines between collisions. See \[DGL1, DGL2, G, Ho\] for related results. The simplest superprocesses, for example, super-Brownian motion, are continuum limits of branching systems in which the branching mechanism is independent of the positions of particles. There has been considerable activity studying models with interactions. Many articles are devoted to models with catalysts, see, e.g., \[DF, De\]. Various other models with interactions are discussed in \[AT, BHM, EP, P3\]. See in particular \[P4\] and references therein. We will study a model similar to that introduced by Harris, in that we will start with linear Brownian motion as the spatial process. We will attempt to build a corresponding superprocess with historical paths that do not cross over although they may touch each other. Our construction is based on a sequence of discrete approximations. Consider for every $`\epsilon (0,1]`$ a branching particle system which starts initially with $`N_\epsilon `$ particles located respectively at $`x_1^\epsilon \mathrm{}x_{N_\epsilon }^\epsilon `$. Particles move independently in space according to linear Brownian motion and are subject to critical binary branching at rate $`\epsilon ^1`$. To be specific, the lifetimes of the particles are exponential with parameter $`\epsilon ^1`$ and when a particle dies it gives rise to $`0`$ or $`2`$ new particles with probability $`1/2`$. Let us now introduce our basic assumptions. Let $$\mu _\epsilon :=\epsilon \underset{j=1}{\overset{N_\epsilon }{}}\delta _{x_j^\epsilon }$$ and assume that there is a finite measure $`\mu `$ on $`R`$ such that $$\mu _\epsilon \underset{\epsilon 0}{\overset{(\mathrm{w})}{}}\mu ,$$ $`(1.1)`$ where the notation (w) indicates weak convergence in the space $`M_f(R)`$ of finite measures on $`R`$. In addition, if $`\mathrm{supp}\mu `$ denotes the topological support of $`\mu `$, we assume that $$\mathrm{supp}\mu _\epsilon \underset{\epsilon 0}{}\mathrm{supp}\mu ,$$ $`(1.2)`$ in the sense of the Hausdorff metric on compact subsets of $`R`$ (in particular, we assume that $`\mathrm{supp}\mu `$ is compact). Let $`X_t^\epsilon `$ denote the random measure equal to $`\epsilon `$ times the sum of the Dirac point masses at the positions of particles alive at time $`t`$. Then, $$(X_t^\epsilon ,t0)\underset{\epsilon 0}{\overset{(\mathrm{d})}{}}(X_t,t0),$$ $`(1.3)`$ where the limit process is super-Brownian motion in $`R`$ with branching rate $`\gamma =1`$ (throughout this work we consider only this branching rate) and initial value $`\mu `$, and the convergence holds in distribution in the Skorohod space $`𝐃(R_+,M_f(R))`$. The convergence (1.3) is the standard approximation of super-Brownian motion (see e.g. \[P4\]). Note that assumption (1.2) is not needed for (1.3) but it guarantees that the graph of $`X^\epsilon `$ also converges in distribution to the graph of $`X`$ (see Lemma 2.3 below), a property that plays an important role in our arguments. For each particle alive at time $`t`$, we can consider its historical path, which is the element of $`𝐂([0,t],R)`$ obtained by concatenating the trajectories of the ancestors of the given particle up to time $`t`$. Denote by $`Y_t^\epsilon `$ the historical measure equal to $`\epsilon `$ times the sum of the Dirac point masses at the historical paths of the particles alive at time $`t`$ ($`Y_t^\epsilon `$ is thus a random measure on the set $`𝐂([0,t],R)`$ of continuous mappings from $`[0,t]`$ into $`R`$). Then the convergence (1.3) can be reinforced as $$(Y_t^\epsilon ,t0)\underset{\epsilon 0}{\overset{(\mathrm{d})}{}}(Y_t,t0),$$ $`(1.4)`$ where the limit process is now historical super-Brownian motion started at $`\mu `$. For every $`\epsilon >0`$, we can use the original branching particle system to construct a new system with reflection. The branching mechanism (critical binary branching at rate $`\epsilon ^1`$) is the same as in the original system, but the particle paths in the new system reflect against each other. A precise construction is given in Section 3, but let us give an informal description. The reflected system is such that for every $`t0`$, the set of positions of particles at time $`t`$ is the same as in the original system, and in particular the branching times are the same. During the time interval between $`0`$ and the first branching time, the vector of positions of the particles labeled $`1,2,\mathrm{},N_\epsilon `$ in the reflected system is the increasing rearrangement of the vector of positions of the particles in the original system. Suppose that at the first branching time, denoted by $`\xi `$, a particle dies and gives rise to 2 children. If the location of this particle is the $`j`$-th coordinate in the increasing rearrangement of the vector of positions at time $`\xi `$, we will say that in the reflected system particle $`j`$ has given rise to two children labeled $`j1`$ and $`j2`$. Then on the interval between $`\xi `$ and the second branching time, the vector of positions of the particles labeled $`1,\mathrm{},j1,j1,j2,j+1,\mathrm{},N_\epsilon `$ in the reflected system is again the increasing rearrangement of the vector of positions of the particles in the original system. We can easily continue this construction by induction. Denote by $`\stackrel{~}{X}_t^\epsilon `$ and $`\stackrel{~}{Y}_t^\epsilon `$ the analogues of $`X_t^\epsilon `$ and $`Y_t^\epsilon `$ for the the system with reflection. We have $`\stackrel{~}{X}_t^\epsilon =X_t^\epsilon `$ since the set of positions of particles is the same at every time $`t`$ in both systems. On the other hand, $`\stackrel{~}{Y}_t^\epsilon `$ is typically very different from $`Y_t^\epsilon `$. Indeed, the following property holds for any two paths $`w`$, $`w^{}`$ in the support of $`\stackrel{~}{Y}_t^\epsilon `$: Either $`w(r)w^{}(r)`$ for every $`0rt`$, or $`w(r)w^{}(r)`$ for every $`0rt`$. The main purpose of this work is to try to understand the limiting behavior of the branching particle system with reflection as $`\epsilon 0`$. Our primary objective was to get an analogue of the convergence (1.4) when the processes $`Y^\epsilon `$ are replaced by $`\stackrel{~}{Y}^\epsilon `$, giving information about the individual paths in the system with reflection. We did not completely succeed in this task, but we can prove the following result, where $`𝒲`$ denotes the set of all stopped paths, or equivalently the union over all $`t0`$ of the sets $`𝐂([0,t],R)`$. Theorem 1.1. Let $``$ be a sequence of positive numbers converging to $`0`$. The laws of the processes $`(\stackrel{~}{Y}_t^\epsilon ,t0)`$ for $`\epsilon `$ are tight in the space of all probability measures on the Skorohod space $`𝐃([0,\mathrm{}),M_f(𝒲))`$. Furthermore any limiting distribution is supported on $`𝐂([0,\mathrm{}),M_f(𝒲))`$. Hence, by extracting a subsequence if necessary, we can assume that the sequence of processes $`\stackrel{~}{Y}^\epsilon `$ converges in distribution towards a process $`\stackrel{~}{Y}`$ with continuous paths with values in $`M_f(𝒲)`$. Note that, for every $`t0`$, the measure $`\stackrel{~}{Y}_t`$ is supported on $`𝐂([0,t],R)`$. Although the question of uniqueness of the limit remains unsolved, we are able to derive several results on the path behavior of the process $`\stackrel{~}{Y}`$. First note that, since $`\stackrel{~}{X}_t^\epsilon =X_t^\epsilon `$ for every $`t0`$, the convergence (1.1) implies that the $`M_f(R)`$-valued process $`\stackrel{~}{X}`$ defined by $$\stackrel{~}{X}_t,\phi =\stackrel{~}{Y}_t(dw)\phi (w(t))$$ is a super-Brownian motion started at $`\mu `$. In particular, it is known (see \[KS\], \[R\]) that a.s. for every $`t>0`$ the measure $`\stackrel{~}{X}_t(dy)`$ has a density denoted by $`x_t(y)`$, and that there exists a jointly continuous modification of $`(x_t(y),t>0,yR)`$. The next result shows that for any $`\gamma >0`$, a typical oscillation of a reflecting historical path is not greater than $`(\mathrm{\Delta }t)^{\frac{3}{4}\gamma }`$, and hence much smaller than a typical Brownian oscillation $`(\mathrm{\Delta }t)^{\frac{1}{2}}`$. This result is consistent with the Harris \[H\] estimate, if we translate the large-time asymptotics to small-time asymptotics. Theorem 1.2. Almost surely for every $`t>0`$ and every $`r(0,t)`$, for every path $`w\mathrm{supp}\stackrel{~}{Y}_t`$, the condition $`x_r(w(r))>0`$ implies that, for every $`\gamma >0`$, $$\underset{\delta 0}{lim\; sup}\frac{\|w(r+\delta )w(r)\|}{\delta ^{\frac{3}{4}\gamma }}=0.$$ A more precise version of Theorem 1.2 is given in Section 5 (Theorem 5.10). It is not hard to check that if we fix $`t>0`$ and $`r(0,t)`$ (fixing $`r`$ is in fact enough), the condition $`x_r(w(r))>0`$, and thus the conclusion of the theorem, will hold for every path $`w\mathrm{supp}\stackrel{~}{Y}_t`$, a.s. Alternatively, for every fixed $`t>0`$, the conclusion of Theorem 1.2 holds for a set of values of $`r(0,t)`$ of full Lebesgue measure, for every $`w\mathrm{supp}\stackrel{~}{Y}_t`$. We believe that $`\delta ^{\frac{3}{4}}`$ is the “typical” size for the oscillation of a historical reflected path although we have no lower bound justifying this claim. We also study the behavior of reflected historical paths at a branching point. If $`w`$ and $`w^{}`$ are two reflected historical paths that coincide up to time $`r>0`$ (meaning informally that the corresponding “particles” have the same ancestor up to time $`r`$), we show that the distance between $`w(r+\delta )`$ and $`w^{}(r+\delta )`$ grows linearly as a function of $`\delta `$, up to logarithmic corrections. The precise statement is as follows. Theorem 1.3. Let $`t>0`$. If $`w`$ and $`w^{}`$ are two distinct elements of $`𝐂([0,t],R)`$, we set $$\gamma _{w,w^{}}=inf\{r0:w(r)w^{}(r)\}.$$ Then a.s. for any two distinct paths $`w,w^{}\mathrm{supp}\stackrel{~}{Y}_t`$ such that $`\gamma _{w,w^{}}>0`$, we have $$\underset{\delta 0}{lim\; sup}\frac{\|w(\gamma _{w,w^{}}+\delta )w^{}(\gamma _{w,w^{}}+\delta )\|}{2\delta \mathrm{log}\|\mathrm{log}\delta \|}=x_{\gamma _{w,w^{}}}(w(\gamma _{w,w^{}}))>0$$ and, for every $`\gamma >0`$, $$\underset{\delta 0}{lim}\frac{\|w(\gamma _{w,w^{}}+\delta )w^{}(\gamma _{w,w^{}}+\delta )\|}{\delta \|\mathrm{log}\delta \|^{1\gamma }}=\mathrm{}.$$ Our proofs rely on several known results on super-Brownian motion. In particular, we use the Brownian snake idea \[L2\] in an essential way, both in the proofs and for giving more precise versions of the results. For instance, as a key step towards Theorem 1.1, we get a uniform continuity result (Theorem 4.1) for the historical paths of the approximating branching particle systems with reflection. The proof of this result requires some precise information about the genealogical structure of the approximating systems, which seems to be more easily accessible via the snake approach (cf Lemma 2.1 below). For an introduction to the theory of superprocesses (measure-valued diffusions) and historical processes, the reader may consult \[Da, Dy, DP, L2, P4\]. The paper is organized as follows. Section 2 describes the specific coding that we use to represent the genealogical structure of the approximating branching particle systems. This section also contains a few important preliminary results. Section 3 presents the construction of the systems with reflection. Tightness results are given in Section 4, including a more precise form of Theorem 1.1. Section 5 contains the proof of Theorem 1.2, and is the most technical part of the paper. Finally, Theorem 1.3 is proved in Section 6. We are grateful to Carl Mueller, Ed Perkins, Tokuzo Shiga and Roger Tribe for very useful advice. 2. Coding discrete trees. We will describe a method that provides a coding of the genealogy of the branching particle systems introduced in Section 1, in a consistent way for all values of the parameter $`\epsilon (0,1]`$. This method involves embedding branching trees in a path of reflected Brownian motion, and is based on \[L1\] (see also \[NP\]). 2.1 Markov chains embedded in reflected Brownian motion. Let $`\beta =(\beta _s,s0)`$ be distributed as twice a reflected Brownian motion on $`R_+`$: $$(\beta _s,s0)\stackrel{\left(\mathrm{d}\right)}{=}(2\|\mathrm{B}_\mathrm{s}\|,\mathrm{s}0),$$ where $`B`$ is a standard linear Brownian motion, with $`B_0=0`$. The reason for the factor $`2`$ will be clear later. We denote by $`(L_s^x,x0,s0)`$ the jointly continuous family of local times of $`\beta `$, normalized in such a way that, for every nonnegative Borel function $`\phi `$ on $`R_+`$, $$_0^t\phi (\beta _s)𝑑s=_{R_+}\phi (x)L_t^x𝑑x.$$ Also set $`\tau _r=inf\{s0:L_s^0>r\}`$, for every $`r>0`$. For every $`\epsilon (0,1]`$, we introduce a sequence of stopping times $`(T_k^\epsilon ,k=0,1,\mathrm{})`$ defined inductively as follows: $$\begin{array}{cc}\hfill T_0^\epsilon & =inf\{s0:\beta _s=2\epsilon \},\hfill \\ \hfill T_{2k+1}^\epsilon & =inf\{uT_{2k}^\epsilon :\underset{T_{2k}^\epsilon su}{sup}\beta _s\beta _u=2\epsilon \},\hfill \\ \hfill T_{2k+2}^\epsilon & =inf\{uT_{2k+1}^\epsilon :\beta _u\underset{T_{2k+1}^\epsilon su}{inf}\beta _s=2\epsilon \}.\hfill \end{array}$$ It is simple to check that the variables $`T_0^\epsilon ,T_1^\epsilon T_0^\epsilon ,T_2^\epsilon T_1^\epsilon ,\mathrm{}`$ are independent and identically distributed. To see this, note that if $`(\gamma _t,t0)`$ is a reflected Brownian motion with initial value $`\gamma _0=b0`$, the process $$\gamma _t\underset{0st}{inf}\gamma _s$$ is again a reflected Brownian motion, with initial value $`0`$, and also observe that $`\beta _{T_{2k}^\epsilon }2\epsilon `$ for every $`k`$. As $`E(T_0^\epsilon )=\epsilon ^2`$, standard arguments show that for every $`K>0`$ $$\underset{sK}{sup}\left\|T_{[s/\epsilon ^2]}^\epsilon s\right\|\underset{\epsilon 0}{\overset{\mathrm{a}.\mathrm{s}.}{}}0.$$ $`(2.1)`$ (First establish this convergence along the sequence $`\epsilon _n=n^2`$ and then use monotonicity arguments.) Thus, $$\underset{sK}{sup}\left\|\beta _{T_{[s/\epsilon ^2]}^\epsilon }\beta _s\right\|\underset{\epsilon 0}{\overset{\mathrm{a}.\mathrm{s}.}{}}0.$$ $`(2.2)`$ For $`k=0,1,\mathrm{}`$, set $$\begin{array}{cc}\hfill S_{2k}^\epsilon & =\beta _{T_{2k}^\epsilon }2\epsilon ,\hfill \\ \hfill S_{2k+1}^\epsilon & =\beta _{T_{2k+1}^\epsilon }.\hfill \end{array}$$ It is easy to verify that $`(S_k^\epsilon ,k=0,1,2,\mathrm{})`$ is a time-inhomogeneous Markov chain with values in $`R_+`$, whose law can be described as follows (see \[L1\] Section 3 for details): $`S_0^\epsilon =0`$ and $`S_{2k+1}^\epsilon `$ has the same distribution as $`S_{2k}^\epsilon +U`$, where $`U`$ is an exponential variable with mean $`2\epsilon `$, independent of $`S_{2k}^\epsilon `$, $`S_{2k+2}^\epsilon `$ has the same distribution as $`(S_{2k+1}^\epsilon V)_+`$ where $`V`$ is exponential with mean $`2\epsilon `$, independent of $`S_{2k+1}^\epsilon `$. From (2.2), we have a.s. for every $`K>0`$, $$\underset{sK}{sup}\left\|S_{[s/\epsilon ^2]}^\epsilon \beta _s\right\|\underset{\epsilon 0}{\overset{\mathrm{a}.\mathrm{s}.}{}}0.$$ We then define a continuous-time process $`(\beta _s^\epsilon ,s0)`$ by setting $$\beta _{k\epsilon ^2}^\epsilon =S_k^\epsilon \text{ for }k=0,1,2,\mathrm{}$$ and by interpolating linearly on intervals of the form $`[k\epsilon ^2,(k+1)\epsilon ^2]`$. It is obvious that we also have $$\underset{sK}{sup}\left\|\beta _s^\epsilon \beta _s\right\|\underset{\epsilon 0}{\overset{\mathrm{a}.\mathrm{s}.}{}}0.$$ $`(2.3)`$ 2.2 The correspondence between excursions and trees With each excursion of $`\beta ^\epsilon `$ away from $`0`$, we can associate a marked tree representing the genealogical structure of a Galton-Watson branching process with critical binary branching at rate $`\epsilon ^1`$, starting with one individual (the ancestor) at time $`0`$. Here a marked tree consists of the set $`𝒯`$ of edges (i.e., particles), which is a subset of $$𝐔:=\underset{n=0}{\overset{\mathrm{}}{}}\{1,2\}^n\text{(by convention, }\{1,2\}^0=\{\mathrm{}\}),$$ and the family $`(\mathrm{}_u,u𝒯)`$ of lengths of edges (i.e., lifetimes of particles). Figure 1. This correspondence is explained in Fig. 1 for the first excursion of $`\beta ^\epsilon `$ away from $`0`$. Informally, if $`(i\epsilon ^2,j\epsilon ^2)`$ is the interval corresponding to an excursion of $`\beta ^\epsilon `$, the lifetime $`\mathrm{}_{\mathrm{}}`$ of the individual at the root of the associated tree is the minimum of $`\beta ^\epsilon `$ over $`[(i+1)\epsilon ^2,(j1)\epsilon ^2]`$ and this individual has two children if and only if $`j1>i+1`$. In that case, by decomposing the excursion restricted to $`[(i+1)\epsilon ^2,(j1)\epsilon ^2]`$ at the time of its minimum over this interval, we get two new excursions, each of which codes the genealogical structure of descendants of one of the ancestor’s children. The construction of the tree is then completed by induction. Note that each time of the form $`k\epsilon ^2`$ in the interval $`(i\epsilon ^2,j\epsilon ^2)`$ corresponds to one edge of the tree (for instance the time of the minimum over $`[(i+1)\epsilon ^2,(j1)\epsilon ^2]`$ corresponds to $`\mathrm{}`$, see Fig. 1). We refer to \[L1\] Section 2 for a more precise description and a proof that this construction yields the family tree of a Galton-Watson branching process with critical binary branching at rate $`\epsilon ^1`$. (We can now explain the factor $`2`$ in the definition of $`\beta `$: We want the branching rate to be $`\epsilon ^1`$ and not $`(\epsilon /2)^1`$.) There is a one-to-one correspondence between excursions of $`\beta ^\epsilon `$ away from $`0`$ and excursions of $`\beta `$ away from $`0`$ with height greater than $`2\epsilon `$: If $`k\epsilon ^2`$ is the beginning of an excursion of $`\beta ^\epsilon `$, then $`T_k^\epsilon `$ is the hitting time of $`2\epsilon `$ by the corresponding excursion of $`\beta `$. As in Section 1, consider for every $`\epsilon (0,1]`$ an integer $`N_\epsilon 1`$ and assume that the family $`(\epsilon N_\epsilon ,\epsilon (0,1])`$ is bounded and that $`\epsilon N_\epsilon `$ converges to $`a0`$ as $`\epsilon 0`$ (this follows from (1.1) with $`a=\mu ,1`$). Let $`\tau ^\epsilon `$ denote the $`N_\epsilon `$-th return of $`\beta ^\epsilon `$ to $`0`$. From the previous observations, (2.1) and the standard approximation of Brownian local times by upcrossing numbers, we have $$\underset{\epsilon 0}{lim}\tau ^\epsilon =\tau _a,\text{ a.s.}$$ We will write $`\tau =\tau _a`$ for simplicity. On the time interval $`[0,\tau ^\epsilon ]`$, the process $`\beta ^\epsilon `$ makes $`N_\epsilon `$ independent excursions away from $`0`$. These excursions can be viewed as representing the genealogical structure of the branching particle system introduced in Section 1. The set of edges, denoted by $`𝒯_\epsilon `$, is then a random subset of $`\{1,\mathrm{},N_\epsilon \}\times 𝐔`$ and conditionally on $`𝒯_\epsilon `$, the corresponding lengths are independent exponentials with mean $`\epsilon `$. The function $`(\beta _s^\epsilon ,s[0,\tau ^\epsilon ])`$ can be reconstructed from this collection of marked trees as shown by Fig. 1. Notice that for this reconstruction to be possible, it is essential to order the trees and the different edges of every single tree. 2.3 Discrete and continuous local times One reason for considering the processes $`\beta ^\epsilon `$ comes from their relation with the upcrossing numbers of $`\beta `$. We first define the (discrete) local times of $`\beta ^\epsilon `$. For every $`x0`$ and $`s0`$, we define $$L_s^{\epsilon ,x}=\epsilon Card\{r[0,s):\beta _r^\epsilon =x\mathrm{and}\beta _u^\epsilon >x\mathrm{for}u(r,r+\delta ],\mathrm{for}\mathrm{some}\delta >0\}.$$ In other words, $`\epsilon ^1L_s^{\epsilon ,x}`$ is the number of upcrossings of $`\beta ^\epsilon `$ above level $`x`$ before time $`s`$. Let $`M_s^\epsilon (x)`$ denote the number of upcrossings of $`\beta `$ from $`x`$ to $`x+2\epsilon `$ completed before time $`s`$. More precisely, $`M_s^\epsilon (x)`$ is the number of pairs $`(u,v)`$ with $`0u<v<s`$, such that $`\beta _u=x`$, $`\beta _r>x`$ for every $`r(u,v)`$ and $`v=inf\{r>u:\beta _r>x+2\epsilon \}`$. Then, a.s. for every $`x0`$ and every integer $`k1`$, we have $$L_{(2k1)\epsilon ^2}^{\epsilon ,x}=L_{2k\epsilon ^2}^{\epsilon ,x}=\epsilon M_{T_{2k}^\epsilon }^\epsilon (x)=\epsilon M_{T_{2k1}^\epsilon }^\epsilon (x).$$ $`(2.4)`$ This identity is easily verified by induction on $`k`$ (the sequence of stopping times $`(T_k^\epsilon )`$ was designed for this property to hold). See also Proposition 7 of \[L1\]. Lemma 2.1. We have a.s. $$\underset{\epsilon 0}{lim}\left(\underset{s0}{sup}\underset{x0}{sup}\|L_{s\tau ^\epsilon }^{\epsilon ,x}L_{s\tau }^x\|\right)=0.$$ Proof. We first observe that a.s. $$\underset{\epsilon 0}{lim}\left(\underset{s0}{sup}\underset{x0}{sup}\|\epsilon M_{s\tau ^\epsilon }^\epsilon (x)L_{s\tau }^x\|\right)=0.$$ $`(2.5)`$ For a fixed value of $`x`$, this is nothing but the classical approximation of Brownian local time by upcrossing numbers, and excursion theory provides precise estimates for the rate of convergence. Using these estimates and monotonicity properties, it is then an easy task to prove (2.5), i.e., the uniform version of the claim. The statement of the lemma is now a simple consequence of (2.1), (2.4) and (2.5). Remark. As an immediate consequence of Lemma 2.1 and the joint continuity of Brownian local times, we get that $$\underset{\epsilon ,\delta 0}{lim}\left(\underset{s0}{sup}\underset{\|xx^{}\|\delta }{\underset{x,x^{}0}{sup}}\|L_{s\tau ^\epsilon }^{\epsilon ,x}L_{s\tau ^\epsilon }^{\epsilon ,x^{}}\|\right)=0,\text{a.s.}$$ Later, we will consider for every $`\epsilon (0,1]`$ a process $`\stackrel{~}{\beta }^\epsilon `$ with the same distribution as $`\beta ^\epsilon `$. If $`\stackrel{~}{L}_s^{\epsilon ,x}`$ denote the discrete local times of $`\stackrel{~}{\beta }^\epsilon `$, the last convergence still holds in probability when $`L_s^{\epsilon ,x}`$ is replaced by $`\stackrel{~}{L}_s^{\epsilon ,x}`$ (and $`\tau ^\epsilon `$ by $`\stackrel{~}{\tau }^\epsilon `$, with an obvious notation). 2.4 Branching particle systems and discrete snakes. We now consider the branching particle system of Section 1, starting with $`N_\epsilon `$ particles located respectively at $`x_1^\epsilon ,x_2^\epsilon ,\mathrm{},x_{N_\epsilon }^\epsilon `$. We may and will assume that the genealogy of the descendants of particle $`k`$ (present at $`x_k^\epsilon `$ at time $`0`$) is given by the tree associated with the $`k`$-th excursion of $`\beta ^\epsilon `$ (cf subsection 2.2). We will refer to this system as the $`\epsilon `$-system of branching Brownian motions. For our purposes, it will be convenient to view the collection of paths traced by the branching particles as the range of a path-valued process called the discrete snake. By definition, a stopped path is a continuous mapping $`w:[0,\zeta ]R`$, where $`\zeta =\zeta _w0`$ is called the “lifetime” of $`w`$ (it is convenient to talk about the “lifetime” of a path although for technical reasons the path is stopped rather than killed). Let $`𝒲`$ be the set of all stopped paths. Then $`𝒲`$ is a separable complete metric space for the distance $$d(w,w^{})=\|\zeta _w\zeta _w^{}\|+\underset{t0}{sup}\|w(t\zeta _w)w^{}(t\zeta _w^{})\|.$$ For any $`xR`$, we write $`\underset{¯}{x}`$ for the trivial path such that $`\zeta _{\underset{¯}{x}}=0`$ and $`\underset{¯}{x}(0)=x`$. With every $`s[0,\tau ^\epsilon ]`$ we now associate a stopped path $`W_s^\epsilon 𝒲`$ with lifetime $`\beta _s^\epsilon `$. If $`s[0,\tau ^\epsilon )\epsilon ^2N`$ and $`\beta _s^\epsilon =0`$, then $`s`$ is the starting time of the $`k`$-th excursion of $`\beta ^\epsilon `$ away from $`0`$, for some $`k\{1,\mathrm{},N_\epsilon \}`$. We then set $`W_s^\epsilon =\underset{¯}{x}_k^\epsilon `$. For definiteness, we also set $`W_{\tau ^\epsilon }^\epsilon =\underset{¯}{x}_{N_\epsilon }^\epsilon `$. If $`s[0,\tau ^\epsilon )\epsilon ^2N`$ but $`\beta _s^\epsilon >0`$, we can associate with $`s`$ a unique edge of the $`k`$-th tree, $`k`$ being the number of the excursion straddling $`s`$. We then let $`W_s^\epsilon `$ be the historical path of the particle in the system of branching Brownian motions that corresponds to this edge. Notice that the death time of this particle is $`\beta _s^\epsilon `$, and thus $`\zeta _{W_s^\epsilon }=\beta _s^\epsilon `$. Finally if $`s[0,\tau ^\epsilon ]`$ but $`s\epsilon ^2N`$, we find an integer $`j`$ such that $`j\epsilon ^2<s<(j+1)\epsilon ^2`$, and let $`l=j`$ if $`\beta _{j\epsilon ^2}^\epsilon >\beta _{(j+1)\epsilon ^2}^\epsilon `$, but $`l=j+1`$ if $`\beta _{j\epsilon ^2}^\epsilon \beta _{(j+1)\epsilon ^2}^\epsilon `$. Then we let $`W_s^\epsilon `$ be the path $`W_{l\epsilon ^2}^\epsilon `$ stopped at time $`\beta _s^\epsilon `$. It is easy to see that conditionally on $`(\beta _s^\epsilon ,s0)`$ the process $`(W_{k\epsilon ^2}^\epsilon ,0k\tau ^\epsilon /\epsilon ^2)`$ is Markovian. To describe its conditional distribution, let $`k\{0,\mathrm{},\tau ^\epsilon /\epsilon ^2\}`$ and suppose that $`\beta _{(k+1)\epsilon ^2}^\epsilon >0`$ (otherwise $`W_{(k+1)\epsilon ^2}^\epsilon =\underset{¯}{x}_j^\epsilon `$, if $`(k+1)\epsilon ^2`$ is the starting point of the $`j`$-th excursion of $`\beta ^\epsilon `$). If $`\beta _{(k+1)\epsilon ^2}^\epsilon \beta _{k\epsilon ^2}^\epsilon `$ (which occurs if $`k`$ is odd) then $`W_{(k+1)\epsilon ^2}^\epsilon `$ is simply the restriction of $`W_{k\epsilon ^2}^\epsilon `$ to $`[0,\beta _{(k+1)\epsilon ^2}^\epsilon ]`$. On the other hand, if $`\beta _{(k+1)\epsilon ^2}^\epsilon >\beta _{k\epsilon ^2}^\epsilon `$, then $`W_{(k+1)\epsilon ^2}^\epsilon `$ is obtained from $`W_{k\epsilon ^2}^\epsilon `$ by “adding at the tip of $`W_{k\epsilon ^2}^\epsilon `$” a Brownian path of length $`\beta _{(k+1)\epsilon ^2}^\epsilon \beta _{k\epsilon ^2}^\epsilon `$ independent of $`(W_{j\epsilon ^2}^\epsilon ,jk)`$. The following snake property is a consequence of the definition of $`W_s^\epsilon `$: If $`s<s^{}`$ and $`s`$ and $`s^{}`$ belong to the same (open) excursion interval of $`\beta ^\epsilon `$ away from $`0`$, then $`W_s^\epsilon (t)=W_s^{}^\epsilon (t)`$ for every $`t[0,inf_{u[s,s^{}]}\beta _u^\epsilon ]`$. 2.5 Convergence to super-Brownian motion As in Section 1, we let $`X_t^\epsilon `$ be $`\epsilon `$ times the sum of the point masses at the positions of the particles alive at time $`t`$ in the $`\epsilon `$-system. This is equivalent to writing $$X_t^\epsilon =_0^{\tau _\epsilon }𝑑L_s^{\epsilon ,t}\delta _{W_s^\epsilon (t)}.$$ To justify this formula, recall the correspondence between excursions and trees described in Subsection 2.2 and note that each upcrossing time $`s`$ of $`\beta ^\epsilon `$ above level $`t`$ corresponds to one particle alive at time $`t`$, whose position is $`W_s^\epsilon (t)`$. Similarly, the historical process $`Y_t^\epsilon `$ is $$Y_t^\epsilon =_0^{\tau _\epsilon }𝑑L_s^{\epsilon ,t}\delta _{W_s^\epsilon }.$$ Recall our assumptions (1.1) and (1.2) and the convergence result in (1.3). We next prove a result about the uniform modulus of continuity for the paths $`W_s^\epsilon `$. For convenience, we make the convention that $`W_s^\epsilon (t)=W_s^\epsilon (\beta _s^\epsilon )`$ when $`t>\beta _s^\epsilon `$. Lemma 2.2. Let $`\eta (0,\frac{1}{2})`$. Then, $$\underset{\delta 0}{lim}\left(\underset{\epsilon (0,1]}{inf}P\left[\|W_s^\epsilon (t+r)W_s^\epsilon (t)\|r^{\frac{1}{2}\eta },\text{for every }t0,r[0,\delta ],s[0,\tau ^\epsilon ]\right]\right)=1.$$ Remark. This is of course reminiscent of the uniform modulus of continuity for historical paths of super-Brownian motion. This lemma is therefore very close to the results of \[DIP\] and \[DP\], which however use different approximations. Proof. Obviously it is enough to treat the case when $`x_1^\epsilon =\mathrm{}=x_{N_\epsilon }^\epsilon =0`$ for every $`\epsilon `$. We then use an embedding technique that will also play an important role later. Let $`(W_s,s0)`$ be the Brownian snake of \[L2\] driven by the process $`(\beta _s,s0)`$ and with starting point $`\underset{¯}{0}`$. Recall that this is a continuous Markov process with values in $`𝒲_0:=\{w𝒲:w(0)=0\}`$, whose law is characterized by the following properties: $``$ For every $`s0`$, the path $`W_s`$ has lifetime $`\beta _s`$. $``$ Conditionally on $`(\beta _s,s0)`$, the process $`(W_s,s0)`$ is time-inhomogeneous Markov, and its transition kernels are characterized as follows. If $`s<s^{}`$, we have $`W_s^{}(t)=W_s(t)`$ for every $`tm(s,s^{}):=inf_{[s,s^{}]}\beta _r`$, and $`(W_s^{}(m(s,s^{})+r)W_s^{}(m(s,s^{})),0r\beta _s^{}m(s,s^{}))`$ is a Brownian path independent of $`W_s`$. Now, for every $`\epsilon (0,1]`$, we may assume that the spatial motions of the particles are chosen in such a way that, for every $`\epsilon >0`$ and every $`k\{0,1,\mathrm{},\tau ^\epsilon /\epsilon ^2\}`$, $$\begin{array}{cc}\hfill W_{k\epsilon ^2}^\epsilon & =W_{T_k^\epsilon }\text{ if }k\text{ is odd},\hfill \\ \hfill W_{k\epsilon ^2}^\epsilon & =W_{T_k^\epsilon }[0,\beta _{T_k^\epsilon }2\epsilon ]\text{ if }k\text{ is even},\hfill \end{array}$$ $`(2.6)`$ where the notation $`W_{T_k^\epsilon }[0,\beta _{T_k^\epsilon }2\epsilon ]`$ means that the path $`W_{T_k^\epsilon }`$ is restricted to the interval $`[0,\beta _{T_k^\epsilon }2\epsilon ]=[0,\beta _{k\epsilon ^2}^\epsilon ]`$. In fact, it is immediate to verify that the process $`(W_{k\epsilon ^2}^\epsilon ,0k\tau ^\epsilon /\epsilon ^2)`$ defined by (2.6) has (conditionally on $`\beta ^\epsilon `$) the distribution described at the end of Subsection 2.4. Note that the family $`(\tau ^\epsilon ,\epsilon (0,1])`$ is bounded a.s. Then the proof of Lemma 2.2 reduces to checking that, for every $`K>0`$, $$\underset{\delta 0}{lim}P\left[\|W_s(t+r)W_s(t)\|r^{\frac{1}{2}\eta },\text{for every }t0,r[0,\delta ],s[0,K]\right]=1.$$ $`(2.7)`$ This can be easily done using Borel-Cantelli type arguments. Alternatively, we may also use the relations between super-Brownian motion and the Brownian snake \[L2\], and the uniform modulus of continuity of \[DP\]. The graph $`𝒢^\epsilon `$ of the $`\epsilon `$-system of branching particles is defined by $$𝒢^\epsilon =\mathrm{cl}\left(\underset{t0}{}\{t\}\times \mathrm{supp}X_t^\epsilon \right)=\{W_s^\epsilon (t):s[0,\tau ^\epsilon ],0t\beta _s^\epsilon \}.$$ We are interested in weak convergence of $`𝒢^\epsilon `$ towards the graph $`𝒢`$ of $`X`$, which we define as $$𝒢=\mathrm{cl}\left(\underset{t0}{}\{t\}\times \mathrm{supp}X_t\right).$$ We view both $`𝒢^\epsilon `$ and $`𝒢`$ as random elements of the space of all compact subsets of $`R_+\times R`$, which is equipped with the Hausdorff metric. Lemma 2.3. We have the joint convergence $$((X_t^\epsilon ,t0),𝒢^\epsilon )\underset{\epsilon 0}{\overset{(\mathrm{d})}{}}((X_t,t0),𝒢).$$ Proof. We first consider the case when $`x_1^\epsilon =\mathrm{}=x_{N_\epsilon }^\epsilon =0`$ for every $`\epsilon `$. Then we can suppose that the processes $`(W_s^\epsilon ,s[0,\tau ^\epsilon ])`$ are constructed via the embedding technique described in the preceding proof. From (2.1) and (2.6), we get $$(W_{s\tau ^\epsilon }^\epsilon ,s0)\underset{\epsilon 0}{\overset{(\mathrm{a}.\mathrm{s}.)}{}}(W_{s\tau },s0)$$ $`(2.8)`$ in the sense of uniform convergence. Using Lemma 2.1, we get $$X_t^\epsilon =_0^{\tau ^\epsilon }𝑑L_s^{\epsilon ,t}\delta _{W_s^\epsilon (t)}\underset{\epsilon 0}{\overset{(\mathrm{a}.\mathrm{s}.)}{}}_0^\tau 𝑑L_s^t\delta _{W_s(t)}=X_t$$ uniformly in $`t`$. (The formula for $`X_t`$ is the Brownian snake representation of super-Brownian motion, see \[L2\].) Furthermore, the convergence (2.8) also implies that $$𝒢^\epsilon =\{W_s^\epsilon (t):s\tau ^\epsilon ,t\beta _s^\epsilon \}\underset{\epsilon 0}{\overset{(\mathrm{a}.\mathrm{s}.)}{}}\{W_s(t):s\tau ,t\beta _s\},$$ and the limit is easily identified with the graph $`𝒢`$ of $`X`$. Therefore we get the statement of the lemma in the special case $`x_1^\epsilon =\mathrm{}=x_{N_\epsilon }^\epsilon =0`$. Before proceeding to the general case, let us make one more observation. Fix $`\delta >0`$ and write $`(V_s^\epsilon ,s0)`$ for a process distributed as an excursion of $`W^\epsilon `$ away from $`\underset{¯}{0}`$ conditioned to have height greater than $`\delta `$. (Alternatively, $`(V_s^\epsilon ,s0)`$ codes the historical paths of the $`\epsilon `$-system starting with one particle at the origin and conditioned to be non-extinct at time $`\delta `$.) It follows from the convergence (2.8) that we have also $$(V_s^\epsilon ,s0)\underset{\epsilon 0}{\overset{(\mathrm{d})}{}}(V_s,s0),$$ where the limiting process is an excursion of $`W`$ conditioned to have height greater than $`\delta `$. As in the first part of the proof, it follows that the graphs of $`V^\epsilon `$ (defined analogously to $`𝒢^\epsilon `$) also converge in distribution towards the graph of $`V`$. Furthermore, this convergence holds jointly with that of the measure-valued processes $`𝒳_t^\epsilon `$ associated with $`V^\epsilon `$ in the same way as $`X_t^\epsilon `$ was associated with $`W^\epsilon `$. Let us consider now the general case. Because of Lemma 2.2 and assumption (1.2), it is enough to prove that for any fixed $`\delta >0`$, $`𝒢_\epsilon ([\delta ,\mathrm{})\times R)`$ converges in distribution to $`𝒢([\delta ,\mathrm{})\times R)`$ (and that this convergence holds jointly with that of $`X^\epsilon `$). Let $`A_\epsilon `$ stand for the set of indices $`j\{1,\mathrm{},N_\epsilon \}`$ such that the $`j`$-th excursion of $`\beta ^\epsilon `$ has a height greater than $`\delta `$. Note that the events $`\{jA_\epsilon \}`$ are independent with the same probability $`2\epsilon (2\epsilon +\delta )^1`$. It follows that the random measure $$\underset{jA_\epsilon }{}\delta _{x_j^\epsilon }$$ converges weakly to a Poisson measure with intensity $`\frac{2}{\delta }\mu `$. Note that, conditionally on $`A_\epsilon `$, $`𝒢_\epsilon [\delta ,\mathrm{})\times R`$ has the same distribution as $$\underset{jA_\epsilon }{}\left(((0,x_j^\epsilon )+𝒢_{j,\epsilon })[\delta ,\mathrm{})\times R\right)$$ where $`𝒢_{j,\epsilon }`$ are independent copies of the graph of $`V^\epsilon `$. If follows that the random sets $`𝒢_\epsilon [\delta ,\mathrm{})\times R`$ converge in distribution to $$\underset{jJ}{}\left(((0,x_j)+𝒢_{(j)})[\delta ,\mathrm{})\times R\right),$$ where $`_{jJ}\delta _{x_j}`$ is a Poisson point measure on $`R`$ with intensity $`\frac{2}{\delta }\mu `$, and, conditionally on this random measure, the random sets $`𝒢_{(j)}`$ are independent and distributed according to the law of the graph of $`V`$. The canonical representation of superprocesses allows us to identify this limiting distribution with that of $`𝒢([\delta ,\mathrm{})\times R)`$. Furthermore, using the joint convergence of $`(V^\epsilon ,𝒳^\epsilon )`$, it is easy to verify that the convergence holds jointly with that of $`X_\epsilon `$. 3. Branching particle systems with reflection 3.1 Reflection for deterministic paths The purpose of this section is to explain, first in a deterministic setting, the construction of reflected systems. We consider a deterministic branching particle system in $`R`$ analogous to the ones considered above. At time $`0`$, we have $`N`$ particles located at $`x_1,\mathrm{},x_N`$. Each particle moves in $`R`$ and gives birth at its death to $`0`$ or $`2`$ new particles. As in Subsection 2.2, denote by $`𝒯`$ the genealogical forest of the population, which is a subset of $`\{1,\mathrm{},N\}\times 𝐔`$. Each element $`v=(k,u)`$ in $`𝒯`$ corresponds to a particle with birth time $`\xi _v`$ and death time $`\zeta _v`$ (as in Section 2, we could alternatively consider the life durations $`\mathrm{}_v:=\zeta _v\xi _v`$ but in this subsection and the next one it is more convenient to deal with the birth and death times). The spatial motion of $`v`$ is a continuous function $`f_v:[\xi _v,\zeta _v]R`$ and $`f_v^{}(\xi _v^{})=f_v(\zeta _v)`$ if $`v^{}`$ is a child of $`v`$ (then $`\xi _v^{}=\zeta _v`$). The historical path of $`v`$ is the continuous function $`w_v:[0,\zeta _v]R`$ such that, for every $`t[0,\zeta _v)`$, $`w_v(t)`$ is the position at time $`t`$ of the ancestor of $`v`$ alive at that time. We assume that the death times $`\zeta _v`$, $`v𝒯`$ are all distinct, that the system becomes extinct after a finite number of generations and that when a particle dies there is no other particle at the same location: For every $`v𝒯`$, $`f_v(\zeta _v)f_v^{}(\zeta _v)`$ for every $`v^{}𝒯`$ such that $`\xi _v^{}\zeta _v<\zeta _v^{}`$. We turn to the construction of the reflected system. This system is such that the number and positions of the particles alive at every time $`t`$ are the same as in the original system (thus each death time for the reflected system is also a death time for the reflected system). However the genealogical forest $`\stackrel{~}{𝒯}`$ will be different, as will be the spatial motions $`\stackrel{~}{f}_u,u\stackrel{~}{𝒯}`$ or the birth and death times $`\stackrel{~}{\xi }_u,\stackrel{~}{\zeta }_u,u\stackrel{~}{𝒯}`$. Set $`R_0=0`$ and denote by $`R_1<R_2<\mathrm{}<R_M`$ the successive death times in the original system. For every $`k\{1,\mathrm{},M\}`$, let $`𝒯_{(k)}`$ be the set of (labels of) particles that are alive on the interval $`[R_{k1},R_k)`$. We use induction on $`k`$ to define sets $`\stackrel{~}{𝒯}_{(k)}`$, which will represent the particles alive on the interval $`[R_{k1},R_k)`$ in the reflected system, and the corresponding spatial motions. To begin with, we have $`\stackrel{~}{𝒯}_{(1)}=\{1,\mathrm{},N\}`$, and we define $`\stackrel{~}{f}_j(t)`$ for every $`t[0,R_1]`$ and every $`j\stackrel{~}{𝒯}_{(1)}`$ by requiring $`(\stackrel{~}{f}_1(t),\mathrm{},\stackrel{~}{f}_N(t))`$ to be the increasing rearrangement of $`(f_1(t),\mathrm{},f_N(t))`$. Note that the mappings $`\stackrel{~}{f}_1,\mathrm{},\stackrel{~}{f}_N`$ are continuous. Suppose that for some $`k\{1,\mathrm{},M1\}`$, we have defined $`\stackrel{~}{𝒯}_{(k)}`$ and the corresponding paths $`(\stackrel{~}{f}_u(t),t[R_{k1},R_k])`$, for $`u\stackrel{~}{𝒯}_{(k)}`$, in such a way that $`Card\stackrel{~}{𝒯}_{(k)}=Card𝒯_{(k)}`$, and, for every $`t[R_{k1},R_k]`$: $``$ The mapping $`\stackrel{~}{𝒯}_{(k)}u\stackrel{~}{f}_u(t)`$ is increasing with respect to the lexicographical order on $`\stackrel{~}{𝒯}_{(k)}`$. $``$ The values of $`\stackrel{~}{f}_u(t)`$ for $`u\stackrel{~}{𝒯}_{(k)}`$ (counted with their multiplicities) are the same as those of $`f_u(t)`$ for $`u𝒯_{(k)}`$. By definition, one of the particles in $`𝒯_{(k)}`$, say $`u_{(k)}`$, dies at time $`R_k`$. Then there is exactly one $`\stackrel{~}{u}_{(k)}\stackrel{~}{𝒯}_{(k)}`$ such that $`\stackrel{~}{f}_{\stackrel{~}{u}_{(k)}}(R_k)=f_{u_{(k)}}(R_k)`$. We set $$\stackrel{~}{𝒯}_{(k+1)}=\left(\stackrel{~}{𝒯}_{(k)}\backslash \{\stackrel{~}{u}_{(k)}\}\right)\{\stackrel{~}{u}_{(k)}1,\stackrel{~}{u}_{(k)}2\}$$ if $`u_{(k)}`$ has two children in the original system, and $$\stackrel{~}{𝒯}_{(k+1)}=\stackrel{~}{𝒯}_{(k)}\backslash \{\stackrel{~}{u}_{(k)}\}$$ if not. Furthermore, let $`u_1^{k+1},\mathrm{},u_{N_{k+1}}^{k+1}`$ be the elements of $`\stackrel{~}{𝒯}_{(k+1)}`$ listed in lexicographical order. We define $`\stackrel{~}{f}_u(t)`$ for every $`t[R_k,R_{k+1}]`$ and every $`u\stackrel{~}{𝒯}_{(k+1)}`$ by requiring that $`(\stackrel{~}{f}_{u_1^{k+1}}(t),\mathrm{},\stackrel{~}{f}_{u_{N_{k+1}}^{k+1}}(t))`$ is the increasing rearrangement of $`(f_u(t),u𝒯_{(k+1)})`$. Notice that when $`u\stackrel{~}{𝒯}_{(k)}\stackrel{~}{𝒯}_{(k+1)}`$ the definition of $`\stackrel{~}{f}_u(R_k)`$ is consistent with the previous step. Finally, the genealogical forest of the reflected system is $$\stackrel{~}{𝒯}=\underset{k=1}{\overset{M}{}}\stackrel{~}{𝒯}_{(k)}.$$ The birth and death times $`\stackrel{~}{\xi }_u,\stackrel{~}{\zeta }_u`$ as well as the (continuous) spatial motions $`\stackrel{~}{f}_u`$ in the reflected system are defined by the requirement of consistency with the construction of $`\stackrel{~}{𝒯}_{(k)}`$’s. Note the two fundamental properties: $``$ At each time $`t0`$, the positions of the particles (counted with their multiplicities) are the same in the original and the reflected system. $``$ If $`u,v\stackrel{~}{𝒯}`$ with $`uv`$ ($``$ denotes the lexicographical order) then $`f_u(t)f_v(t)`$ for every $`t[\stackrel{~}{\xi }_u,\stackrel{~}{\zeta }_u][\stackrel{~}{\xi }_v,\stackrel{~}{\zeta }_v]`$. Historical paths $`\stackrel{~}{w}_u`$, $`u\stackrel{~}{𝒯}`$ for the reflected system are defined in a way analogous to the original one. If $`u,v\stackrel{~}{𝒯}`$ and $`uv`$ then $`\stackrel{~}{w}_u(t)\stackrel{~}{w}_v(t)`$ for every $`t[0,\stackrel{~}{\zeta }_u\stackrel{~}{\zeta }_v]`$. 3.2 A technical lemma Let $`M\{1,\mathrm{},N\}`$, and consider a branching system consisting only of the particles labeled $`1,\mathrm{},M`$ at time $`0`$ and their descendants. The new genealogy is described by the forest $$𝒯^{}:=𝒯(\{1,\mathrm{},M\}\times U).$$ From this new branching particle system, we can construct a reflected system by the procedure described in Subsection 3.1. We denote by $`\stackrel{~}{𝒯}^{}`$ the genealogical forest for this new reflected system, and by $`\stackrel{~}{w}_v^{}`$, $`v\stackrel{~}{𝒯}^{}`$ the associated historical paths. In general, the historical paths $`\stackrel{~}{w}_v^{}`$ will be very different from those obtained by reflecting the original system. Under special assumptions however, we can say that some of the paths $`\stackrel{~}{w}_v^{}`$ will also be (reflected) historical paths in the original system. Lemma 3.1. Let $`t>0`$ and let $`I`$ be a bounded interval in $`R`$. Suppose that $`w_v(r)I`$ for every $`v𝒯\backslash 𝒯^{}`$ and $`r[0,t]`$. If $`v\stackrel{~}{𝒯}`$ is such that $`\stackrel{~}{\zeta }_vt`$ and $`\stackrel{~}{w}_v(r)I`$ for every $`r[0,t]`$, then there exists $`v^{}\stackrel{~}{𝒯}^{}`$ such that $`\stackrel{~}{\zeta }_v^{}^{}t`$ and $`\stackrel{~}{w}_v^{}^{}(r)=\stackrel{~}{w}_v(r)`$ for every $`r[0,t]`$. The converse also holds: If $`v^{}\stackrel{~}{𝒯}^{}`$ is such that $`\stackrel{~}{\zeta }_v^{}^{}t`$ and $`\stackrel{~}{w}_v^{}^{}(r)I`$ for every $`r[0,t]`$, then there exists $`v\stackrel{~}{𝒯}`$ such that $`\stackrel{~}{\zeta }_vt`$ and $`\stackrel{~}{w}_v(r)=\stackrel{~}{w}_v^{}^{}(r)`$ for every $`r[0,t]`$. In other words, the first assertion means that the path $`\stackrel{~}{w}_v`$, or rather its restriction to $`[0,t]`$, will still be a historical path for the new reflected system. We leave an easy proof of the lemma to the reader. 3.3 Reflected branching particle systems For every $`\epsilon (0,1]`$, we can apply the construction of Subsection 3.1 to the $`\epsilon `$-system of branching Brownian motions. Note that the assumptions that we imposed on the deterministic system hold with probability one for this random system. We write $`𝒯_\epsilon `$ for the genealogical forest of the $`\epsilon `$-system, and $`(\mathrm{}_u^\epsilon ,u𝒯_\epsilon )`$ for the lifetimes of particles. The notation $`\stackrel{~}{𝒯}_\epsilon `$ and $`(\stackrel{~}{\mathrm{}}_u^\epsilon ,u\stackrel{~}{𝒯}_\epsilon )`$ has a similar meaning for the corresponding reflected system, which we call the $`\epsilon `$-reflected system. Observe that $`(𝒯_\epsilon ,(\mathrm{}_u^\epsilon ,u𝒯_\epsilon ))`$ and $`(\stackrel{~}{𝒯}_\epsilon ,(\stackrel{~}{\mathrm{}}_u^\epsilon ,u\stackrel{~}{𝒯}_\epsilon ))`$ have the same distribution. This is so because the spatial motions and branching structure for the $`\epsilon `$-system of branching Brownian motions are independent (a tedious rigorous justification could be given, but we feel that the result is sufficiently obvious to allow us to omit it). Furthermore, $`Card\stackrel{~}{𝒯}_\epsilon =Card𝒯_\epsilon `$. We noticed at the end of Subsection 2.2 that the process $`(\beta _s^\epsilon ,s[0,\tau ^\epsilon ])`$ can be reconstructed as a measurable function of the marked trees $`(𝒯_\epsilon ,(\mathrm{}_u^\epsilon ,u𝒯_\epsilon ))`$. Hence, we can also code the branching structure of the $`\epsilon `$-reflected system by a random process $`(\stackrel{~}{\beta }_s^\epsilon ,s[0,\stackrel{~}{\tau }^\epsilon ])`$ which has the same distribution as $`(\beta _s^\epsilon ,s[0,\tau ^\epsilon ])`$. The fact that $`Card\stackrel{~}{𝒯}_\epsilon =Card𝒯_\epsilon `$ implies that the time $`\tau _\epsilon `$ is also the end of the $`N_\epsilon `$-th excursion of $`\stackrel{~}{\beta }^\epsilon `$ away from $`0`$, and thus $`\stackrel{~}{\tau }^\epsilon =\tau ^\epsilon `$. The discrete local times of $`\stackrel{~}{\beta }^\epsilon `$ (cf. Subsection 2.3) are denoted by $`(\stackrel{~}{L}_s^{\epsilon ,x},xR_+,s[0,\tau ^\epsilon ])`$. Finally, we can code the historical paths of the $`\epsilon `$-reflected system by a discrete snake $`(\stackrel{~}{W}_s^\epsilon ,s[0,\tau ^\epsilon ])`$ in a way analogous to what we did in Subsection 2.4. Recall that we assume $`x_1^\epsilon \mathrm{}x_{N_\epsilon }^\epsilon `$. As in Section 2, if $`s\epsilon ^2N[0,\tau ^\epsilon )`$ and $`\stackrel{~}{\beta }_s^\epsilon =0`$, we set $`\stackrel{~}{W}_s^\epsilon =\underset{¯}{x}_k^\epsilon `$ if $`s`$ is the beginning of the $`k`$-th excursion of $`\stackrel{~}{\beta }^\epsilon `$ away from $`0`$ (and $`\stackrel{~}{W}_{\tau _\epsilon }^\epsilon =\underset{¯}{x}_{N_\epsilon }^\epsilon `$). Otherwise, if $`s\epsilon ^2N[0,\tau ^\epsilon )`$ and $`\stackrel{~}{\beta }_s^\epsilon >0`$, then $`(s,\stackrel{~}{\beta }_s^\epsilon )`$ can be associated with a unique edge $`u`$ of the forest $`\stackrel{~}{𝒯}_\epsilon `$, and we let $`\stackrel{~}{W}_s^\epsilon `$ be equal to $`\stackrel{~}{w}_u^\epsilon `$, the historical path of $`u`$. If $`s\epsilon ^2N`$, we use the same interpolation as in Section 2. A fundamentally important property of the process $`(\stackrel{~}{W}_s^\epsilon ,s[0,\tau ^\epsilon ])`$, from the point of view of our project, is that for $`s<s^{}`$, $$\stackrel{~}{W}_s^\epsilon (t)\stackrel{~}{W}_s^{}^\epsilon (t),t[0,\stackrel{~}{\beta }_s^\epsilon \stackrel{~}{\beta }_s^{}^\epsilon ].$$ $`(3.1)`$ This follows from our construction and the end of Subsection 3.1. As in the case of $`W_s^\epsilon `$, we see that if $`s<s^{}`$, then $$\stackrel{~}{W}_s^\epsilon (t)=\stackrel{~}{W}_s^{}^\epsilon (t),t[0,\underset{u[s,s^{}]}{inf}\stackrel{~}{\beta }_u^\epsilon ].$$ Because at every time the locations of particles are the same in the reflected system and in the original one, the random measure $$\stackrel{~}{X}_t^\epsilon =_0^{\tau _\epsilon }𝑑\stackrel{~}{L}_s^{\epsilon ,t}\delta _{\stackrel{~}{W}_s^\epsilon (t)}$$ coincides with $`X_t^\epsilon `$. Things are however very different for the historical measure $$\stackrel{~}{Y}_t^\epsilon =_0^{\tau _\epsilon }𝑑\stackrel{~}{L}_s^{\epsilon ,t}\delta _{\stackrel{~}{W}_s^\epsilon }.$$ 4. Tightness of the reflected system 4.1 Uniform continuity of the reflected paths Our first goal is to derive an important uniform continuity property for the individual paths of the $`\epsilon `$-reflected system (Theorem 4.1 below). From the intuitive point of view, reflected paths should have smaller oscillations than “free” paths and so this property seems to be a straightforward consequence of Lemma 2.2. However the intuition about the relationship between moduli of continuity of free and reflected paths is only correct as long as we do not have any deaths. To be specific consider $`p`$ paths $`w_{(1)},\mathrm{},w_{(p)}`$ all defined on the time interval $`[0,1]`$, and let $`\stackrel{~}{w}_{(1)},\mathrm{},\stackrel{~}{w}_{(p)}`$ be the corresponding system of reflected paths. Then, if we assume that $`\|w_{(i)}(t)w_{(i)}(t^{})\|\phi (\|tt^{}\|)`$ for every $`i=1,\mathrm{},p`$ and $`t,t^{}[0,1]`$ and for some nondecreasing function $`\phi `$, an easy argument shows that the same bound holds when the paths $`w_{(i)}`$ are replaced by $`\stackrel{~}{w}_{(i)}`$. It turns out that a similar assertion about moduli of continuity is false if paths may have different lifetimes. Fig. 2 shows a system of two paths. In the original system, the oscillations of paths over the intervals where they are defined are equal to $`z_1y_1`$ and $`z_2y_2`$. One of the paths in the reflected system goes from $`y_1`$ to $`z_2`$ and so has an oscillation larger than the oscillations of the original paths. In this article we consider Brownian particles which die at different times so we cannot use known estimates for the modulus of continuity of the original (non-reflecting) historical paths in a direct way. We will use them later in a different but quite essential way. Figure 2. Recall our notation $`\stackrel{~}{w}_u^\epsilon `$, $`u\stackrel{~}{𝒯}_\epsilon `$, for the historical paths of the $`\epsilon `$-reflected system. By convention, $`\stackrel{~}{w}_u^\epsilon (t)=\stackrel{~}{w}_u^\epsilon (\stackrel{~}{\zeta }_u^\epsilon )`$ if $`t\stackrel{~}{\zeta }_u^\epsilon `$. Theorem 4.1. For every $`\eta >0`$, $$\underset{\delta 0}{lim}\left(\underset{\epsilon 0}{lim\; sup}P[\underset{\|tt^{}\|\delta }{\underset{t,t^{}0}{sup}}\underset{u\stackrel{~}{𝒯}_\epsilon }{sup}\|\stackrel{~}{w}_u^\epsilon (t)\stackrel{~}{w}_u^\epsilon (t^{})\|>\eta ]\right)=0.$$ Proof. Let $`(\delta _{(p)},\epsilon _{(p)})`$ be a sequence in $`(0,1]^2`$ converging to $`0`$. We will prove that there exists a subsequence $`(\delta _{(p)}^{},\epsilon _{(p)}^{})`$ such that: $$\underset{p\mathrm{}}{lim}\left(\underset{\|tt^{}\|\delta _{(p)}^{}}{\underset{t,t^{}0}{sup}}\underset{u\stackrel{~}{𝒯}_{\epsilon _{(p)}^{}}}{sup}\|\stackrel{~}{w}_u^{\epsilon _{(p)}^{}}(t)\stackrel{~}{w}_u^{\epsilon _{(p)}^{}}(t^{})\|\right)=0$$ $`(4.1)`$ in probability. Clearly, the statement of Theorem 4.1 is a consequence of this fact. We first explain how we choose the sequence $`(\delta _{(p)}^{},\epsilon _{(p)}^{})`$. By Lemma 2.3 and the Skorohod representation theorem (\[EK\] Theorem 3.1.8), we may, for every $`p1`$, replace the pair $`(X^{\epsilon _{(p)}},𝒢_{\epsilon _{(p)}})`$ by a new pair with the same distribution (for which we keep the same notation), in such a way that $$(X^{\epsilon _{(p)}},𝒢_{\epsilon _{(p)}})\underset{p\mathrm{}}{\overset{(\mathrm{a}.\mathrm{s}.)}{}}(X,𝒢),$$ where $`X`$ is a super-Brownian motion started at $`\mu `$ and $`𝒢`$ denotes its graph. Note that the genealogical forest $`\stackrel{~}{𝒯}_{\epsilon _{(p)}}`$, the process $`\stackrel{~}{\beta }^{\epsilon _{(p)}}`$ and the historical paths $`\stackrel{~}{w}_u^{\epsilon _{(p)}}`$, $`u\stackrel{~}{𝒯}_{\epsilon _{(p)}}`$, are reconstructed as measurable functions of the new process $`X^{(\epsilon _p)}`$, and that it suffices to prove (4.1) for the new historical paths. As a consequence of the remark following Lemma 2.1, we have $$\underset{p\mathrm{}}{lim}\left(\underset{s0}{sup}\underset{\|tt^{}\|\delta _{(p)}}{\underset{t,t^{}0}{sup}}\|\stackrel{~}{L}_{s\tau ^{\epsilon _{(p)}}}^{\epsilon _{(p)},t}\stackrel{~}{L}_{s\tau ^{\epsilon _{(p)}}}^{\epsilon _{(p)},t^{}}\|\right)=0$$ $`(4.2)`$ in probability. We choose the subsequence $`(\delta _{(p)}^{},\epsilon _{(p)}^{})`$ so that the convergence (4.2) holds almost surely along this subsequence. We will argue by contradiction to prove (4.1). If (4.1) does not hold, then on a set $`A`$ of positive probability, we can find a number $`\eta >0`$ and a (random) subsequence $`p_k\mathrm{}`$ such that, if $`\epsilon _k:=\epsilon _{(p_k)}^{}`$ and $`\delta _k:=\delta _{(p_k)}^{}`$, $$\underset{\|tt^{}\|\delta _k}{\underset{t,t^{}0}{sup}}\underset{u\stackrel{~}{𝒯}_{\epsilon _k}}{sup}\|\stackrel{~}{w}_u^{\epsilon _k}(t)\stackrel{~}{w}_u^{\epsilon _k}(t^{})\|>\eta .$$ $`(4.3)`$ From now on until the end of the proof, we will assume that the event $`A`$ holds. By (4.3), for every $`k1`$, there exist $`u_k\stackrel{~}{𝒯}_{\epsilon _k}`$, $`t_k,t_k^{}0`$ with $`\|t_kt_k^{}\|\delta _k`$, such that $$\|\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k)\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k^{})\|>\eta .$$ Clearly, we can assume that $`t_kt_k^{}\stackrel{~}{\zeta }_{u_k}^{\epsilon _k}`$. Recall that the graphs $`𝒢_{\epsilon _{(p)}}`$ converge to $`𝒢`$ in the Hausdorff metric. In particular, the set of all pairs $`(t_k,\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k))`$ and $`(t_k^{},\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k^{}))`$ is relatively compact. By passing to a subsequence, if necessary, we may assume that $`t_k,t_k^{}t_{\mathrm{}}`$, $`\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k)x_1`$ and $`\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k^{})x_2`$ as $`k\mathrm{}`$. We have $`\|x_1x_2\|\eta `$, and we take $`x_2>x_1`$ for definiteness. We also know that $`\stackrel{~}{X}_t^{\epsilon _{(p)}}=X_t^{\epsilon _{(p)}}`$ converges to $`X_t`$ a.s. as $`p\mathrm{}`$, uniformly on compact subsets of $`R_+`$. Hence, both sequences $`\stackrel{~}{X}_{t_k}^{\epsilon _k}`$ and $`\stackrel{~}{X}_{t_k^{}}^{\epsilon _k}`$ converge to $`X_t_{\mathrm{}}`$, and $$\begin{array}{cc}& \underset{k\mathrm{}}{lim\; sup}\stackrel{~}{X}_{t_k}^{\epsilon _k}((\mathrm{},\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k)])X_t_{\mathrm{}}((\mathrm{},x_1]),\hfill \\ & \underset{k\mathrm{}}{lim\; inf}\stackrel{~}{X}_{t_k^{}}^{\epsilon _k}((\mathrm{},\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k^{})))X_t_{\mathrm{}}((\mathrm{},x_2)).\hfill \end{array}$$ $`(4.4)`$ We claim that $$\underset{k\mathrm{}}{lim\; inf}\left(\stackrel{~}{X}_{t_k}^{\epsilon _k}\left((\mathrm{},\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k)]\right)\stackrel{~}{X}_{t_k^{}}^{\epsilon _k}\left((\mathrm{},\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k^{}))\right)\right)0.$$ $`(4.5)`$ To see this, we use the discrete snake representation of Subsection 3.3. Write $`s_k[0,\tau _{\epsilon _k})\epsilon ^2N`$ for the time associated with the edge $`u_k`$ of $`\stackrel{~}{𝒯}_{\epsilon _k}`$ in this representation. By construction, $`\stackrel{~}{w}_{u_k}^{\epsilon _k}=\stackrel{~}{W}_{s_k}^{\epsilon _k}`$, and (3.1) implies $$\stackrel{~}{X}_{t_k}^{\epsilon _k}((\mathrm{},\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k)])=_0^{\tau _{\epsilon _k}}𝑑\stackrel{~}{L}_s^{\epsilon _k,t_k}\mathbf{\hspace{0.17em}1}_{\{\stackrel{~}{W}_s^{\epsilon _k}(t_k)\stackrel{~}{W}_{s_k}^{\epsilon _k}(t_k)\}}\stackrel{~}{L}_{s_k}^{\epsilon _k,t_k}.$$ Similarly, we get $$\stackrel{~}{X}_{t_k^{}}^{\epsilon _k}((\mathrm{},\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k^{})))\stackrel{~}{L}_{s_k}^{\epsilon _k,t_k^{}}.$$ Hence, $$\stackrel{~}{X}_{t_k}^{\epsilon _k}((\mathrm{},\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k)])\stackrel{~}{X}_{t_k^{}}^{\epsilon _k}((\mathrm{},\stackrel{~}{w}_{u_k}^{\epsilon _k}(t_k^{})))\stackrel{~}{L}_{s_k}^{\epsilon _k,t_k}\stackrel{~}{L}_{s_k}^{\epsilon _k,t_k^{}}.$$ $`(4.6)`$ On the other hand, $$\|\stackrel{~}{L}_{s_k}^{\epsilon _k,t_k}\stackrel{~}{L}_{s_k}^{\epsilon _k,t_k^{}}\|\underset{s[0,\tau _{\epsilon _k}]}{sup}\underset{\|tt^{}\|\delta _k}{\underset{t,t^{}0}{sup}}\|\stackrel{~}{L}_s^{\epsilon _k,t}\stackrel{~}{L}_s^{\epsilon _k,t^{}}\|.$$ By the convergence in (4.2), which holds a.s. along the subsequence $`(\delta _{(p)}^{},\epsilon _{(p)}^{})`$, the right hand side tends to $`0`$ as $`k\mathrm{}`$. This and (4.6) give the claim (4.5). From (4.5) and (4.4), we get $`X_t_{\mathrm{}}((\mathrm{},x_1])X_t_{\mathrm{}}((\mathrm{},x_2))`$ and thus (recall that $`x_1<x_2`$), $`X_t_{\mathrm{}}((x_1,x_2))=0`$. This a priori does not imply that $`\{t_{\mathrm{}}\}\times (x_1,x_2)𝒢=\mathrm{}`$ as there could be a “local extinction” of $`X`$ at time $`t_{\mathrm{}}`$ in $`(x_1,x_2)`$. However, by Theorem 1.4 of Perkins \[P1\], there can be at most one local extinction at a given time, so we can choose $`x_1^{}`$ and $`x_2^{}`$ with $`x_1<x_1^{}<x_2^{}<x_2`$ such that $`\{t_{\mathrm{}}\}\times [x_1^{},x_2^{}]𝒢=\mathrm{}`$. Since $`𝒢`$ is closed, we have also $`[t_{\mathrm{}}\delta ,t_{\mathrm{}}+\delta ]\times [x_1^{},x_2^{}]𝒢=\mathrm{}`$ for $`\delta >0`$ sufficiently small. However, by construction, for $`k`$ sufficiently large the paths $`\stackrel{~}{w}_{u_k}^{\epsilon _k}`$ and thus also the graph $`𝒢_{\epsilon _k}`$ must intersect $`[t_{\mathrm{}}\delta ,t_{\mathrm{}}+\delta ]\times [x_1^{},x_2^{}]`$. This gives a contradiction since we know that $`𝒢_{\epsilon _k}`$ converge to $`𝒢`$. This contradiction completes the proof of Theorem 4.1. 4.2 Tightness of reflected discrete snakes From now on, we restrict our attention to values of $`\epsilon `$ belonging to a fixed sequence $``$ decreasing to $`0`$. For convenience, we extend the definition of the discrete snakes $`\stackrel{~}{W}^\epsilon `$ by taking $`\stackrel{~}{W}_s^\epsilon =\stackrel{~}{W}_{\tau ^\epsilon }^\epsilon =\underset{¯}{x}_{N_\epsilon }^\epsilon `$ (and thus $`\stackrel{~}{\beta }_s^\epsilon =0`$) for $`s>\tau ^\epsilon `$. Proposition 4.2. The laws of the processes $`\stackrel{~}{W}^\epsilon `$, $`\epsilon `$, are tight in the space of all probability measures on $`𝐃([0,\mathrm{}),𝒲)`$. Furthermore, if $`(\stackrel{~}{W}_s,s0)`$ is a weak limit point of this sequence of processes, we have the following properties. (i) If $`\stackrel{~}{\beta }_s:=\zeta _{\stackrel{~}{W}_s}`$, the process $`(\stackrel{~}{\beta }_s,s0)`$ has the same distribution as $`(\beta _{s\tau },s0)`$. (ii) Almost surely for every $`ss^{}`$ we have $`\stackrel{~}{W}_s(t)\stackrel{~}{W}_s^{}(t)`$ for every $`t[0,\stackrel{~}{\beta }_s\stackrel{~}{\beta }_s^{}]`$. (iii) The set of discontinuities of the mapping $`s\stackrel{~}{W}_s`$ is contained in the zero set of $`\stackrel{~}{\beta }`$. Furthermore, if $`s<s^{}`$ belong to the same connected component of the complement of the zero set, we have $$\stackrel{~}{W}_s(t)=\stackrel{~}{W}_s^{}(t)\text{for every }t[0,\underset{r[s,s^{}]}{inf}\stackrel{~}{\beta }_r].$$ Proof. The hard part of the proof is to show tightness. To this end we rely on the classical criteria (see e.g. Corollary 3.7.4 of \[EK\]). We first observe that the compact containment condition is a straightforward consequence of Theorem 4.1. In fact, if $`\eta >0`$ is fixed, then for every integer $`p1`$, Theorem 4.1 and the construction of the discrete snake $`\stackrel{~}{W}^\epsilon `$ allow us to find $`\delta _p>0`$ such that, for $`\epsilon `$ small enough, $$P\left[\underset{s0}{sup}\underset{\|tt^{}\delta _p}{\underset{t,t^{}0}{sup}}\|\stackrel{~}{W}_s^\epsilon (t)\stackrel{~}{W}_s^\epsilon (t^{})\|>2^p\right]\eta \mathrm{\hspace{0.17em}2}^{p1}.$$ $`(4.7)`$ (Here and later, we make the convention that $`\stackrel{~}{W}_s^\epsilon (t)=\stackrel{~}{W}_s^\epsilon (\stackrel{~}{\beta }_s^\epsilon )`$ for $`t>\stackrel{~}{\beta }_s^\epsilon `$.) It is easy to see that an even stronger assertion holds, namely, (4.7) is true for all $`\epsilon `$; this can be achieved by taking $`\delta _p`$ even smaller if necessary—note that for any fixed value of $`\epsilon `$ we need only consider a finite number of historical paths. Then let $`H`$ be a compact subset of $`R_+`$ containing $`\mathrm{supp}\mu _\epsilon `$ for $`\epsilon `$, and let $`A>0`$ be a constant. The set $$\begin{array}{cc}\hfill K:=& \{w𝒲:w(0)H,\zeta _wA,\hfill \\ & \text{and }\|w(t)w(t^{})\|2^p\text{ for every }t,t^{}[0,\zeta _w]\text{ with }\|tt^{}\|\delta _p\text{ and every }p1\}\hfill \end{array}$$ is compact, and it follows from (4.7) that $$P[\stackrel{~}{W}_s^\epsilon K\text{ for some }s0]<\eta $$ provided that $`A`$ is chosen large enough. Recall the definition of the distance $`d`$ from Subsection 2.4. We set $$\theta (\epsilon ,\delta )=\underset{(s_i)}{inf}\left\{\underset{i}{sup}\underset{s,s^{}[s_{i1},s_i)}{sup}d(W_s^\epsilon ,W_s^{}^\epsilon )\right\},$$ where the infimum is over all finite sequences $`0=s_0<s_1<\mathrm{}<s_{m1}<\tau ^\epsilon s_m`$ such that $`inf\{\|s_is_{i1}\|;1im\}\delta \}`$. As a direct application of Corollary 3.7.4 in \[EK\], the proof of tightness will be complete if we can verify that, for every $`\eta >0`$, we can choose $`\delta >0`$ sufficiently small so that $$\underset{\epsilon 0}{lim\; sup}P[\theta (\epsilon ,\delta )>\eta ]<\eta .$$ $`(4.8)`$ We now fix $`\eta >0`$ and proceed to the proof of (4.8). As a consequence of Theorem 4.1, we can choose $`\rho (0,\eta /5)`$ so small that, for every $`\epsilon `$, $$P\left[\underset{s0}{sup}\underset{\|tt^{}\|\rho }{\underset{t,t^{}0}{sup}}\|\stackrel{~}{W}_s^\epsilon (t)\stackrel{~}{W}_s^\epsilon (t^{})\|\frac{\eta }{5}\right]1\frac{\eta }{5}.$$ $`(4.9)`$ Then, by the tightness of the laws of $`\stackrel{~}{\beta }^\epsilon `$ (cf (2.3)), we can choose $`\kappa >0`$ small enough so that, for every $`\epsilon `$, $$P[\underset{\|ss^{}\|\kappa }{\underset{s,s^{}0}{sup}}\|\stackrel{~}{\beta }_s^\epsilon \stackrel{~}{\beta }_s^{}^\epsilon \|\rho ]1\frac{\eta }{5}.$$ $`(4.10)`$ We denote by $`E_\epsilon `$ the intersection of the events considered in (4.9) and (4.10), so that the probability of the complement of $`E_\epsilon `$ is bounded above by $`2\eta /5`$. Set $`\gamma =\eta /5`$. Since $`\mu `$ is a finite measure with compact support, we can easily find an integer $`M_\gamma `$ and a finite sequence of reals $`y_1<z_1y_2<z_2\mathrm{}y_{M_\gamma }<z_{M_\gamma }`$, such that: $``$ $`z_iy_i<\gamma `$ for every $`i=1,\mathrm{},M_\gamma `$, $``$ $`{\displaystyle \underset{i=1}{\overset{M_\gamma }{}}}[y_i,z_i)`$ contains a neighborhood of $`\mathrm{supp}\mu `$, $``$ $`\mu (\{y_i\})=\mu (\{z_i\})=0`$, and $`\mu ([y_i,z_i))>0`$ for every $`i=1,\mathrm{}M_\gamma `$. By the last condition, $`a_\gamma :=inf\{\mu ([y_i,z_i)),i=1,\mathrm{}M_\gamma \}>0`$. Furthermore, if $`\epsilon `$ is small enough, $$\mathrm{supp}\mu _\epsilon \underset{i=1}{\overset{M_\gamma }{}}[y_i,z_i)$$ and $$Card\{j:x_j^\epsilon [y_i,z_i)\}>\frac{a_\gamma }{2\epsilon }1,$$ for every $`i=1,\mathrm{},M_\gamma `$. From now on, we assume that $`\epsilon `$ is small enough so that the last two conditions hold, and we set $$n_i^\epsilon =inf\{j:x_j^\epsilon [y_i,z_i)\},i=1,\mathrm{}M_\gamma .$$ Denote by $`\stackrel{~}{\tau }_k^\epsilon `$ the $`k`$-th return of $`\stackrel{~}{\beta }^\epsilon `$ to the origin. We also set $`\sigma _i^\epsilon :=\stackrel{~}{\tau }_{n_i^\epsilon }^\epsilon `$ and $`\sigma _{M_\gamma +1}^\epsilon :=\stackrel{~}{\tau }_{N_\epsilon }^\epsilon =\tau ^\epsilon `$. Note that each of the variables $`\sigma _{i+1}^\epsilon \sigma _i^\epsilon `$ is bounded below in distribution by $`\stackrel{~}{\tau }_{[a_\gamma /2\epsilon ]}^\epsilon `$, and recall that for every $`c>0`$, $`\stackrel{~}{\tau }_{[c/\epsilon ]}^\epsilon `$ converges in distribution to $`\tau _c`$. Since $`\tau _c>0`$ a.s., we may choose $`\delta (0,\kappa /2)`$ so small that, for $`\epsilon `$ small, $$P[\sigma _{i+1}^\epsilon \sigma _i^\epsilon >2\delta \text{ for every }i\{1,\mathrm{},M_\gamma \}]>1\frac{\eta }{5}.$$ $`(4.11).`$ Write $`E_\epsilon ^{}`$ for the intersection of the set $`E_\epsilon `$ with the event considered in (4.11). Notice that on $`E_\epsilon ^{}`$ we can choose a finite sequence $`0=s_0^\epsilon <s_1^\epsilon <\mathrm{}<s_{K_\epsilon }^\epsilon =\tau ^\epsilon `$ in such a way that $`\delta s_j^\epsilon s_{j1}^\epsilon 2\delta <\kappa `$, for every $`j\{1,\mathrm{},K_\epsilon \}`$, and each interval $`[s_{j1}^\epsilon ,s_j^\epsilon )`$ is contained in exactly one interval $`[\sigma _{k1}^\epsilon ,\sigma _k^\epsilon )`$. We use the sequence $`(s_i^\epsilon )`$ to get an upper bound on $`\theta (\epsilon ,\delta )`$ on the event $`E_\epsilon ^{}`$. First observe that for $`j\{1,\mathrm{},K_\epsilon \}`$, $$\underset{s,s^{}[s_{j1}^\epsilon ,s_j^\epsilon )}{sup}d(\stackrel{~}{W}_s^\epsilon ,\stackrel{~}{W}_s^{}^\epsilon )\underset{s,s^{}[s_{j1}^\epsilon ,s_j^\epsilon )}{sup}\|\stackrel{~}{\beta }_s^\epsilon \stackrel{~}{\beta }_s^{}^\epsilon \|+\underset{s,s^{}[s_{j1}^\epsilon ,s_j^\epsilon )}{sup}\underset{t0}{sup}\|\stackrel{~}{W}_s^\epsilon (t)\stackrel{~}{W}_s^{}^\epsilon (t)\|.$$ The first term on the right hand side is bounded above by $`\rho \eta /5`$ by the definition of $`E_\epsilon `$ (cf (4.10)) and the property $`s_j^\epsilon s_{j1}^\epsilon <\kappa `$. To bound the second term, let $`s,s^{}[s_{j1}^\epsilon ,s_j^\epsilon )`$ and consider first the case when $$m^\epsilon (s,s^{}):=\underset{r[s,s^{}]}{inf}\stackrel{~}{\beta }_r^\epsilon >0.$$ Then $`\stackrel{~}{W}_s^\epsilon (t)=\stackrel{~}{W}_s^{}^\epsilon (t)`$ for every $`t[0,m^\epsilon (s,s^{})]`$, and thus $$\begin{array}{cc}& \underset{t0}{sup}\|\stackrel{~}{W}_s^\epsilon (t)\stackrel{~}{W}_s^{}^\epsilon (t)\|\hfill \\ & \underset{m^\epsilon (s,s^{})t\stackrel{~}{\beta }_s^\epsilon }{sup}\|\stackrel{~}{W}_s^\epsilon (t)\stackrel{~}{W}_s^\epsilon (m^\epsilon (s,s^{}))\|+\underset{m^\epsilon (s,s^{})t\stackrel{~}{\beta }_s^{}^\epsilon }{sup}\|\stackrel{~}{W}_s^{}^\epsilon (t)\stackrel{~}{W}_s^{}^\epsilon (m^\epsilon (s,s^{}))\|\hfill \\ & \frac{2\eta }{5}\hfill \end{array}$$ again by the definition of $`E_\epsilon `$ (cf (4.9) and (4.10)). The case $`m^\epsilon (s,s^{})=0`$ is analogous, but we now get the additional term $`\|\stackrel{~}{W}_s^\epsilon (0)\stackrel{~}{W}_s^{}^\epsilon (0)\|`$. However, by construction, $`s`$ and $`s^{}`$ belong to the same interval $`[\sigma _{k1}^\epsilon ,\sigma _k^\epsilon )`$ and thus $`\stackrel{~}{W}_s^\epsilon (0)`$ and $`\stackrel{~}{W}_s^{}^\epsilon (0)`$ belong to the same $`[y_k,z_k)`$, which implies that $`\|\stackrel{~}{W}_s^\epsilon (0)\stackrel{~}{W}_s^{}^\epsilon (0)\|\gamma =\eta /5`$. Finally, for every $`j\{1,\mathrm{},K_\epsilon \}`$, we get the bound $$\underset{s,s^{}[s_{j1}^\epsilon ,s_j^\epsilon )}{sup}d(\stackrel{~}{W}_s^\epsilon ,\stackrel{~}{W}_s^{}^\epsilon )\frac{4\eta }{5}<\eta $$ on $`E_\epsilon ^{}`$. It follows that, for $`\epsilon `$ small, $$P[\theta (\epsilon ,\delta )\eta ]P[(E_\epsilon ^{})^c]\frac{3\eta }{5}<\eta .$$ This completes the proof of (4.8) and of the tightness of the sequence $`\stackrel{~}{W}^\epsilon `$. The remaining assertions of Proposition 4.2 are easy. (i) is clear since $`\stackrel{~}{\beta }`$ must be the weak limit of $`\stackrel{~}{\beta }^\epsilon `$. (ii) follows from the analogous property for $`\stackrel{~}{W}^\epsilon `$, and a similar argument applies to (iii). 4.3 Tightness of the reflected historical processes Recall that the historical process for the $`\epsilon `$-reflected system is the process with values in $`M_f(𝒲)`$ defined by $$\stackrel{~}{Y}_t^\epsilon =_0^{\tau _\epsilon }𝑑\stackrel{~}{L}_s^{\epsilon ,t}\delta _{\stackrel{~}{W}_s^\epsilon }.$$ It is easy to verify that $`\stackrel{~}{Y}^\epsilon `$ has right-continuous paths with left limits. The following theorem is a slightly more precise version of Theorem 1.1. Theorem 4.3. The sequence of the laws $`\stackrel{~}{}_Y^\epsilon `$ of $`\stackrel{~}{Y}^\epsilon `$, $`\epsilon `$, is tight in the space of probability measures on $`𝐃([0,\mathrm{}),M_f(𝒲))`$ and any limit law is supported on $`𝐂([0,\mathrm{}),M_f(𝒲))`$. Suppose that $`\stackrel{~}{}_Y`$ is the limit of a subsequence of $`\stackrel{~}{}_Y^\epsilon `$. By passing to a further subsequence of $`\epsilon `$’s, if necessary, we may assume that the laws $`\stackrel{~}{}_W^\epsilon `$ of $`\stackrel{~}{W}^\epsilon `$ converge to a law $`\stackrel{~}{}_W`$. Then one can construct on some probability space processes $`\stackrel{~}{Y}`$ and $`\stackrel{~}{W}`$ with distributions $`\stackrel{~}{}_Y`$ and $`\stackrel{~}{}_W`$, resp., related by $$\stackrel{~}{Y}_t=_0^{\stackrel{~}{\tau }}𝑑\stackrel{~}{L}_s^t\delta _{\stackrel{~}{W}_s},$$ where $`(\stackrel{~}{L}_s^t,t0,s0)`$ denote the local times of the process $`\stackrel{~}{\beta }_s:=\zeta _{\stackrel{~}{W}_s}`$, and $`\stackrel{~}{\tau }=inf\{s0:\stackrel{~}{L}_s^0=a\}`$. Proof. By Proposition 4.2, the laws of $`\stackrel{~}{W}^\epsilon `$, $`\epsilon `$ are tight. Hence, from any subsequence of $``$, we can extract a further subsequence $`_0`$ along which $`\stackrel{~}{W}^\epsilon `$ converges in distribution. We can in fact obtain more. For every $`\epsilon >0`$ and $`t0`$, denote by $`\mathrm{\Gamma }_t^\epsilon `$, $`\stackrel{~}{\mathrm{\Gamma }}_t^\epsilon `$ the random measures on $`R_+`$ defined by $$\mathrm{\Gamma }_t^\epsilon ,\phi =_0^{\tau \epsilon }𝑑L_s^{\epsilon ,t}\phi (s),\stackrel{~}{\mathrm{\Gamma }}_t^\epsilon ,\phi =_0^{\stackrel{~}{\tau }\epsilon }𝑑\stackrel{~}{L}_s^{\epsilon ,t}\phi (s).$$ (We have $`\tau ^\epsilon =\stackrel{~}{\tau }^\epsilon `$ but we prefer to keep a different notation here.) Also define $`\mathrm{\Gamma }_t`$ by: $$\mathrm{\Gamma }_t,\phi =_0^\tau 𝑑L_s^t\phi (s).$$ As a consequence of Lemma 2.1, we know that $$(\beta _{\tau ^\epsilon }^\epsilon ,\mathrm{\Gamma }^\epsilon )\underset{\epsilon 0}{}(\beta _\tau ,\mathrm{\Gamma })$$ uniformly on $`[0,\mathrm{})^2`$, a.s. If we replace the pair $`(\beta _{\tau ^\epsilon }^\epsilon ,\mathrm{\Gamma }^\epsilon )`$ by $`(\stackrel{~}{\beta }_{\stackrel{~}{\tau }^\epsilon }^\epsilon ,\stackrel{~}{\mathrm{\Gamma }}^\epsilon )`$ this convergence still holds in distribution in $`𝐂(R_+,R)\times 𝐃(R_+,M_f(R_+))`$. From this observation and standard arguments, we have the joint convergence $$(\stackrel{~}{W}^\epsilon ,\stackrel{~}{\beta }^\epsilon ,\stackrel{~}{\mathrm{\Gamma }}^\epsilon )\underset{\epsilon 0,\epsilon _0}{\overset{(\mathrm{d})}{}}(\stackrel{~}{W},\stackrel{~}{\beta },\stackrel{~}{\mathrm{\Gamma }})$$ $`(4.12)`$ where $$\stackrel{~}{\mathrm{\Gamma }}_t,\phi =_0^{\stackrel{~}{\tau }}𝑑\stackrel{~}{L}_s^t\phi (s),$$ with the notation introduced in the theorem. By the Skorohod representation theorem, we can replace for every $`\epsilon _0`$ the triplet $`(\stackrel{~}{W}^\epsilon ,\stackrel{~}{\beta }^\epsilon ,\stackrel{~}{\mathrm{\Gamma }}^\epsilon )`$ by a new triplet having the same distribution, in such a way that the convergence (4.12) now holds almost surely. Without risk of confusion, we keep the same notation for the new triplets. We claim that we have then $$\stackrel{~}{Y}_t^\epsilon =\stackrel{~}{\mathrm{\Gamma }}_t^\epsilon (ds)\delta _{\stackrel{~}{W}_s^\epsilon }\underset{\epsilon 0,\epsilon _0}{}\stackrel{~}{\mathrm{\Gamma }}_t(ds)\delta _{\stackrel{~}{W}_s}=\stackrel{~}{Y}_t$$ $`(4.13)`$ uniformly on compact subsets of $`R_+`$, a.s. Clearly Theorem 4.3 follows from (4.13) and the fact that the limiting process $`\stackrel{~}{Y}`$ that appears in (4.13) is continuous. Both (4.13) and the latter fact are immediate consequences of the convergence (4.12) (now assumed to hold a.s.) and the following “elementary” lemma, whose proof is left to the reader. Lemma 4.4. Let $`(\gamma ^n,nN)`$ be a sequence in $`𝐃(R_+,M_f(R_+))`$. Assume that $`\gamma _t^n`$ converges as $`n\mathrm{}`$ to $`\gamma _t`$, uniformly on every compact of $`R_+`$, that $`t\gamma _t`$ is continuous and that the measure $`\gamma _t`$ is diffuse, for every $`tR_+`$. Let $`E`$ be a Polish space and let $`(f_n,nN)`$ be a sequence in $`𝐃(R_+,E)`$ that converges to $`f`$ in $`𝐃(R_+,E)`$. For every integer $`nN`$ and every $`tR_+`$, let $`\nu _t^nM_f(E)`$ be defined by $$\nu _t^n=\gamma _t^n(ds)\delta _{f_n(s)}.$$ Then $`\nu _t^n`$ converges as $`n\mathrm{}`$, uniformly on compact subsets of $`R_+`$, to the measure $`\nu _t`$ defined by $$\nu _t=\gamma _t(ds)\delta _{f(s)}.$$ Furthermore, the mapping $`t\nu _t`$ is continuous. Remark. We do not know whether the limit law of the sequence $`\stackrel{~}{}_Y^\epsilon `$ in Theorem 4.3 is unique. A positive answer would give the convergence in distribution of the processes $`\stackrel{~}{Y}^\epsilon `$. We can also formulate the problem in terms of the reflected snake. Is there a unique (in law) process $`\stackrel{~}{W}`$ satisfying properties (i) – (iii) of Proposition 4.2 and such that $$t_0^{\stackrel{~}{\tau }}𝑑\stackrel{~}{L}_s^t\delta _{\stackrel{~}{W}_s(t)}$$ is a super-Brownian motion started at $`\mu `$ ? 5. Path properties of the reflected historical process 5.1 Preliminaries Throughout this section, we consider a process $`\stackrel{~}{Y}`$ which is a weak limit of the processes $`\stackrel{~}{Y}^\epsilon `$ as $`\epsilon 0`$. According to Theorem 4.3, we may and will assume that $`\stackrel{~}{Y}`$ is constructed together with the reflected Brownian snake $`\stackrel{~}{W}`$, in such a way that, for every $`t0`$, $$\stackrel{~}{Y}_t=_0^{\stackrel{~}{\tau }}𝑑\stackrel{~}{L}_s^t\delta _{\stackrel{~}{W}_s}$$ where $`(\stackrel{~}{L}_s^t,t0,s0)`$ denote the local times of the process $`\stackrel{~}{\beta }_s:=\zeta _{\stackrel{~}{W}_s}`$, which is (twice) a reflected Brownian motion stopped at time $`\stackrel{~}{\tau }=inf\{s0:\stackrel{~}{L}_s^0=a\}`$. The process $$X_t=_0^{\stackrel{~}{\tau }}𝑑\stackrel{~}{L}_s^t\delta _{\stackrel{~}{W}_s(t)}$$ is the weak limit of the processes $`X^\epsilon =\stackrel{~}{X}^\epsilon `$ and therefore must be a super-Brownian motion started at $`\mu `$. Let us recall the two key properties of the reflected snake $`\stackrel{~}{W}`$ (cf Proposition 4.2): $``$ Monotonicity property: Almost surely for every $`ss^{}`$ we have $`\stackrel{~}{W}_s(t)\stackrel{~}{W}_s^{}(t)`$ for every $`t[0,\stackrel{~}{\beta }_s\stackrel{~}{\beta }_s^{}]`$. $``$ Snake property: The set of discontinuities of the mapping $`s\stackrel{~}{W}_s`$ is contained in the zero set of $`\stackrel{~}{\beta }`$. Furthermore, if $`s<s^{}`$ belong to the same connected component of the complement of the zero set, we have $$\stackrel{~}{W}_s(t)=\stackrel{~}{W}_s^{}(t)\text{for every }t[0,\underset{r[s,s^{}]}{inf}\stackrel{~}{\beta }_r].$$ In order to state a useful preliminary result, we introduce some notation. Let us fix $`t>0`$, and denote by $`(a_i^t,b_i^t)`$, $`iI_t`$ the excursion intervals of $`\stackrel{~}{\beta }`$ above level $`t`$ (equivalently, these are the connected components of the open set $`\{s0:\stackrel{~}{\beta }_s>t\}`$). Note that the index set $`I_t`$ may be empty. For each $`iI_t`$, denote by $`e_i^t`$ the corresponding excursion $$e_i^t(s)=\stackrel{~}{\beta }_{(a_i^t+s)b_i^t}t,s0.$$ By the snake property of $`\stackrel{~}{W}`$, we have $$\stackrel{~}{W}_s(t)=\stackrel{~}{W}_{a_i^t}(t)=:z_i^t,s[a_i^t,b_i^t].$$ We denote by $`n(de)`$ the Itô measure of positive Brownian excursions. We normalize the measure $`n(de)`$ by declaring that the Poisson point process of excursions from $`0`$, i.e., the family of points $`(L_{a_i^0}^0,e_i^0)`$, has intensity $`dsn(de)`$. Proposition 5.1. Conditionally on $`X_t`$, the point measure $$\underset{iI_t}{}\delta _{(z_i^t,e_i^t)}$$ is Poisson with intensity $`X_t(dz)n(de)`$. Consequently, for every Borel subset $`A`$ of $`R`$, the process $$rZ_r^{t,A}=\stackrel{~}{Y}_{t+r}(dw)\mathbf{\hspace{0.17em}1}_A(w(t))$$ is a Feller diffusion started at $`X_t(A)`$. We recall that the Feller diffusion is a diffusion process $`Z`$ on $`R_+`$ whose transition kernels are characterized by the Laplace transform: $`E[\mathrm{exp}(\lambda Z_t)\|Z_0=z]=\mathrm{exp}(zu_t(\lambda ))`$ where $$u_t(\lambda )=\frac{\lambda }{1+\frac{1}{2}\lambda t}.$$ The total mass process $`X_t,1=\stackrel{~}{L}_{\stackrel{~}{\tau }}^t`$ is a Feller diffusion started at $`a`$. Proof. We denote by $`\tau _r^{(t)}`$ the right-continuous inverse of the function $`r\stackrel{~}{L}_r^t`$. Note that $`\tau _r^{(t)}<\mathrm{}`$ iff $`r<\stackrel{~}{L}_{\stackrel{~}{\tau }}^t=X_t,1`$. We can rewrite the definition of $`X_t`$ as $$X_t,\phi =_0^{\stackrel{~}{L}_{\stackrel{~}{\tau }}^t}𝑑r\phi (\stackrel{~}{W}_{\stackrel{~}{\tau }_r^{(t)}}(t)).$$ $`(5.1)`$ We also set for every $`r0`$, $$A_r^{(t)}=_0^r𝑑u\mathbf{\hspace{0.17em}1}_{\{\stackrel{~}{\beta }_u>t\}}$$ and we let $`\gamma _r^{(t)}`$ be the right-continuous inverse of the function $`rA_r^{(t)}`$. Finally we set $`\stackrel{~}{\beta }_r^{(t)}=\stackrel{~}{\beta }_{\gamma _r^{(t)}}t`$, for every $`r[0,A_{\stackrel{~}{\tau }}^{(t)})`$. We then claim that, conditionally on $`\{\stackrel{~}{L}_{\stackrel{~}{\tau }}^t=x\}`$, the process $`(\stackrel{~}{\beta }_r^{(t)},0r<A_{\stackrel{~}{\tau }}^{(t)})`$ is a reflected Brownian motion started at $`0`$ and killed at the first hitting time of $`x`$ by its local time at level $`0`$, and is independent of the process $`(\stackrel{~}{W}_{\stackrel{~}{\tau }_r^{(t)}}(t),0r<\stackrel{~}{L}_{\stackrel{~}{\tau }}^t)`$. Except for the independence statement, this is a familiar property of linear Brownian motion: See e.g. Section VI.2 of \[RY\]. To get the independence property, observe that the analogue of the process $`\stackrel{~}{\beta }^{(t)}`$ for the $`\epsilon `$-reflected system codes (in the sense of Section 2) the genealogy of the descendants of particles at time $`t`$. On the other hand, if $`\tau ^{\epsilon ,(t)}`$ denotes the right-continuous inverse of $`\stackrel{~}{L}^{\epsilon ,t}`$, the process $`(\stackrel{~}{W}_{\tau _r^{\epsilon ,(t)}}^\epsilon (t),r0)`$ just enumerates in increasing order the positions of the particles alive at $`t`$. The required independence is thus clear at the discrete level of the $`\epsilon `$-reflected system, and it is preserved under the passage to the limit (4.12). To complete the proof, write $`\mathrm{}_i^{(t)}`$ for the local time at $`0`$ of $`\stackrel{~}{\beta }^{(t)}`$ at the beginning, or the end, of excursion $`e_i^t`$. Note that $`\tau _{\mathrm{}_i^{(t)}}^{(t)}=b_i^t`$ and thus $$z_i^{(t)}=\stackrel{~}{W}_{\tau _{\mathrm{}_i^{(t)}}^{(t)}}(t).$$ $`(5.2)`$ The point measure $`\delta _{(\mathrm{}_i^{(t)},e_i^t)}`$ is the excursion process of the process $`\stackrel{~}{\beta }^{(t)}`$. Hence, conditionally on $`\{\stackrel{~}{L}_{\stackrel{~}{\tau }}^t=x\}`$, this point measure is Poisson with intensity $`1_{[0,x)}(\mathrm{})d\mathrm{}n(de)`$ and is independent of $`(\stackrel{~}{W}_{\stackrel{~}{\tau }_r^{(t)}}(t),0r<\stackrel{~}{L}_{\stackrel{~}{\tau }}^t)`$. The first part of the lemma then follows from this property, (5.2) and (5.1) (which just says that $`X_t`$ is the image of the measure $`1_{[0,\stackrel{~}{L}_{\stackrel{~}{\tau }}^{(t)})}(\mathrm{})d\mathrm{}`$ under the mapping $`\mathrm{}\stackrel{~}{W}_{\tau _{\mathrm{}}^{(t)}}(t)`$). To get the second assertion of the lemma, note that by the definition of $`\stackrel{~}{Y}_{t+r}`$, $$Z_r^{t,A}=\underset{iI_t}{}\mathrm{𝟏}_{\{z_i^{(t)}A\}}\mathrm{}^r(e_i^{(t)}),$$ where $`\mathrm{}^r(e_i^{(t)})`$ denotes the total local time of excursion $`e_i^{(t)}`$ at level $`r`$. By the first part of the proposition, conditionally on $`X_t`$, the random measure $$\underset{iI_t}{}\mathrm{𝟏}_{\{z_i^{(t)}A\}}\delta _{e_i^{(t)}}$$ is Poisson with intensity $`X_t(A)n(de)`$. Hence, conditionally on $`\{X_t(A)=x\}`$, the process $`(Z_r^{t,A},r0)`$ has the same law as $`(L_{\tau _x}^r,r0)`$, and the desired result follows from the celebrated Ray-Knight theorem on Brownian local time. Remark. We could easily sharpen the statement of Proposition 5.1 by conditioning on $`\stackrel{~}{Y}_t`$, or even on $`(\stackrel{~}{Y}_u,ut)`$ rather than on $`X_t`$. We will not need these refinements. 5.2 A priori estimates By \[KS\] or \[R\], we know that, almost surely for every $`t>0`$, the measure $`X_t`$ has a continuous density $`x_t(y)`$ with respect to Lebesgue measure on $`R`$, and the family $`(x_t(y),t>0,yR)`$ is jointly continuous. Some of our results will be proved under the following additional assumption: Assumption (H). The measure $`\mu `$ has a continuous density $`x_0(y)`$ with respect to Lebesgue measure. Under (H), the family $`(x_t(y),t0,yR)`$ is jointly continuous (see Theorem 8.3.2 in \[Da\]). In order to simplify the statements of the results in this subsection we introduce a constant $`\alpha `$. All the results hold for $`\alpha =0`$, assuming (H). Without this assumption, the results hold for any fixed strictly positive $`\alpha `$. For every $`t0`$, $`r>0`$ and $`zR`$, we set $$\psi _{t,t+r}(z)=sup\{\stackrel{~}{W}_s(t+r):\stackrel{~}{\beta }_st+r\text{ and }\stackrel{~}{W}_s(t)<z\},$$ with the usual convention $`sup\mathrm{}=\mathrm{}`$. We also consider the symmetric quantity: $$\widehat{\psi }_{t,t+r}(z)=inf\{\stackrel{~}{W}_s(t+r):\stackrel{~}{\beta }_st+r\text{ and }\stackrel{~}{W}_s(t)>z\},$$ Proposition 5.2. Let $`\eta (0,\frac{1}{2})`$ and $`c>0`$. Then, almost surely, one can choose $`\delta _0>0`$ small enough so that, for every $`\delta (0,\delta _0)`$, $`t\alpha `$ and $`zR`$, the condition $`x_t(z)c`$ implies $$\psi _{t,t+\delta }(z)z\delta ^{\frac{1}{2}\eta }.$$ Proof. For every $`t0`$ and $`zR`$ set $$\gamma ^{t,z}=inf\{s0:\stackrel{~}{\beta }_st\text{ and }\stackrel{~}{W}_s(t)z\},$$ with the convention $`inf\mathrm{}=\stackrel{~}{\tau }`$. Using the formula for $`X_t`$ in terms of $`\stackrel{~}{W}`$, and then the monotonicity property, we get $$X_t((\mathrm{},z])=_0^{\stackrel{~}{\tau }}𝑑\stackrel{~}{L}_s^t\mathbf{\hspace{0.17em}1}_{\{\stackrel{~}{W}_s(t)<z\}}=_0^{\stackrel{~}{\tau }}𝑑\stackrel{~}{L}_s^t\mathbf{\hspace{0.17em}1}_{\{s<\gamma ^{t,z}\}}=\stackrel{~}{L}_{\gamma ^{t,z}}^t.$$ On the other hand, if $`s<\gamma ^{t,z}`$ and $`\stackrel{~}{\beta }_st+\delta `$, we have $`\stackrel{~}{W}_s(t)<z`$ and $`\stackrel{~}{W}_s(t+\delta )\psi _{t,t+\delta }(z)`$. Therefore, $$X_{t+\delta }((\mathrm{},\psi _{t,t+\delta }(z)])=_0^{\stackrel{~}{\tau }}𝑑\stackrel{~}{L}_s^{t+\delta }\mathbf{\hspace{0.17em}1}_{\{\stackrel{~}{W}_s(t+\delta )\psi _{t,t+\delta }(z)\}}\stackrel{~}{L}_{\gamma ^{t,z}}^{t+\delta }.$$ Thanks to the Hölder continuity of Brownian local time in the time variable, we can choose $`\delta _1>0`$ so small that, for every $`\delta (0,\delta _1]`$, $`t0`$ and $`zR`$, $$\stackrel{~}{L}_{\gamma ^{t,z}}^{t+\delta }\stackrel{~}{L}_{\gamma ^{t,z}}^t\delta ^{\frac{1}{2}\eta }.$$ By combining all these facts we obtain for every $`\delta (0,\delta _1]`$, $`t0`$ and $`zR`$, $$X_{t+\delta }((\mathrm{},\psi _{t,t+\delta }(z)])X_t((\mathrm{},z])\delta ^{\frac{1}{2}\eta }.$$ $`(5.3)`$ Note that the set $`\{(t,y):x_t(y)>0\}`$ is contained in the graph of $`X`$ and is thus relatively compact. By uniform continuity, we can choose $`\delta _2>0`$ small enough so that, for every $`t\alpha `$ and $`zR`$, the condition $`x_t(z)c`$ implies that $`x_{t+\delta }(y)>\frac{c}{2}`$ for all $`\delta [0,\delta _2]`$ and $`y[z\delta _2,z+\delta _2]`$. In particular, if $`0<r<\delta _2`$ and $`\delta [0,\delta _2]`$, $$X_{t+\delta }((\mathrm{},zr])<X_{t+\delta }((\mathrm{},z])\frac{c}{2}r.$$ $`(5.4)`$ The proof of the following simple estimate for super-Brownian motion is postponed to the appendix. Lemma 5.3. Almost surely there exists $`\delta _3>0`$ such that, for every $`t\alpha `$, $`zR`$ and $`\delta (0,\delta _3)`$, $$\|X_{t+\delta }((\mathrm{},z])X_t((\mathrm{},z])\|\delta ^{\frac{1}{2}\eta }.$$ $`(5.5)`$ To complete the proof of Proposition 5.2, choose $`\delta _0(0,\delta _1\delta _2\delta _3)`$ and also such that $`\frac{4}{c}\delta _0^{\frac{1}{2}\eta }<\delta _2`$. Then, if $`t\alpha `$ and $`zR`$ are such that $`x_t(z)c`$, (5.3) and (5.5) give for $`\delta (0,\delta _0)`$, $$X_{t+\delta }((\mathrm{},\psi _{t,t+\delta }(z)])X_{t+\delta }((\mathrm{},z])2\delta ^{\frac{1}{2}\eta }.$$ Using (5.4) with $`r=\frac{4}{c}\delta ^{\frac{1}{2}\eta }`$, we get $$X_{t+\delta }((\mathrm{},\psi _{t,t+\delta }(z)])>X_{t+\delta }((\mathrm{},z\frac{4}{c}\delta ^{\frac{1}{2}\eta }]),$$ which implies $$\psi _{t,t+\delta }(z)z\frac{4}{c}\delta ^{\frac{1}{2}\eta }.$$ By replacing $`\eta `$ with $`\eta ^{}(0,\eta )`$ we can get rid of the factor $`\frac{4}{c}`$. We can immediately use Proposition 5.2 to derive some useful results on continuity properties of the paths $`\stackrel{~}{W}_s`$. Note that, if $`s(0,\stackrel{~}{\tau })`$ is such that $`\stackrel{~}{\beta }_st+r`$ and $`\stackrel{~}{W}_s(t)z`$, the monotonicity property of the reflected snake implies that $`\stackrel{~}{W}_s(t+r)\psi _{t,t+r}(z)`$. Using Proposition 5.2 and the symmetric result for $`\widehat{\psi }_{t,t+r}(z)`$, we get the following corollary. Recall that we take $`\alpha =0`$ if (H) is assumed to hold and $`\alpha >0`$ otherwise. Corollary 5.4. Let $`\eta (0,\frac{1}{2})`$ and $`c>0`$. Then almost surely we can choose $`\delta _0`$ small enough so that, for every $`t\alpha `$ and every $`s(0,\stackrel{~}{\tau })`$ such that $`\stackrel{~}{\beta }_s>t`$ and $`x_t(\stackrel{~}{W}_s(t))c`$, we have for every $`r[t,(t+\delta _0)\stackrel{~}{\beta }_s]`$, $$\|\stackrel{~}{W}_s(r)\stackrel{~}{W}_s(t)\|(rt)^{\frac{1}{2}\eta }.$$ 5.3 The key technical lemma Our aim is to refine the a priori estimates that were derived in the previous subsection. To this end, we will need a crucial technical lemma (Lemma 5.7 below), whose proof requires coming back to the approximating branching particle systems. Recall the notation $`(W^\epsilon ,\beta ^\epsilon ,Y^\epsilon )`$ of the previous sections. A much simplified version of the arguments of Section 4 yields the convergence in distribution $$(W^\epsilon ,Y^\epsilon )\underset{\epsilon 0}{\overset{(\mathrm{d})}{}}(W,Y),$$ where $`W`$ is a minor modification of the Brownian snake of \[L2\] (to be precise, $`W`$ is obtained by concatenating a Poisson point process of Brownian snake excursions with intensity $`\mu (dy)N_y`$, in the notation of \[L2\]) and $`Y`$ is the historical super-Brownian motion connected to $`W`$ via the formula $$Y_t=_0^\tau 𝑑L_s^t\delta _{W_s},$$ where $`(L_s^t,t0,s0)`$ are the local times of the lifetime process $`\beta _s=\zeta _{W_s}`$, which is a reflected Brownian motion stopped at time $`\tau =inf\{s0:L_s^0=a\}`$. (Our notation is slightly inconsistent with the previous sections, where $`\beta `$ was not stopped, but this should cause no confusion.) On the other hand (cf the proof of Theorem 4.3), we may and will assume that there is a sequence $`_0`$ of values of $`\epsilon `$ such that $$(\stackrel{~}{W}^\epsilon ,\stackrel{~}{Y}^\epsilon )\underset{\epsilon 0,\epsilon _0}{\overset{(\mathrm{d})}{}}(\stackrel{~}{W},\stackrel{~}{Y}).$$ By a compactness argument, and replacing the sequence $`_0`$ by a subsequence if necessary, we have also $$(W^\epsilon ,Y^\epsilon ,\stackrel{~}{W}^\epsilon ,\stackrel{~}{Y}^\epsilon )\underset{\epsilon 0,\epsilon _0}{\overset{(\mathrm{d})}{}}(W,Y,\stackrel{~}{W},\stackrel{~}{Y}).$$ By the Skorohod representation theorem, we can for every $`\epsilon _0`$ find a 4-tuple which has the same distribution as $`(W^\epsilon ,Y^\epsilon ,\stackrel{~}{W}^\epsilon ,\stackrel{~}{Y}^\epsilon )`$ (and for which we keep the same notation), in such a way that the previous convergence now holds a.s.: $$(W^\epsilon ,Y^\epsilon ,\stackrel{~}{W}^\epsilon ,\stackrel{~}{Y}^\epsilon )\underset{\epsilon 0,\epsilon _0}{\overset{(\mathrm{a}.\mathrm{s}.)}{}}(W,Y,\stackrel{~}{W},\stackrel{~}{Y}).$$ $`(5.6)`$ From now on we will restrict our attention to values of $`\epsilon `$ in the sequence $`_0`$ and assume that (5.6) holds. From the equality $`\stackrel{~}{X}^\epsilon =X^\epsilon `$, we also have $$_0^\tau 𝑑L_s^t\delta _{W_s(t)}=_0^{\stackrel{~}{\tau }}𝑑\stackrel{~}{L}_s^t\delta _{\stackrel{~}{W}_s(t)}=X_t,$$ and we see that $`\tau `$ coincides with $`\stackrel{~}{\tau }`$. We introduce the following more restrictive version of Assumption (H): Assumption (H’). The measure $`\mu `$ has a continuous density $`x_0(y)`$, which is Hölder continuous with exponent $`\frac{1}{2}\delta `$, for every $`\delta >0`$. As in the case of Assumption (H), in order to be able to use a single statement for a result with or without Assumption (H’), we take $`\alpha =0`$ if (H’) holds and otherwise we let $`\alpha `$ be a fixed strictly positive constant. We also fix a constant $`c(0,1)`$. Let $`\eta ,\eta ^{},\rho `$ be three positive constants, with $`0<\eta <\eta ^{}<1/4`$ and $`\rho (0,\frac{1}{2})`$. For every $`\delta (0,1)`$, we denote by $`E(\delta )`$ the event on which the following three conditions hold. A. For every $`s0`$, $`t[0,\beta _s]`$, and $`r[t,(t+\delta )\beta _s]`$, $$\|W_s(r)W_s(t)\|\frac{1}{2}(rt)^{\frac{1}{2}\eta }.$$ B. For every $`t\alpha `$ and $`s0`$ such that $`\stackrel{~}{\beta }_s>t`$ and $`x_t(\stackrel{~}{W}_s(t))c`$, we have for every $`r[t,(t+\delta )\stackrel{~}{\beta }_s]`$, $$\|\stackrel{~}{W}_s(r)\stackrel{~}{W}_s(t)\|(rt)^{\frac{1}{2}\eta }.$$ C. For every $`t\alpha `$, $`zR`$, and $`y[z\delta ^{\frac{1}{2}\eta ^{}},z+\delta ^{\frac{1}{2}\eta ^{}}]`$, $$\|x_t(z)x_t(y)\|\|zy\|^{\frac{1}{2}\rho }.$$ Note that the sets $`E(\delta )`$ are decreasing in $`\delta `$. We have $`P[_nE(2^n)]=1`$. The fact that properties A and B hold for $`\delta `$ small enough follows from the Hölder continuity properties of the Brownian snake paths (cf (2.7)) and Corollary 5.4 respectively. For property C, see Theorem 8.3.2 in \[Da\] when $`\alpha >0`$. When $`\alpha =0`$ (then (H’) is in force), the desired Hölder continuity of the densities is easily obtained from formula (8.3.5b) of \[Da\] by using the techniques of \[KS\]. Throughout this subsection, we fix $`\delta (0,1)`$, $`t\alpha `$ and $`zR`$. We plan to improve the estimates obtained on $`\psi _{t,t+\delta }(z)`$ in the previous subsection. We set $$\gamma =\delta ^{\frac{1}{2}\eta ^{}}$$ and we assume that $`\delta `$ has been chosen small enough so that $`\gamma >4\delta ^{\frac{1}{2}\eta }`$. Then, for every $`r[t,t+\delta ]`$, we set $$X_r^{}=Y_r(dw)\mathbf{\hspace{0.17em}1}_{\{w(t)(z\gamma ,z+\gamma )\}}\delta _{w(r)}.$$ The random measure $`X_r^{}`$ corresponds, for the historical super-Brownian motion $`Y`$, to the contribution of those particles alive at time $`r`$ whose ancestor at time $`t`$ lies in the interval $`(z\gamma ,z+\gamma )`$. Note that $`X_t^{}`$ is simply the restriction of $`X_t`$ to $`(z\gamma ,z+\gamma )`$. Our goal is to compare $`X_{t+\delta }^{}((\mathrm{},\psi _{t,t+\delta }(z)])`$ to $`X_t^{}((\mathrm{},z])`$ in the same way as we compared $`X_{t+\delta }((\mathrm{},\psi _{t,t+\delta }(z)])`$ to $`X_t((\mathrm{},z])`$ in (5.3) above. Unfortunately, the argument has to be significantly more complicated. We set for every $`\epsilon >0`$, $$\psi _{t,t+\delta }^\epsilon (z)=sup\{\stackrel{~}{W}_s^\epsilon (t+\delta ):\stackrel{~}{\beta }_s^\epsilon t+\delta \text{ and }\stackrel{~}{W}_s^\epsilon (t)<z\},$$ which represents for the $`\epsilon `$-reflected system the right-most position among those particles alive at time $`t+\delta `$ which are descendants of the particles located to the left of $`z`$ at time $`t`$. Lemma 5.5. We have $$\psi _{t,t+\delta }(z)=\underset{\epsilon 0}{lim}\psi _{t,t+\delta }^\epsilon (z)\text{a.s.}$$ Proof. This is basically a consequence of the convergence of $`\stackrel{~}{W}^\epsilon `$ towards $`\stackrel{~}{W}`$, which entails the convergence of $`\stackrel{~}{\beta }^\epsilon `$ to $`\stackrel{~}{\beta }`$. We also use the fact that in the definition of $`\psi _{t,t+\delta }(z)`$, i.e., $$\psi _{t,t+\delta }(z)=sup\{\stackrel{~}{W}_s(t+\delta ):\stackrel{~}{\beta }_st+\delta \text{ and }\stackrel{~}{W}_s(t)<z\},$$ we can replace the weak inequality $`\stackrel{~}{\beta }_st+\delta `$ by a strict one, and/or the strict inequality $`\stackrel{~}{W}_s(t)<z`$ by a weak one. To justify this, note that: (a) Almost surely, every $`s`$ such that $`\stackrel{~}{\beta }_s=t+\delta `$ is the limit of a sequence $`s_n`$ such that $`\stackrel{~}{\beta }_{s_n}>t+\delta `$ (simply because $`t+\delta `$ cannot be a local maximum of $`\stackrel{~}{\beta }`$). (b) With probability 1, there is no value of $`s`$ such that $`\stackrel{~}{W}_s(t)=z`$ and $`\stackrel{~}{\beta }_st+\delta `$ (this immediately follows from Proposition 5.1). We leave details to the reader. We now introduce a different approximation of $`\psi _{t,t+\delta }(z)`$. We consider in the (non-reflected) $`\epsilon `$-system those particles which are located at time $`t`$ in the interval $`(z\gamma ,z+\gamma )`$, and the descendants of these particles after time $`t`$. With this branching particle system (evolving over the time interval $`[t,\mathrm{})`$), we can associate a reflected system in the way explained in Subsection 3.1. We denote by $`\psi _{t,t+\delta }^{,\epsilon }(z)`$ the position in this new reflected system of the right-most particle at time $`t+\delta `$, among those particles which are descendants of the particles located to the left of $`z`$ at time $`t`$. For every $`r>0`$, we set $`\underset{¯}{x}(t,z,r)=inf\{x_t(y):\|yz\|r\}`$ and $`\overline{x}(t,z,r)=sup\{x_t(y):\|yz\|r\}`$. Lemma 5.6. We have $$P\left[\left(\underset{\epsilon 0}{lim\; sup}\{\psi _{t,t+\delta }^{,\epsilon }(z)\psi _{t,t+\delta }^\epsilon (z)\}\right)E(\delta )\{\underset{¯}{x}(t,z,\delta ^{1/2})c\}\right]2\mathrm{exp}(2c\delta ^{1/2}).$$ Proof. We introduce the following events: $$\mathrm{\Lambda }^+=\{s0:\stackrel{~}{\beta }_s>t+\delta \text{ and }\stackrel{~}{W}_s(t)(z\delta ^{1/2},z)\},$$ and $$\mathrm{\Lambda }^{}=\{s0:\stackrel{~}{\beta }_s>t+\delta \text{ and }\stackrel{~}{W}_s(t)(z,z+\delta ^{1/2})\}.$$ We first verify that a.s., $$\left(\left(\underset{\epsilon 0}{lim\; sup}\{\psi _{t,t+\delta }^{,\epsilon }(z)\psi _{t,t+\delta }^\epsilon (z)\}\right)E(\delta )\{\underset{¯}{x}(t,z,\delta ^{1/2})c\}\right)(\mathrm{\Lambda }^+\mathrm{\Lambda }^{})^c.$$ $`(5.7)`$ Suppose that $`\mathrm{\Lambda }^+\mathrm{\Lambda }^{}E(\delta )\{\underset{¯}{x}(t,z,\delta ^{1/2})c\}`$ holds. Then, there exists $`s_10`$ such that $`\stackrel{~}{\beta }_{s_1}>t+\delta `$ and $`W_{s_1}(t)(z\delta ^{1/2},z)`$. From property B in the definition of $`E(\delta )`$ we also have $`\|\stackrel{~}{W}_{s_1}(r)z\|<2\delta ^{\frac{1}{2}\eta }`$ for every $`r[t,t+\delta ]`$. Similarly, there exists $`s_20`$ such that $`\stackrel{~}{\beta }_{s_2}>t+\delta `$, $`W_{s_2}(t)(z,z+\delta ^{1/2})`$ and $`\|\stackrel{~}{W}_{s_2}(r)z\|<2\delta ^{\frac{1}{2}\eta }`$ for every $`r[t,t+\delta ]`$. By the convergence (5.6), the same properties hold for $`\epsilon >0`$ small enough, if we replace $`\stackrel{~}{W}_{s_i}`$ and $`\stackrel{~}{\beta }_{s_i}`$ by $`\stackrel{~}{W}_{s_i}^\epsilon `$ and $`\stackrel{~}{\beta }_{s_i}^\epsilon `$ respectively. On the other hand, by property A of the definition of $`E(\delta )`$ and the convergence (5.6), we have also for $`\epsilon `$ small enough, for every $`s`$ such that $`\beta _s^\epsilon t`$ and every $`r[t,(t+\delta )\stackrel{~}{\beta }_s^\epsilon ]`$, $$\|W_s^\epsilon (r)W_s^\epsilon (t)\|\delta ^{\frac{1}{2}\eta }.$$ In particular, if $`s`$ is such that $`\beta _s^\epsilon t`$ and $`\|W_s^\epsilon (t)z\|\gamma 4\delta ^{\frac{1}{2}\eta }`$, we have for every $`r[t,(t+\delta )\stackrel{~}{\beta }_s^\epsilon ]`$, $$\|W_s^\epsilon (r)z\|>2\delta ^{\frac{1}{2}\eta }.$$ We have shown that, on the event $`\mathrm{\Lambda }^+\mathrm{\Lambda }^{}E(\delta )\{\underset{¯}{x}(t,z,\delta ^{1/2})c\}`$, provided that $`\epsilon `$ is small enough: $``$There exist $`s_1`$ and $`s_2`$ such that $`\stackrel{~}{\beta }_{s_1}>t+\delta `$, $`\stackrel{~}{\beta }_{s_2}>t+\delta `$ and $$\begin{array}{cc}& W_{s_1}^\epsilon (t)(z\delta ^{1/2},z),\|\stackrel{~}{W}_{s_1}^\epsilon (r)z\|<2\delta ^{\frac{1}{2}\eta },r[t,t+\delta ]\hfill \\ & W_{s_2}^\epsilon (t)(z,z+\delta ^{1/2}),\|\stackrel{~}{W}_{s_2}^\epsilon (r)z\|<2\delta ^{\frac{1}{2}\eta },r[t,t+\delta ].\hfill \end{array}$$ $``$For every $`s0`$ such that $`\beta _s^\epsilon t`$ and $`\|W_s^\epsilon (t)z\|\gamma `$, $$\|W_s^\epsilon (r)z\|>2\delta ^{\frac{1}{2}\eta },r[t,(t+\delta )\stackrel{~}{\beta }_s^\epsilon ].$$ These properties allow us to apply Lemma 3.1. In the context of that lemma, the original system is the $`\epsilon `$-system considered after time $`t`$, the new (restricted) system consists of the descendants of the particles which are located at time $`t`$ in the interval $`(z\gamma ,z+\gamma )`$, and we take $`I=(z2\delta ^{\frac{1}{2}\eta },z+2\delta ^{\frac{1}{2}\eta })`$. Lemma 3.1 and the previous properties imply that the restrictions of the paths $`\stackrel{~}{W}_{s_1}^\epsilon `$ and $`\stackrel{~}{W}_{s_2}^\epsilon `$ to $`[t,t+\delta ]`$ still appear as restrictions of reflected historical paths in the new system. Note that in the definition of $`\psi _{t,t+\delta }^\epsilon (z)`$, respectively of $`\psi _{t,t+\delta }^{,\epsilon }(z)`$, we may restrict our attention to those reflected historical paths between times $`t`$ and $`t+\delta `$ in the original system, resp. in the new system, whose value at time $`t`$ lies in the interval $`[\stackrel{~}{W}_{s_1}^\epsilon (t),z)`$ (this is so because of the monotonicity property of reflected historical paths). Any such path is bounded below and above by $`\stackrel{~}{W}_{s_1}^\epsilon `$ and $`\stackrel{~}{W}_{s_2}^\epsilon `$ respectively, on the time interval $`[t,t+\delta ]`$. By Lemma 3.1 again, the class of paths that we need to consider is exactly the same for both the original system and the new one. This is enough to conclude that $`\psi _{t,t+\delta }^{,\epsilon }(z)=\psi _{t,t+\delta }^\epsilon (z)`$, and we get our claim (5.7). It follows from (5.7) that the probability considered in the lemma is bounded above by $$P[(\mathrm{\Lambda }^+\mathrm{\Lambda }^{})^c\{\underset{¯}{x}(t,z,\delta ^{1/2})c\}].$$ By the construction of $`\stackrel{~}{Y}`$, we have $$_0^{\stackrel{~}{\tau }}𝑑\stackrel{~}{L}_s^{t+\delta }\mathbf{\hspace{0.17em}1}_{\{\stackrel{~}{W}_s(t)(z\delta ^{1/2},z)\}}=\stackrel{~}{Y}_{t+\delta }(dw)\mathbf{\hspace{0.17em}1}_{\{w(t)(z\delta ^{1/2},z)\}}.$$ Hence the event $`\mathrm{\Lambda }^+`$ certainly holds if $$\stackrel{~}{Y}_{t+\delta }(dw)\mathbf{\hspace{0.17em}1}_{\{w(t)(z\delta ^{1/2},z)\}}>0.$$ It follows that $$P[(\mathrm{\Lambda }^+)^c\{\underset{¯}{x}(t,z,\delta ^{1/2})c\}]P[\{\stackrel{~}{Y}_{t+\delta }(dw)\mathbf{\hspace{0.17em}1}_{\{w(t)(z\delta ^{1/2},z)\}}=0\}\{\underset{¯}{x}(t,z,\delta ^{1/2})c\}],$$ and a similar bound holds if we replace $`\mathrm{\Lambda }^+`$ by $`\mathrm{\Lambda }^{}`$. By Proposition 5.1, the last quantity is bounded above by the probability that a Feller diffusion started at $`c\delta ^{1/2}`$ vanishes at time $`\delta `$. This probability is equal to $`\mathrm{exp}(2c\delta ^{1/2})`$, which completes the proof. We can now state the key lemma. We fix still another constant $`\eta ^{\prime \prime }(\eta ^{},1/4)`$. Lemma 5.7. There exist two positive constants $`C`$ and $`\kappa `$, that depend only on $`c,\eta ,\eta ^{},\eta ^{\prime \prime }`$ and $`\rho `$, such that $$\begin{array}{cc}& P[\{\|X_{t+\delta }^{}((\mathrm{},\psi _{t,t+\delta }(z)])X_t^{}((\mathrm{},z])\|>\delta ^{\frac{3}{4}\eta ^{\prime \prime }}\}\hfill \\ & E(\delta )\{\underset{¯}{x}(t,z,\delta ^{1/2})c\}\{\overline{x}(t,z,\gamma )c^1\}]C\mathrm{exp}(\delta ^\kappa ).\hfill \end{array}$$ Proof. For every $`r[t,t+\delta ]`$, set $$X_r^{,\epsilon }=Y_r^\epsilon (dw)\mathbf{\hspace{0.17em}1}_{\{w(t)(z\gamma ,z+\gamma )\}}\delta _{w(r)},$$ which represents the contribution at time $`r`$ of the descendants (in the non-reflected system) of particles which are located in $`(z\gamma ,z+\gamma )`$ at time $`t`$. From the convergence of $`Y^\epsilon `$ to $`Y`$, and the fact that $`Y_r(dw)\mathbf{\hspace{0.17em}1}_{\{w(t)=z\pm \gamma \}}=0`$, one can easily show that for every $`r[t,t+\delta ]`$, the measures $`X_r^{,\epsilon }`$ converge weakly to $`X_r^{}`$. In particular, a.s. for every $`yR`$, $$\underset{\epsilon 0}{lim}X_{t+\delta }^{,\epsilon }((\mathrm{},y])=X_{t+\delta }^{}((\mathrm{},y]).$$ From Lemma 5.5 and Lemma 5.6, we get that on the set $`E(\delta )\{\underset{¯}{x}(t,z,\delta ^{1/2}c\}`$, we have the convergence $$\underset{\epsilon 0}{lim}X_{t+\delta }^{,\epsilon }((\mathrm{},\psi _{t,t+\delta }^{,\epsilon }(z)])=X_{t+\delta }^{}((\mathrm{},\psi _{t,t+\delta }(z)]),$$ except possibly on a set of measure at most $`2\mathrm{exp}(2c\delta ^{1/2})`$. However, by the definition of $`\psi _{t,t+\delta }^{,\epsilon }(z)`$, and the monotonicity property of reflected systems, the quantity $`X_{t+\delta }^{,\epsilon }((\mathrm{},\psi _{t,t+\delta }^{,\epsilon }(z)])`$ is equal to $`\epsilon `$ times the number of descendants at time $`t+\delta `$ of the particles present at time $`t`$ in $`(z\gamma ,z)`$, for the $`\epsilon `$-reflected system constructed over the time interval $`[t,\mathrm{})`$ from the particles present at time $`t`$ in $`(z\gamma ,z+\gamma )`$. Since the law of the branching evolution is the same for the reflected system as for the original one, we see that conditionally on $`\{X_t^{,\epsilon }((\mathrm{},z])=\epsilon k\}`$, the variable $`X_{t+\delta }^{,\epsilon }((\mathrm{},\psi _{t,t+\delta }^{,\epsilon }(z)])`$ is distributed as $`\epsilon Z_\delta ^{(\epsilon ,k)}`$, where $`Z^{(\epsilon ,k)}`$ denotes a Galton-Watson process with critical binary branching at rate $`\epsilon ^1`$ and initial value $`k`$. Recall that $`X_t^{,\epsilon }((\mathrm{},z])`$ converges a.s. to $`X_t^{}((\mathrm{},z])`$. By standard limit theorems for Galton-Watson processes, $$(X_t^{,\epsilon }((\mathrm{},z]),X_{t+\delta }^{,\epsilon }((\mathrm{},\psi _{t,t+\delta }^{,\epsilon }(z)]))\underset{\epsilon 0}{\overset{(\mathrm{d})}{}}(X_t^{}((\mathrm{},z]),U),$$ where conditionally on $`X_t^{}((\mathrm{},z])=u`$, the variable $`U`$ is distributed as the value at time $`\delta `$ of a Feller diffusion started at $`u`$. Note that $`X_t^{}((\mathrm{},z])=X_t((z\gamma ,z])`$ and that on the set $`\{\overline{x}(t,z,\gamma )c^1\}`$ we have $`X_t^{}((\mathrm{},z])c^1\gamma =c^1\delta ^{\frac{1}{2}\eta ^{}}`$. Elementary estimates on the Feller diffusion, using only the form of the Laplace transform of the semigroup (see the appendix for very similar estimates) show that $$P\left[\{X_t^{}((\mathrm{},z])c^1\delta ^{\frac{1}{2}\eta ^{}}\}\{\|UX_t^{}((\mathrm{},z])\|\delta ^{\frac{3}{4}\eta ^{\prime \prime }}\}\right]C^{}\mathrm{exp}(\delta ^\kappa ^{}),$$ where the constants $`C^{}`$ and $`\kappa ^{}>0`$ depend only on $`c,\eta ^{}`$ and $`\eta ^{\prime \prime }`$. To complete the proof of the lemma, we write $$\begin{array}{cc}& P[\{\|X_{t+\delta }^{}((\mathrm{},\psi _{t,t+\delta }(z)])X_t^{}((\mathrm{},z])\|>\delta ^{\frac{3}{4}\eta ^{\prime \prime }}\}\hfill \\ & E(\delta )\{\underset{¯}{x}(t,z,\delta ^{1/2})c\}\{\overline{x}(t,z,\gamma )c^1\}]\hfill \\ & 2\mathrm{exp}(2c\delta ^{1/2})+P[\{\underset{\epsilon 0}{lim\; inf}\|X_{t+\delta }^{,\epsilon }((\mathrm{},\psi _{t,t+\delta }^{,\epsilon }(z)])X_t^{}((\mathrm{},z])\|>\delta ^{\frac{3}{4}\eta ^{\prime \prime }}\}\hfill \\ & \{X_t^{}((\mathrm{},z])c^1\delta ^{\frac{1}{2}\eta ^{}}\}]\hfill \\ & 2\mathrm{exp}(2c\delta ^{1/2})+\underset{\epsilon 0}{lim\; inf}P[\{\|X_{t+\delta }^{,\epsilon }((\mathrm{},\psi _{t,t+\delta }^{,\epsilon }(z)])X_t^{}((\mathrm{},z])\|>\delta ^{\frac{3}{4}\eta ^{\prime \prime }}\}\hfill \\ & \{X_t^{}((\mathrm{},z])c^1\delta ^{\frac{1}{2}\eta ^{}}\}]\hfill \\ & 2\mathrm{exp}(2c\delta ^{1/2})+P\left[\{X_t^{}((\mathrm{},z])c^1\delta ^{\frac{1}{2}\eta ^{}}\}\{\|UX_t^{}((\mathrm{},z])\|\delta ^{\frac{3}{4}\eta ^{\prime \prime }}\}\right]\hfill \\ & 2\mathrm{exp}(2c\delta ^{1/2})+C^{}\mathrm{exp}(\delta ^\kappa ^{}).\hfill \end{array}$$ 5.4 The main result We keep the notation introduced in the previous subsection. The reals $`t\alpha `$, $`\delta (0,1)`$ and $`zR`$ are fixed for the moment. Lemma 5.8. Assume that $`\eta ^{\prime \prime }>\frac{3}{2}\eta ^{}+\frac{1}{2}\rho `$. There exist two constants $`\overline{C}`$ and $`\overline{\kappa }>0`$, that depend only on $`c,\eta ^{},\eta ^{\prime \prime }`$ and $`\rho `$, such that $$P\left[\{X_{t+\delta }^{}((\mathrm{},z])X_t^{}((\mathrm{},z])+\delta ^{\frac{3}{4}\eta ^{\prime \prime }}\}E(\delta )\{\overline{x}(t,z,\gamma )c^1\}\right]\overline{C}\mathrm{exp}(\delta ^{\overline{\kappa }}).$$ The proof of this lemma is an application of standard techniques in the theory of super-Brownian motion. See the appendix for a detailed argument. Proposition 5.9. Under the assumptions of Lemma 5.8, there exist two constants $`C_0`$ and $`\kappa _0>0`$, that depend only on $`c,\eta ,\eta ^{},\eta ^{\prime \prime }`$ and $`\rho `$, such that $$\begin{array}{cc}& P[\{\psi _{t,t+\delta }(z)<z\frac{2}{c}\delta ^{\frac{3}{4}\eta ^{\prime \prime }}\}E(\delta )\hfill \\ & \{c\underset{¯}{x}(t,z,\gamma )\overline{x}(t,z,\gamma )c^1\}\{\underset{¯}{x}(t+\delta ,z,\gamma )>c\}]C_0\mathrm{exp}(\delta ^{\kappa _0}).\hfill \end{array}$$ Proof. Our argument is very similar to the proof of Proposition 5.2. We will assume that the event $`E(\delta )\{c\underset{¯}{x}(t,z,\gamma )\overline{x}(t,z,\gamma )c^1\}`$ holds. By Lemmas 5.7 and 5.8, we have on this set $$X_{t+\delta }^{}((\mathrm{},z])X_t^{}((\mathrm{},z])+\delta ^{\frac{3}{4}\eta ^{\prime \prime }}X_{t+\delta }^{}((\mathrm{},\psi _{t,t+\delta }(z)])+2\delta ^{\frac{3}{4}\eta ^{\prime \prime }}$$ $`(5.8)`$ except possibly on a set of probability at most $`C\mathrm{exp}(\delta ^\kappa )+\overline{C}\mathrm{exp}(\delta ^{\overline{\kappa }})`$. On the other hand, condition A in the definition of $`E(\delta )`$ (and the fact that $`\gamma >4\delta ^{\frac{1}{2}\eta }`$) ensures that the measures $`X_{t+\delta }^{}`$ and $`X_{t+\delta }`$ coincide over the interval $`(z\frac{\gamma }{2},z+\frac{\gamma }{2})`$. Hence, on the event $`\{\underset{¯}{x}(t+\delta ,z,\gamma )>c\}`$, we get $$X_{t+\delta }^{}((\mathrm{},z])2\delta ^{\frac{3}{4}\eta ^{\prime \prime }}>X_{t+\delta }^{}((\mathrm{},z\frac{2}{c}\delta ^{\frac{3}{4}\eta ^{\prime \prime }}]),$$ provided that $`\delta `$ is small enough so that $`\frac{2}{c}\delta ^{\frac{3}{4}\eta ^{\prime \prime }}\frac{\gamma }{2}`$. On the set where (5.8) holds, we get $$X_{t+\delta }^{}((\mathrm{},\psi _{t,t+\delta }(z)])>X_{t+\delta }^{}((\mathrm{},z\frac{2}{c}\delta ^{\frac{3}{4}\eta ^{\prime \prime }}]),$$ and the desired result follows. We now come to the main result of this section, which is a refinement of Corollary 5.4. Recall our conventions concerning $`\alpha `$—this constant is equal $`0`$ if (H’) is assumed to hold and otherwise $`\alpha `$ is a fixed strictly positive constant. Theorem 5.10. Let $`\lambda >0`$ and $`c(0,1)`$. Then a.s. we can choose $`\delta _0`$ small enough so that, for every $`t\alpha `$ and every $`s(0,\tau )`$ such that $`\stackrel{~}{\beta }_s>t`$ and $`x_t(\stackrel{~}{W}_s(t))c`$, we have for every $`r[t,(t+\delta _0)\stackrel{~}{\beta }_s]`$, $$\|\stackrel{~}{W}_s(r)\stackrel{~}{W}_s(t)\|(rt)^{\frac{3}{4}\lambda }.$$ Proof. We can choose $`\eta ,\eta ^{},\eta ^{\prime \prime }`$ with $`0<\eta <\eta ^{}<\eta ^{\prime \prime }<\lambda `$ and $`\rho (0,\frac{1}{2})`$ such that the assumptions of Lemma 5.8 hold. We then apply the estimate of Proposition 5.9 with $`\delta =2^n`$ ($`n`$ large enough) to all reals $`t[\alpha ,n]`$, $`z[n,n]`$ of the form $`t=k2^n`$, $`z=j2^n`$. We have already observed that $`P[_nE(2^n)]=1`$. Furthermore, if we assume that $`cx_t(z)c^1`$ we will have $`\underset{¯}{x}(t,z,2^{n(\frac{1}{2}\eta ^{})})c/2`$, $`\overline{x}(t,z,2^{n(\frac{1}{2}\eta ^{})})2/c`$, and $`\underset{¯}{x}(t+2^n,z,2^{n(\frac{1}{2}\eta ^{})})c/2`$, for all $`n`$ sufficiently large (depending on $`\omega `$ but not on $`t`$ and $`z`$). Then, by combining the estimate of Proposition 5.9 with the Borel-Cantelli lemma, we obtain the following property: There exists an integer $`n_0(\omega )`$ such that for every $`nn_0(\omega )`$, for every $`t=k2^n`$, $`z=j2^n`$ with $`t[\alpha ,n]`$, $`z[n,n]`$, the condition $`cx_t(z)c^1`$ implies $$\psi _{t,t+2^n}(z)z(2^n)^{\frac{3}{4}\lambda }.$$ Since the densities $`x_r(y)`$ are bounded over $`[\alpha ,\mathrm{})\times R`$, a simple argument shows that we can drop the condition $`x_t(z)c^1`$ in the previous assertion. Then, if $`s0`$ is such that $`\stackrel{~}{\beta }_st+2^n`$, where $`t`$ is of the form $`t=k2^n`$, we let $`z=j2^n`$ be such that $`z<\stackrel{~}{W}_s(t)z+2^n`$. If $`n`$ is large enough (again independently of the choice of $`s`$ and $`t`$), the condition $`x_t(\stackrel{~}{W}_s(t))2c`$ will imply $`x_t(z)>c`$. Then, by the definition of $`\psi _{t,t+\delta }(z)`$ and the preceding estimate, $$\stackrel{~}{W}_s(t+2^n)\psi _{t,t+2^n}(z)\stackrel{~}{W}_s(t)2^n(2^n)^{\frac{3}{4}\lambda }.$$ Thanks to this observation and a symmetry argument, we obtain that a.s. for $`n`$ large enough, for every $`t\alpha `$ of the form $`t=k2^n`$ and every $`s0`$ such that $`\stackrel{~}{\beta }_st+2^n`$ and $`x_t(\stackrel{~}{W}_s(t))2c`$, $$\|\stackrel{~}{W}_s(t+2^n)\stackrel{~}{W}_s(t)\|2(2^n)^{\frac{3}{4}\lambda }.$$ The statement of Theorem 5.10 now follows easily thanks to the usual chaining argument. Theorem 1.2 is an immediate consequence of Theorem 5.10. Note that, by the representation formula for $`\stackrel{~}{Y}`$ in terms of $`\stackrel{~}{W}`$, the set $`\mathrm{supp}\stackrel{~}{Y}_t`$ is contained in $`\{\stackrel{~}{W}_s;\stackrel{~}{\beta }_s=t\}`$, for every $`t>0`$, a.s. The comments following the statement of Theorem 1.2 are justified by Proposition 5.1. 6. Branching points In this last section, we prove Theorem 1.3. As in Section 5, we assume that the process $`\stackrel{~}{Y}`$ is constructed together with the reflected Brownian snake $`\stackrel{~}{W}`$, in such a way that we have the representation formula $$\stackrel{~}{Y}_t=_0^{\stackrel{~}{\tau }}𝑑\stackrel{~}{L}_s^t\delta _{\stackrel{~}{W}_s}.$$ We need a preliminary lemma. If $`s_1<s_2`$, we set $`m(s_1,s_2)=inf_{s[s_1,s_2]}\stackrel{~}{\beta }_s`$. Lemma 6.1. Almost surely, for any $`t>0`$ and any $`s_1<s_2`$ such that $`\stackrel{~}{\beta }_{s_1}=\stackrel{~}{\beta }_{s_2}=t`$ and $`0<m(s_1,s_2)<t`$, we have $$x_{m(s_1,s_2)}(\stackrel{~}{W}_{s_1}(m(s_1,s_2)))>0.$$ Proof. Let $`\alpha >0`$ and let $`A1`$ be an integer. Write $`E_A`$ for the event $`E_A=\{𝒢[0,A]\times [A,A]\}`$, where $`𝒢`$ is as above the graph of $`X`$. It is enough to prove that a.s. on $`E_A`$, the following holds: (P) For any $`t>\alpha `$ and $`s_1<s_2`$ such that $`\stackrel{~}{\beta }_{s_1}=\stackrel{~}{\beta }_{s_2}=t`$ and $`0<m(s_1,s_2)<t\alpha `$, we have $`x_{m(s_1,s_2)}(\stackrel{~}{W}_{s_1}(m(s_1,s_2)))>0`$. We first introduce some notation. Let $`e`$ be an excursion, that is a continuous function $`e:R_+R_+`$ such that $`e(s)>0`$ iff $`0<s<\sigma (e)`$, for some $`\sigma (e)>0`$. Set $$T_\alpha (e)=inf\{s0:e(s)=\alpha \}$$ and, if $`T_\alpha (e)<\mathrm{}`$, $$\begin{array}{cc}\hfill L_\alpha (e)& =sup\{t0:e(t)=\alpha \},\hfill \\ \hfill M_\alpha (e)& =\underset{T_\alpha (e)sL_\alpha (e)}{inf}e(s)\hfill \end{array}$$ By convention we take $`M_\alpha (e)=0`$ if $`T_\alpha (e)=\mathrm{}`$. Let $`r>0`$. Recall the notation $`I_r`$ and $`e_i^r,z_i^r`$, $`iI_r`$ introduced before Proposition 5.1, and for every $`c>0`$ and $`\delta (0,\alpha )`$, set $$N_r^\delta (\alpha ,c)=\underset{iI_r}{}\mathrm{𝟏}_{\{x_r(z_i^r)c\}}\mathbf{\hspace{0.17em}1}_{\{0<M_\alpha (e_i^r)\delta \}}.$$ Proposition 5.1 allows us to conclude that, $$\begin{array}{cc}\hfill E[N_r^\delta (\alpha ,c)\mathbf{\hspace{0.17em}1}_{E_A}]& E\left[\underset{iI_r}{}\mathrm{𝟏}_{\{\|z_i^r\|A\}}\mathbf{\hspace{0.17em}1}_{\{x_r(z_i^r)c\}}\mathbf{\hspace{0.17em}1}_{\{0<M_\alpha (e_i^r)\delta \}}\right]\hfill \\ & =E\left[_A^A𝑑zx_t(z)\mathbf{\hspace{0.17em}1}_{\{x_r(z)c\}}n(0<M_\alpha (e)\delta )\right]\hfill \\ & 2cA\delta \alpha ^2,\hfill \end{array}$$ using the easy formula $`n(0<M_\alpha (e)\delta )=\delta \alpha ^2`$. We apply this estimate with $`\delta =1/k`$ ($`k`$ large enough) and $`r=j/k`$ for all $`j=1,2,\mathrm{},Ak`$. It follows that $$E\left[\mathrm{𝟏}_{E_A}\underset{j=1}{\overset{\mathrm{}}{}}N_{j/k}^{1/k}(\alpha ,c)\right]2cA^2\alpha ^2.$$ In particular, if $`E_k(\alpha ,c,A)`$ denotes the event $`\{j1:N_{j/k}^{1/k}(\alpha ,c)1\}E_A`$, we have $$P\left[\underset{k\mathrm{}}{lim\; inf}E_k(\alpha ,c,A)\right]2cA^2\alpha ^2.$$ $`(6.1)`$ Suppose that property (P) fails. Then, we may find $`t>\alpha `$ and $`s_1<s_2`$ such that $`\stackrel{~}{\beta }_{s_1}=\stackrel{~}{\beta }_{s_2}=t`$ and $`0<m(s_1,s_2)<t\alpha `$, and furthermore $`x_{m(s_1,s_2)}(\stackrel{~}{W}_{s_1}(m(s_1,s_2)))=0`$. We take $`j`$ such that $`j/k<m(s_1,s_2)(j+1)/k`$, and observe that $`x_{j/k}(\stackrel{~}{W}_{s_1}(j/k))<c`$ for all $`k`$ sufficiently large, by the joint continuity of densities. Hence by considering the excursion of $`\stackrel{~}{\beta }`$ above level $`j/k`$ that contains $`s_1`$, we see that $`N_{j/k}^{1/k}(\alpha ,c)1`$ for all $`k`$ large. Therefore, if $`F(\alpha ,A)`$ denotes the event on which (P) fails, we have $$P[F(\alpha ,A)E_A]P[\underset{k\mathrm{}}{lim\; inf}E_k(\alpha ,c,A))]2cA^2\alpha ^2.$$ Since $`c`$ was arbitrary, we have $`P[F(\alpha ,A)E_A]=0`$, which completes the proof. Proof of Theorem 1.3. The representation formula for $`\stackrel{~}{Y}_t`$ implies that $$\mathrm{supp}\stackrel{~}{Y}_t=\{\stackrel{~}{W}_s:\stackrel{~}{\beta }_s=t\}.$$ (Note that the set on the right hand side is closed, by the continuity properties of $`\stackrel{~}{W}`$.) Hence if $`w_1`$ and $`w_2`$ belong to $`\mathrm{supp}\stackrel{~}{Y}_t`$ and $`w_1w_2`$, we can find $`s_1`$ and $`s_2`$ such that $`\stackrel{~}{\beta }_{s_1}=\stackrel{~}{\beta }_{s_2}=t`$, and $`\stackrel{~}{W}_{s_1}=w_1`$, $`\stackrel{~}{W}_{s_2}=w_2`$. With no loss of generality, we can assume $`s_1<s_2`$. We claim that $$m(s_1,s_2)=inf\{r[0,t]:w_1(r)w_2(r)\}.$$ $`(6.2)`$ The inequality $`m(s_1,s_2)inf\{r[0,t]:w_1(r)w_2(r)\}`$ is immediate from the snake property (when $`m(s_1,s_2)=0`$ there is nothing to prove). On the other hand, if we assume that there is a rational $`r(m(s_1,s_2),t)`$ such that $`\stackrel{~}{W}_{s_1}(r)=\stackrel{~}{W}_{s_2}(r)`$, then the monotonicity property implies $`\stackrel{~}{W}_s(r)=\stackrel{~}{W}_{s_1}(r)`$ for every $`s[s_1,s_2]`$ such that $`\stackrel{~}{\beta }_sr`$. Hence, $$X_r=_0^{\stackrel{~}{\tau }}𝑑\stackrel{~}{L}_s^r\delta _{\stackrel{~}{W}_s(r)}_{s_1}^{s_2}𝑑\stackrel{~}{L}_s^r\delta _{\stackrel{~}{W}_s(r)}=(\stackrel{~}{L}_{s_2}^r\stackrel{~}{L}_{s_1}^r)\delta _{W_{s_1}(r)},$$ which gives a contradiction since $`\stackrel{~}{L}_{s_2}^r\stackrel{~}{L}_{s_1}^r>0`$ by standard properties of linear Brownian motion. From now on, we assume $`m(s_1,s_2)>0`$. Note that we have also $`m(s_1,s_2)<t`$ since we assumed that $`w_1w_2`$. By Lemma 6.1, we have $`x_{m(s_1,s_2)}(\stackrel{~}{W}_{s_1}(m(s_1,s_2)))>0`$. By monotonicity (and the fact that the measure $`X_r`$ gives no mass to singletons), we get for every $`r(m(s_1,s_2),t)`$, $$_{s_1}^{s_2}𝑑\stackrel{~}{L}_s^r=X_r((\stackrel{~}{W}_{s_1}(r),\stackrel{~}{W}_{s_2}(r)))=_{\stackrel{~}{W}_{s_1}(r)}^{\stackrel{~}{W}_{s_2}(r)}𝑑zx_r(z),$$ and by the continuity of densities, it follows that $$\underset{rm(s_1,s_2)}{lim}\frac{\stackrel{~}{W}_{s_2}(r)\stackrel{~}{W}_{s_1}(r)}{\stackrel{~}{L}_{s_2}^r\stackrel{~}{L}_{s_1}^r}=x_{m(s_1,s_2)}(\stackrel{~}{W}_{s_1}(m(s_1,s_2)))>0.$$ $`(6.3)`$ Thanks to (6.3), the behavior of $`w_1(r)w_2(r)`$ as $`r\gamma _{w_1,w_2}=m(s_1,s_2)`$ is reduced to that of $`\stackrel{~}{L}_{s_2}^r\stackrel{~}{L}_{s_1}^r`$. Write $`s_0`$ for the (unique) time in $`(s_1,s_2)`$ such that $`\stackrel{~}{\beta }_{s_0}=m(s_1,s_2)`$. Standard results on Brownian path decompositions show that, for events that depend only on the asymptotic $`\sigma `$-field at time $`0`$, the processes $`\{\stackrel{~}{\beta }_{s_0u}\stackrel{~}{\beta }_{s_0},u[0,s_0s_1]\}`$ and $`\{\stackrel{~}{\beta }_{s_0+u}\stackrel{~}{\beta }_{s_0},u[0,s_2s_0]\}`$ behave as two independent 3-dimensional Bessel processes. It follows from this and the Ray-Knight theorem that the process $`\delta \stackrel{~}{L}_{s_2}^{m(s_1,s_2)+\delta }\stackrel{~}{L}_{s_1}^{m(s_1,s_2)+\delta }`$ has the same local path properties (for $`\delta `$ close to 0) as the sum of two independent squares of 2-dimensional Bessel processes, which is the square of a 4-dimensional Bessel process. If $`\delta R_\delta `$ is the square of a 4-dimensional Bessel process, the law of the iterated logarithm shows that $$\underset{\delta 0}{lim\; sup}\frac{R_\delta }{2\delta \mathrm{log}\|\mathrm{log}\delta \|}=1.$$ On the other hand, from the well-known rate of escape for Brownian motion in space (Theorem 6 in \[DE\] combined with time-inversion), we have for $`\alpha >0`$, $$\underset{\delta 0}{lim}\frac{R_\delta }{\delta \|\mathrm{log}\delta \|^{1\alpha }}=\mathrm{}.$$ We have just argued that the same properties hold if we replace $`R_\delta `$ with $`\stackrel{~}{L}_{s_2}^{m(s_1,s_2)+\delta }\stackrel{~}{L}_{s_1}^{m(s_1,s_2)+\delta }`$. This and (6.3) imply Theorem 1.3. Appendix Proof of Lemma 5.3. For a fixed value of $`z`$, the estimate of Lemma 5.3 follows from \[P2\]. As we need uniformity in $`z`$, we will provide a detailed argument. Recall the notation from Subsection 5.2, and especially the conventions concerning the constant $`\alpha `$. Recall that $`𝒢`$ denotes the graph of $`X`$ and for every integer $`A1`$ consider the event $$E_A=\{𝒢[0,A]\times [A,A];\underset{t\alpha ,yR}{sup}x_t(y)A\}.$$ Note that $`P[E_A]1`$ as $`A\mathrm{}`$. (We use assumption (H) when $`\alpha =0`$.) The key step of the proof is to show the following inequality for all $`t\alpha `$ and $`zR`$, $$P[\{\|X_{t+\delta }((\mathrm{},z])X_t((\mathrm{},z])\|\delta ^{\frac{1}{2}\eta }\}E_A]C\mathrm{exp}(\delta ^\kappa )$$ $`(A1),`$ where the constants $`C`$ and $`\kappa >0`$ may depend on $`A`$ but not on $`t,z`$ and $`\delta `$. To prove (A1), we may apply the Markov property at time $`t`$ and reduce the problem to the case $`t=0`$. More precisely it is enough to consider a super-Brownian motion $`\mathrm{\Gamma }=(\mathrm{\Gamma }_t,t0)`$ with initial value $`\mathrm{\Gamma }_0(dz)=g(z)dz`$, with a function $`g`$ bounded above by $`A`$ and such that $`g(z)𝑑z2A^2`$, and to prove that for every $`\delta (0,1)`$, $$P[\|\mathrm{\Gamma }_\delta ((\mathrm{},0])\mathrm{\Gamma }_0((\mathrm{},0])\|\delta ^{\frac{1}{2}\eta }]C\mathrm{exp}(\delta ^\kappa ).$$ $`(A2)`$ Let us first bound $`P[\mathrm{\Gamma }_\delta ((\mathrm{},0])\mathrm{\Gamma }_0((\mathrm{},0])\delta ^{\frac{1}{2}\eta }]`$. We know that for every $`\lambda >0`$, $$E[\mathrm{exp}(\lambda \mathrm{\Gamma }_\delta ((\mathrm{},0]))]=\mathrm{exp}(\mathrm{\Gamma }_0,u_\delta ),$$ where $`u_t(z)`$ solves the integral equation $$u_t(z)+\frac{1}{2}E_z\left[_0^tu_{tr}(B_r)^2𝑑r\right]=\lambda P_z[B_t0],$$ if $`B`$ is a linear Brownian motion started at $`z`$ under $`P_z`$. The integral equation gives the bound $$u_t(z)\lambda P_z[B_t0]\frac{\lambda ^2}{2}t.$$ We use this bound in the following estimates, $$\begin{array}{cc}& P[\mathrm{\Gamma }_\delta ((\mathrm{},0])\mathrm{\Gamma }_0((\mathrm{},0])\delta ^{\frac{1}{2}\eta }]\hfill \\ & \mathrm{exp}(\lambda \delta ^{\frac{1}{2}\eta }+\lambda \mathrm{\Gamma }_0((\mathrm{},0]))E[\mathrm{exp}(\lambda \mathrm{\Gamma }_\delta ((\mathrm{},0]))]\hfill \\ & \mathrm{exp}(\lambda \delta ^{\frac{1}{2}\eta }+\frac{\lambda ^2}{2}\delta \mathrm{\Gamma }_0,1)\mathrm{exp}\left(\lambda \left(\mathrm{\Gamma }_0((\mathrm{},0])𝑑zg(z)P_z[B_\delta 0]\right)\right).\hfill \end{array}$$ Note that for every $`\epsilon >0`$, $$\left\|_{\mathrm{}}^0𝑑zg(z)𝑑zg(z)P_z[B_\delta 0]\right\|C_\epsilon \delta ^{\frac{1}{2}\epsilon },$$ with a constant $`C_\epsilon `$ depending only on $`\epsilon `$ and $`A`$. By choosing $`\lambda =\gamma ^{\frac{1}{2}+\epsilon }`$ with $`0<\epsilon <\eta `$, we arrive at the desired estimate for $`P[\mathrm{\Gamma }_\delta ((\mathrm{},0])\mathrm{\Gamma }_0((\mathrm{},0])\delta ^{\frac{1}{2}\eta }]`$. Slightly different arguments apply to $`P[\mathrm{\Gamma }_\delta ((\mathrm{},0])\mathrm{\Gamma }_0((\mathrm{},0])+\delta ^{\frac{1}{2}\eta }]`$. In fact, it is easier to observe that $$\begin{array}{cc}& P[\mathrm{\Gamma }_\delta ((\mathrm{},0])\mathrm{\Gamma }_0((\mathrm{},0])+\delta ^{\frac{1}{2}\eta }]\hfill \\ & P[\mathrm{\Gamma }_\delta ,1\mathrm{\Gamma }_0,1+\frac{1}{2}\delta ^{\frac{1}{2}\eta }]+P[\mathrm{\Gamma }_\delta ((0,\mathrm{}))\mathrm{\Gamma }_0((0,\mathrm{}))\frac{1}{2}\delta ^{\frac{1}{2}\eta }].\hfill \end{array}$$ $`(A3)`$ We have just shown how to bound the second term on the right hand side of (A3). As for the first term, we need simply recall that $`\mathrm{\Gamma }_t,1`$ is a Feller diffusion and use the fact that for $`\lambda (0,\frac{2}{\delta })`$ $$E[\mathrm{exp}(\lambda \mathrm{\Gamma }_\delta ,1)]=\mathrm{exp}\left(\frac{\lambda \mathrm{\Gamma }_0,1}{1\frac{1}{2}\lambda \delta }\right).$$ $`(A4)`$ This immediately leads to the estimate needed to complete the proof of (A2) and (A1). From (A1) and the Borel-Cantelli lemma, we get that a.s. there is an integer $`n_0(\omega )`$ such that, for every $`nn_0`$, for every $`t0`$ of the form $`t=j2^n`$ and every $`zR`$ of the form $`z=k2^n`$, we have $$\|X_{t+2^n}((\mathrm{},z])X_t((\mathrm{},z])\|(2^n)^{\frac{1}{2}\eta }.$$ Note that for every fixed $`z`$, the process $`tX_t((\mathrm{},z])`$ has continuous sample paths a.s. (see e.g. Corollary 6 in \[P2\]). The proof of Lemma 5.3 is easily completed thanks to this observation, the preceding bound and the usual chaining argument. Proof of Lemma 5.8. This is very similar to the proof of (A1) above. Note that the process $`(X_{t+r}^{},0r\delta )`$ is a super-Brownian motion started at $`X_t^{}`$, which is simply the restriction of $`X_t`$ to $`[z\gamma ,z+\gamma ]`$. Thanks to this observation and the definition of $`E(\delta )`$, we see that it is enough to prove the following statement. Let $`\mathrm{\Gamma }=(\mathrm{\Gamma }_r,r0)`$ be super-Brownian motion with initial value $`\mathrm{\Gamma }_0(dz)=g(z)dz`$. Assume that the function $`g`$ vanishes outside $`[\gamma ,\gamma ]`$ and that $`cg(z)c^1`$ and $`\|g(z)g(z^{})\|\|zz^{}\|^{\frac{1}{2}\rho }`$ for all $`z,z^{}[\gamma ,\gamma ]`$. Then, $$P[\mathrm{\Gamma }_\delta ((\mathrm{},0])\mathrm{\Gamma }_0((\mathrm{},0])+\delta ^{\frac{3}{4}\eta ^{\prime \prime }}]\overline{C}\mathrm{exp}(\delta ^{\overline{\kappa }}),$$ $`(A5)`$ where the constants $`\overline{C}`$ and $`\overline{\kappa }`$ depend only on $`c,\eta ^{},\eta ^{\prime \prime }`$ and $`\rho `$. In a way similar to (A3) we first write $$\begin{array}{cc}& P[\mathrm{\Gamma }_\delta ((\mathrm{},0])\mathrm{\Gamma }_0((\mathrm{},0])+\delta ^{\frac{3}{4}\eta ^{\prime \prime }}]\hfill \\ & P[\mathrm{\Gamma }_\delta ,1\mathrm{\Gamma }_0,1+\frac{1}{2}\delta ^{\frac{3}{4}\eta ^{\prime \prime }}]+P[\mathrm{\Gamma }_\delta ((0,\mathrm{}))\mathrm{\Gamma }_0((0,\mathrm{}))\frac{1}{2}\delta ^{\frac{3}{4}\eta ^{\prime \prime }}].\hfill \end{array}$$ Thanks to (A4), we see that, for $`\lambda <2/\delta `$, $$\begin{array}{cc}\hfill P[\mathrm{\Gamma }_\delta ,1\mathrm{\Gamma }_0,1+\frac{1}{2}\delta ^{\frac{3}{4}\eta ^{\prime \prime }}]& \mathrm{exp}(\lambda (\mathrm{\Gamma }_0,1+\frac{1}{2}\delta ^{\frac{3}{4}\eta ^{\prime \prime }}))E[e^{\lambda \mathrm{\Gamma }_\delta ,1}]\hfill \\ & =\mathrm{exp}(\frac{\lambda }{2}\delta ^{\frac{3}{4}\eta ^{\prime \prime }})\mathrm{exp}\left(\frac{\lambda ^2\mathrm{\Gamma }_0,1\delta /2}{1\frac{1}{2}\lambda \delta }\right).\hfill \end{array}$$ Since $`\mathrm{\Gamma }_0,12c^1\gamma =2c^1\delta ^{\frac{1}{2}\eta ^{}}`$, we get a bound of the desired form by taking $`\lambda =\delta ^{\frac{3}{4}+\epsilon }`$ with $`\eta ^{\prime \prime }>\epsilon >\eta ^{}`$. For the other term, we proceed as in the proof of Lemma 5.3: $$P[\mathrm{\Gamma }_\delta ((0,\mathrm{}))\mathrm{\Gamma }_0((0,\mathrm{}))\frac{1}{2}\delta ^{\frac{3}{4}\eta ^{\prime \prime }}]\mathrm{exp}(\lambda (\mathrm{\Gamma }_0((0,\mathrm{}))\frac{1}{2}\delta ^{\frac{3}{4}\eta ^{\prime \prime }}))E[e^{\lambda \mathrm{\Gamma }_\delta ((0,\mathrm{}))}],$$ and $`E[e^{\lambda \mathrm{\Gamma }_\delta ((0,\mathrm{}))}]=\mathrm{exp}(\mathrm{\Gamma }_0,u_\delta )`$, with $`u_\delta (y)\lambda P_y[B_\delta >0]\frac{1}{2}\lambda ^2\delta `$. It follows that $$\begin{array}{cc}& P[\mathrm{\Gamma }_\delta ((0,\mathrm{}))\mathrm{\Gamma }_0((0,\mathrm{}))\frac{1}{2}\delta ^{\frac{3}{4}\eta ^{\prime \prime }}]\hfill \\ & \mathrm{exp}(\frac{1}{2}\lambda \delta ^{\frac{3}{4}\eta ^{\prime \prime }}+\frac{\lambda ^2}{2}\delta \mathrm{\Gamma }_0,1)\mathrm{exp}\left(\lambda \left(_0^{\mathrm{}}𝑑zg(z)𝑑zg(z)P_z[B_\delta >0]\right)\right)\hfill \\ & \mathrm{exp}(\frac{1}{2}\lambda \delta ^{\frac{3}{4}\eta ^{\prime \prime }}+c^1\lambda ^2\delta \gamma )\mathrm{exp}(4\lambda \gamma ^{\frac{3}{2}\rho }),\hfill \end{array}$$ where in the last line we used our assumption that $`\|g(z)g(0)\|\|z\|^{\frac{1}{2}\rho }`$ to bound $`_0^{\mathrm{}}𝑑zg(z)𝑑zg(z)P_z[B_\delta >0]`$. In view of the assumptions of Lemma 5.8, we can now choose $`\lambda =\delta ^{\frac{3}{4}\epsilon }`$, with $`\eta ^{\prime \prime }>\epsilon >\frac{3}{2}\eta ^{}+\frac{\rho }{2}`$, and we arrive at a bound of the desired form. This completes the proof. 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Mech. 18, 973–989 (1968/69) Krzysztof Burdzy Department of Mathematics University of Washington Box 354350 Seattle, WA 98195-4350, USA e-mail: burdzy@math.washington.edu Jean-François Le Gall DMA — Ecole Normale Supérieure 45, rue d’Ulm 75230 Paris Cedex 05, France e-mail: legall@dma.ens.fr
warning/0003/astro-ph0003457.html	ar5iv	text	# Mid-infrared ISO spectroscopy of NGC 4945 Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA ## 1 Introduction NGC 4945 is a nearby, large (20$`\mathrm{}\times `$4$`\mathrm{}`$) spiral galaxy seen nearly edge on (i $``$ 78$`\mathrm{°}`$; Ott Ott (1995)). At a recession velocity of 560km/s it is at the mean radial velocity of the Centaurus group (Hesser et al. Hesser (1984)), of which it is believed to be a member. Distance estimates vary between 3.5 and 4.0Mpc (see Bergman et al. Bergman (1992) and Mauersberger et al. Mauersberger (1996) for discussions). In this paper we will adopt a distance of 3.9 Mpc (Bergman et al. Bergman (1992)), which implies that 1$`\mathrm{}`$ is equivalent to 18pc. NGC 4945 is one of the brightest infrared galaxies in the sky: S=24Jy, S=43Jy, S=588Jy, S=1416Jy (Rice et al. Rice (1988)). The total infrared luminosity amounts to L(8-1000$`\mu `$m)=2.95$`\times `$10<sup>10</sup> L, $``$75% of which originates in the central 12$`\mathrm{}\times `$9$`\mathrm{}`$ (Brock et al. Brock (1988)). Near infrared observations reveal the nuclear region to be the site of a powerful, yet visually obscured, starburst. Br$`\gamma `$ (Moorwood et al. 1996a ) and Pa$`\alpha `$ (Marconi et al. Marconi (2000)) recombination line maps show the starburst to be concentrated in a circumnuclear disk or ring $``$200pc across (11$`\mathrm{}`$). Further evidence for (a period of) strong star formation comes from the discovery of a conical structure, roughly perpendicular to the galaxy major axis. It is believed to be a cavity, vacated by a starburst-driven superwind (Heckman et al. Heckman (1990); Moorwood et al. 1996a ). The non-detection of $`[`$O iii$`]`$ within the cone and the absence of coronal lines excludes an AGN as the driver of the outflow. Clear evidence for the presence of an AGN comes from hard X-ray observations (Iwasawa et al. Iwasawa (1993); Guainazzi et al. Guainazzi (2000)). The AGN X-ray emission is however heavily absorbed by a column density of 10<sup>24.7</sup>cm<sup>-2</sup>, which obscures the AGN at all optical and infrared wavelengths. Previous authors have attributed most of the IR luminosity to the starburst (e.g. Moorwood & Oliva Moorwood94 (1994); Koornneef & Israel Koornneef96 (1996)). Hard X-ray observations with BeppoSAX indicate that the bolometric luminosity may as well be accounted for by the AGN alone (Guainazzi et al. Guainazzi (2000)). 3cm&6cm ATCA radio maps of the central region of NGC 4945 (Forbes & Norris Forbes (1998)) are dominated by strong nuclear emission, and emission extended along the disc of the galaxy. There is also evidence for some filamentry structure associated with the cavity cleared by the starburst superwind. VLBI observations by Sadler et al. (Sadler (1995)) reveal the existence of a compact radio core. This, as well as the presence of H<sub>2</sub>O megamasers in a Keplerian disc about a $``$10<sup>6</sup> M black hole (Greenhill et al. Greenhill (1997)), are taken as further evidence for the presence of an AGN. Near infrared observations of molecular hydrogen emission in NGC 4945 have been reported by several authors over the last 15 years (e.g. Moorwood & Glass Moorwood84 (1984); Moorwood & Oliva Moorwood88 (1988); Koornneef Koornneef93 (1993); Moorwood & Oliva Moorwood94 (1994); Koornneef & Israel Koornneef96 (1996); Moorwood et al. 1996a ; Quillen et al. Quillen (1999); Marconi et al. Marconi (2000)). While fluxes are known for eight ro-vibrational transitions accesible from the ground (Koornneef & Israel Koornneef96 (1996)), spatial information is available only for the (1-0) S(1) 2.1218$`\mu `$m line. These observations show the H<sub>2</sub> emission to be associated with the hollow cone, not with the starburst traced in hydrogen recombination emission. The absence of a correlation argues against photons as the source of excitation. Instead, the emission is attributed to shock heating of the molecular material at the face of the cavity (Moorwood et al. 1996a ; Marconi et al. Marconi (2000)). Mid-infrared spectroscopy is much less affected by intervening extinction than the UV and optical equivalents, with A($`\lambda `$)/A<sub>V</sub> less than 0.1. Observations of the central region of the galaxy, using the mid-infrared spectrometer SWS (De Graauw et al. deGraauw (1996)) and the spectrophotometer PHT-S (Lemke et al. Lemke (1996)), both aboard ISO (Kessler et al. Kessler (1996)), are therefore very useful to study the nuclear components otherwise hidden by heavy extinction. In Sect. 3.1 we present the results of the search for high excitation emission from the AGN. In Sect. 3.2 we study the properties of the nuclear starburst. In Sect. 3.3 we discuss the dominant nuclear power source. Sect. 3.4 discusses the broad emission and absorption features, tracing the properties of the interstellar medium in and in front of the nucleus. Finally, in Sect. 3.5 we discuss the physical conditions and excitation of the warm molecular hydrogen. ## 2 Observations As part of the Central Program “MPEXGAL”, we have observed the central region of NGC 4945 with the Short Wavelength Spectrometer (SWS) and the spectrophotometer PHT-S on board ISO. ### 2.1 SWS spectroscopy SWS grating line profile scans (SWS02 mode) were obtained on 1996 February 6 for 28 spectral lines in the range of 2.42 to 40.34$`\mu `$m. The spectral resolution in this range varies between R=900 and 2000, corresponding to a velocity resolution of 330–150 km/s. Aperture sizes used range between 14$`\mathrm{}\times `$20$`\mathrm{}`$ and 20$`\mathrm{}\times `$33$`\mathrm{}`$. SWS was centered on the 1.4GHz continuum peak from Ables et al. (Ables (1987)), which coincides with the the position of the L-band peak of Moorwood et al. (1996a ). At the time of observation the position angle of the major axis of the SWS apertures was 35.4$`\mathrm{°}`$, 10$`\mathrm{°}`$ off from that of the galaxy major axis (45$`\mathrm{°}`$). The data reduction was performed using the SWS Interactive Analysis system (SIA; Lahuis et al. Lahuis (1998); Wieprecht et al. Wieprecht (1998)), putting special emphasis on tools to improve cosmic ray spike removal, dark current subtraction and flat fielding. The wavelength calibration of SWS is discussed by Valentijn et al. (1996a ). We used calibration files as of March 1999. The accuracy of the flux calibration is estimated to be 30% (Schaeidt et al. Schaeidt (1996)). The resulting spectra are shown in Fig 1. In total 17 spectral lines were detected. For another 11 lines we derived upper limits, using gaussian profiles of width equal to other lines of the same (or comparable) species, scaled to a peak height corresponding to approximately 3$`\sigma `$ of the noise. Both detections and upper limits are presented in Table 1. ### 2.2 PHT-S spectrophotometry We have obtained two low resolution ($`\lambda `$/$`\mathrm{\Delta }\lambda `$90) ISO-PHT-S spectra, on 1996 October 12 and 1997 August 3, respectively. ISO-PHT-S comprises two low-resolution grating spectrometers covering simultaneously the wavelength range 2.47 to 4.87$`\mu `$m and 5.84 to 11.62$`\mu `$m. The spectrum is registered by two linear arrays of 64 Si:Ga detectors with a common entrance aperture of 24$`\mathrm{}`$ $`\times `$ 24$`\mathrm{}`$. The measurements were carried out in rectangular chopped mode, using a chopper throw of 180$`\mathrm{}`$. The resulting spectra thus are free of contributions from zodiacal light, that would otherwise affect the spectrum. The pure on-source integration times were 512 and 1024 s. The ISO-PHT-S data were reduced using PIA<sup>1</sup><sup>1</sup>1PIA is a joint development by the ESA Astrophysics Division and the ISO-PHT Consortium (Gabriel et al. Gabriel (1997)) version 8.1. Steps in the data reduction included: 1) deglitching on ramp level. 2) subdivision of ramps in two sections of 32 non destructive read-outs. 3) ramp fitting to derive signals. 4) masking of bad signals by eye-inspection. 5) kappa sigma and min/max clipping on remaining signal distribution. 6) determination of average signal per chopper plateau. 7) masking or correction of bad plateaux by eye-inspection. 8) background subtraction using all but the first four plateaux. 9) finally, flux calibration, using the signal dependent spectral response function of Acosta-Pulido (Acosta99 (1999)). This spectral response function avoids some deficiencies of the previous PIA response function, in particular in the 3$`\mu `$m region near the “Ice” feature. The absolute calibration is accurate to within 20%. The two resulting spectra were obtained at slightly different position angles about the nucleus. For the first, the square aperture was aligned with the galaxy major axis (45$`\mathrm{°}`$). For the second, the position angle was 31.1$`\mathrm{°}`$. Fig. 2 shows the averaged ISO-PHT-S spectrum. The on-source integration times were used as weight factors in the computation of the average spectrum. A number of emission lines can be identified in the ISO-PHT-S spectrum. These include 9.66$`\mu `$m H<sub>2</sub> (0-0) S(3), the unresolved blend of 6.99$`\mu `$m $`[`$Ar ii$`]`$ and 6.91$`\mu `$m H<sub>2</sub> (0-0) S(5), and 4.05$`\mu `$m H Br$`\alpha `$. For the 9.66$`\mu `$m H<sub>2</sub> (0-0) S(3) line, not included in the SWS02 line scans, we measure a flux of 5.4$`\times `$10<sup>-20</sup> W/cm<sup>2</sup>, with an uncertainty of 30%. ## 3 Results ### 3.1 AGN not seen at mid-infrared wavelengths NGC 4945 is a peculiar and interesting target for studying the relation of AGN and star formation in galaxies. Clear evidence for hidden AGN activity comes from hard X-ray observations. NGC 4945 is amongst the brighest hard X-ray emitting galaxies and exhibits variability of its 13–200keV flux on timescales of $``$10hrs, which clearly establishes its AGN origin (Iwasawa et al. Iwasawa (1993); Guainazzi et al. Guainazzi (2000); Marconi et al. Marconi (2000)). The AGN X-ray emission is heavily absorbed by a column density of 10<sup>24.7</sup> cm<sup>-2</sup> (corresponding to A<sub>V</sub>$``$2600), a high value, but within the range observed for Seyfert 2 galaxies (e.g. Risaliti et al. Risaliti (1999)). In unified schemes, the X-ray obscuration measures a line of sight towards the very center. Obscuration towards the NLR probes a different line of sight and is usually significantly lower, making the NLR visible in Seyfert 2 galaxies. Mid-infrared high excitation lines are able to penetrate a far larger dust obscuration than their optical and UV counterparts. They are therefore ideally suited as tracers of embedded AGN activity. Mid-infrared emission lines like $`[Ne\text{v}]`$ 14.32$`\mu `$m & 24.32$`\mu `$m, $`[Ne\text{vi}]`$ 7.65$`\mu `$m and $`[O\text{iv}]`$ 25.9$`\mu `$m are prominently present in the spectrum of all Seyferts observed with ISO (Moorwood et al. 1996b ; Genzel et al. Genzel (1998)). On the other hand, the same emission lines are also weakly visible towards, for instance, supernova remnant RCW 103 (Oliva et al. Oliva (1999)). $`[O\text{iv}]`$ emission (at the few percent level compared to $`[Ne\text{ii}]`$) has also been detected in a sample of starburst galaxies (Lutz et al. Lutz98 (1998)), again at a level much weaker than seen in typical AGNs. The origin of the weak level emission in these sources is believed to be shocks. A detection of any of the high excitation lines discussed above does therefore not automatically imply the detection of an AGN in NGC 4945. We do not detect the lines of $`[Ne\text{v}]`$ and $`[Ne\text{vi}]`$. No trace of $`[Ne\text{vi}]`$ is seen in the wing of the nearby PAH emission feature (Fig. 1). From Fig. 1 it might appear that the two $`[Ne\text{v}]`$ lines were indeed detected. However, at the level at which the features appear, instrumental effects play a significant role. In the 14.35$`\mu `$m line scan a strong fringe in the relative spectral response function coincides exactly with the expected position for the $`[Ne\text{v}]`$ line. Depending on the size of the emitting area, the feature may be entirely attributed to this instrumental effect. We therefore chose to state an upper limit for the 14.32$`\mu `$m $`[Ne\text{v}]`$ line. The feature seen in the other $`[Ne\text{v}]`$ scan, centered at 24.4$`\mu `$m, was registered by only two detectors, although 12 detectors scanned over the central wavelength. As visible in Fig. 1, the feature is redshifted with respect to the NGC 4945 systemic velocity. This shift is not observed for any other line we observed. We therefore derive an upper limit for this $`[Ne\text{v}]`$ line too. The only detected high ionization line in NGC 4945 is the 25.9$`\mu `$m $`[O\text{iv}]`$ line. An AGN contribution to this line is possible — to match the limits on higher excitation lines, only part of the $`[O\text{iv}]`$ emission would be related to an AGN. The detection of possibly shock-related $`[O\text{iv}]`$ in many starbursts (Lutz et al. Lutz98 (1998)) cautions, however, that this may be a more plausible origin of $`[O\text{iv}]`$ in NGC 4945. The ratio of 0.033 with respect to $`[Ne\text{ii}]`$+0.44$`\times [Ne\text{iii}]`$ is above average, but well within the range observed for the Lutz et al. (Lutz98 (1998)) starbursts, also considering that the high extinction (Sect. 3.2) will increase the observed ratio relative to the intrinsic one. A population of Wolf-Rayet stars as origin of the $`[O\text{iv}]`$ emission seems unlikely. Lutz et al. (Lutz98 (1998)) have shown that the $`[O\text{iv}]`$ emission would have to originate in widely dispersed small H ii regions and would have to be relatively strong. $`[O\text{iv}]`$ emission at this level has not been observed in local star forming regions. A conservative analysis will hence not attribute the $`[O\text{iv}]`$ emission in NGC 4945 to the AGN nor to a population of Wolf-Rayet stars. The limits on high excitation AGN tracers are consistent with several scenarios, or perhaps more likely a combination of them: * The Narrow Line Region is extremely obscured even in the mid-IR. We derive an A$`{}_{V}{}^{}`$160 (A(7.65$`\mu `$m)=A(14.3$`\mu `$m)=A(24.3$`\mu `$m)$``$4.3) to the NLR from a comparison of Circinus and NGC 4945 $`[Ne\text{v}]`$ and $`[Ne\text{vi}]`$ line strengths, under the assumption that the galaxies’ NLRs are similar. The choice for Circinus is motivated in Table 2. * UV photons from the AGN are absorbed close to the nucleus along all lines of sight * The extreme ultraviolet luminosity of the AGN is lower than in Circinus. In comparison to the Circinus SED, this would imply a large deficiency in UV relative to X-ray flux (Table 2). ### 3.2 Starburst properties Near-infrared broad-band and emission-line imaging has revealed the nucleus of NGC 4945 to be the site of a sizeable starburst, the presence of which is illustrated by the conically shaped starburst superwind-blown cavity traced at many near-infrared wavelengths (Moorwood et al. 1996a ; Marconi et al. Marconi (2000)). Hampered by the large extinction even in the near-infrared, age estimates for the nuclear starburst are sparse and intrinsically uncertain. ISO-SWS offers the possibility for the first time to observe the mid-infrared line ratio $`[Ne\text{iii}]`$ 15.56$`\mu `$m/$`[Ne\text{ii}]`$ 12.81$`\mu `$m. This ratio, which is much less affected by extinction than visible and UV lines, is sensitive to the hardness of the stellar radiation field and hence is a good indicator for the age of the nuclear starburst. We observed the two lines in the same ISO-SWS aperture, which was centered on the nucleus (see Table 1). To estimate the extinction to the NGC 4945 nuclear starburst we use the ratio of the 18.71$`\mu `$m and 33.48$`\mu `$m $`[`$S iii$`]`$ lines. This ratio is commonly used as a density diagnostic for the density range 10<sup>2.5</sup>–10<sup>4.5</sup> cm<sup>-3</sup> and is only mildly dependent on the temperature of the emitting gas. Assuming a typical starburst gas density of 300 cm<sup>-3</sup> (Kunze et al. Kunze (1996); Rigopoulou et al. Rigopoulou96 (1996)), the intrinsic ratio should be $``$0.43 (i.e. the value in the low density limit, computed using the collision strengths of Tayal & Gupta Tayal (1999)). The observed ratio is far lower: 0.14$`\pm `$0.06. We hence deduce a screen extinction of A(18.7$`\mu `$m)=1.7$`{}_{0.5}{}^{}{}_{}{}^{+0.8}`$, which, using the Galactic center extinction law of Draine (Draine (1989); with A(9.7$`\mu `$m)/E(J–K)=0.71), amounts to A<sub>V</sub>=36$`{}_{11}{}^{}{}_{}{}^{+18}`$. This value is to be considered an upper limit in case the $`[`$S iii$`]`$ emitting area is larger than 14$`\mathrm{}\times `$27$`\mathrm{}`$, in which case an aperture effect causes the $`[`$S iii$`]`$ ratio to be a lower limit. Another independent estimate of the extinction is usually obtained from hydrogen recombination line strengths, assuming ‘case-B’ conditions. For NGC 4945 we therefore observed the 4.05$`\mu `$m Br$`\alpha `$ and the 7.46$`\mu `$m Pf$`\alpha `$ line. Both were measured in the same aperture. The ratio Pf$`\alpha `$/Br$`\alpha `$ is 0.25$`\pm `$0.10, whereas ‘case-B’ recombination theory predicts a ratio of 0.32. The extinction at 7.46$`\mu `$m must therefore be similar or slightly larger than at 4.05$`\mu `$m. This indicates that the grain composition is unusual and probably more similar to the composition found in the line of sight towards the Galactic center (Lutz et al. Lutz96 (1996); Lutz Lutz99 (1999)) than found towards other parts of our Galaxy. An extinction towards the NGC 4945 nuclear starburst can therefore at present not be derived from lines in the 4–7$`\mu `$m range. The extinction derived for the nuclear starburst is somewhat larger than the value we derive for the warm molecular hydrogen (see Sect. 3.5). This indicates that the warm molecular hydrogen and nuclear starburst emission are coming from different nuclear components, the latter possibly more enshrouded than the former. With the unusual grain composition in mind, it is striking how well the Galactic Center extinction law fits our molecular hydrogen data, resulting in a smooth excitation diagram, even for the H<sub>2</sub> (0-0) S(3) line in the center of the 9.7$`\mu `$m silicate feature (see Sect. 3.5). We are therefore confident that the extinction correction for the starburst, derived using the $`[`$S iii$`]`$ ratio, is also useful. In order to determine the excitation of the nuclear starburst we apply the extinction correction derived from the $`[S\text{iii}]`$ ratio to the observed $`[Ne\text{iii}]`$/$`[Ne\text{ii}]`$ ratio. The extinction corrections amount to A(12.8$`\mu `$m)=1.51 and A(15.6$`\mu `$m)=1.19. The extinction corrected $`[Ne\text{iii}]`$/$`[Ne\text{ii}]`$ ratio is 0.064$`{}_{0.032}{}^{}{}_{}{}^{+0.037}`$. Thornley et al. (Thornley (2000)) list observed $`[Ne\text{iii}]`$/$`[Ne\text{ii}]`$ ratios for 26 starburst galaxies, all measured in the same ISO-SWS configuration. Clearly, NGC 4945 is among the lowest excitation targets in their sample. Note that the ISO-SWS aperture used is large compared to the typical size scales in starbursts. The ratios listed by Thornley et al. (Thornley (2000)) should therefore be regarded as aperture averaged. For starburst galaxies L<sub>bol</sub>/L<sub>lyc</sub> is another measure of the excitation of star clusters. Depending on the upper mass cut-off, the star formation decay time scale and the age of the clusters, Thornley et al. (Thornley (2000)) modeled L<sub>bol</sub>/L<sub>lyc</sub> to lie between 3 and 200. The measured values for starburst galaxies range between 3 and 50. Below we will determine L<sub>bol</sub>/L<sub>lyc</sub> for the NGC 4945 nuclear starburst. We assume L<sub>bol</sub>=L<sub>IR</sub> (i.e. no AGN contribution to L<sub>IR</sub>) and estimate L<sub>lyc</sub> from the dereddened 4.05$`\mu `$m Br$`\alpha `$ flux. For A(4.05$`\mu `$m)=1.2 (applying the Galactic center law of Draine (Draine (1989)) for A<sub>V</sub>=36$`{}_{11}{}^{}{}_{}{}^{+18}`$) and L<sub>lyc</sub>=670 L<sub>Brα</sub> we find L<sub>lyc</sub>=8$`{}_{4}{}^{}{}_{}{}^{+9}`$$`\times `$10<sup>8</sup>L and L<sub>bol</sub>/L<sub>lyc</sub>=28$`{}_{15}{}^{}{}_{}{}^{+26}`$. Using the 12.81$`\mu `$m $`[Ne\text{ii}]`$ line and the empirical scaling L<sub>lyc</sub>=64 L$`_{Ne\text{ii}}`$ (Genzel et al. Genzel (1998)) a similar result is obtained. Given the variety of possible star forming histories, it is hard to constrain the age of the nuclear starburst (assuming no AGN contribution to L<sub>IR</sub>). However, both excitation diagnostics agree on a low excitation which suggests an evolved burst with an age in excess of 5$`\times `$10<sup>6</sup> years, but would also be consistent with a low IMF upper mass cut-off. Marconi et al. (Marconi (2000)) show that it is possible to construct starburst models which are consistent with their near-infrared observations of NGC 4945, but differ by the total luminosity generated (their Fig. 4). An instantaneous burst would not leave space in the energy budget for a sizable contribution from the hidden AGN, whereas a combination of instantaneous burst and constant star formation would. We would like to point out here that the latter model would be inconsistent with the low $`[Ne\text{iii}]`$/$`[Ne\text{ii}]`$ ratio observed by us. Only their instantaneous burst is in agreement with both the near-infrared and mid-infrared observations. ### 3.3 What powers the nucleus of NGC 4945? The large extinction towards the nuclear starburst and the AGN buried within, makes it very difficult to assess the contributions of either component to the nuclear bolometric luminosity. The optical, near-infrared, mid-infrared and far-infrared spectra of NGC 4945 are entirely consistent with a starburst-like nature: BLR or NLR high-excitation lines are absent; the starburst excitation indicator $`[Ne\text{iii}]`$/$`[Ne\text{ii}]`$ has a starburst-like value; the ratios of 6$`\mu `$m (ISO-PHT-S), S12 or S25 to S60 or S100 fluxes are all very low and consistent with emission from cold dust only. Furthermore, the ratio L<sub>bol</sub>/L<sub>lyc</sub>=28$`{}_{15}{}^{}{}_{}{}^{+26}`$, is perfectly consistent with the low excitation of the starburst as deduced from $`[Ne\text{iii}]`$/$`[Ne\text{ii}]`$. And last, NGC 4945 lies on the radio-far-infrared correlation for starburst galaxies (Forbes & Norris Forbes (1998)). Hence, the starburst might well account for the the entire observed bolometric luminosity. On the other hand, Guainazzi et al. (Guainazzi (2000)), who have observed the AGN in NGC 4945 in hard X-rays, compute the AGN to be able to account for all the bolometric luminosity observed, if it has a typical quasar L<sub>X</sub>/L<sub>bol</sub> ratio. Since there is no such thing as a template AGN spectrum, the conversion factor applied, L<sub>1-10keV</sub>/L$`{}_{\mathrm{bol}}{}^{}`$0.05 (Elvis et al. Elvis (1994)), may have an uncertainty which could easily allow for the NGC 4945 starburst to dominate the bolometric luminosity instead. The same uncertanties surround the accretion rate of the $``$1.6$`\times `$10<sup>6</sup>M black hole inferred from H<sub>2</sub>O maser observations (Greenhill et al. Greenhill (1997)). A high but not implausible rate of 50% of the Eddington rate (L<sub>Edd</sub>$``$4.1$`\times `$10<sup>10</sup> L) would suffice to power the observed bolometric luminosity. Given the wide range of efficiencies inferred for Seyferts, this information does not add anything to identify the dominant power source. In this complex situation with two potentially dominant power sources, the most significant constraint on their relative weight is the total L<sub>bol</sub>/L<sub>lyc</sub> ratio and its implications on the L<sub>bol</sub>/L<sub>lyc</sub> of the starburst component. L$`{}_{}{}^{\mathrm{sb}}{}_{\mathrm{lyc}}{}^{}`$ is directly constrained by observations, but L$`{}_{}{}^{\mathrm{sb}}{}_{\mathrm{bol}}{}^{}`$ changes for different assumptions on the starburst and AGN contributions to the total bolometric luminosity. If there is a significant AGN contribution, (L<sub>bol</sub>/L<sub>lyc</sub>)<sub>sb</sub> will be less than the global value of 28. Values as low as $``$3 which are possible for a zero age massive star population with Salpeter IMF (e.g. Leitherer & Heckman Leitherer (1995)) are strongly inconsistent with the low excitation observed in NGC 4945. Thornley et al. (Thornley (2000)) model $`[`$Ne iii$`]`$/$`[`$Ne ii$`]`$ and L<sub>bol</sub>/L<sub>lyc</sub> ratios of starbursts, taking into account clusters of different ages within the ISO-SWS aperture. An evolving starburst with $`[`$Ne iii$`]`$/$`[`$Ne ii$`]`$=0.064 as in NGC 4945 must have a L<sub>bol</sub>/L<sub>lyc</sub>$``$15 (their Fig. 8). This limit simply reflects the higher L<sub>bol</sub>/L<sub>lyc</sub> of later type O stars and persists if the low excitation is due to an upper mass cut-off rather than evolution. With a lower limit of $``$15 on (L<sub>bol</sub>/L<sub>lyc</sub>)<sub>sb</sub>, the starburst contribution to the bolometric luminosity must be at least $``$50%. We hence conclude that the AGN in NGC 4945 plays a secondary although most likely not insignificant role in the energetics of this nearby galaxy. Extremely small values for the AGN contribution to the bolometric luminosity would imply an unrealistically high ratio of L<sub>X</sub>/L<sub>bol</sub> for the AGN. The very low inferred black hole mass, the very cold mid-infrared to far-infrared colors, and the absence of any clear line of sight towards the AGN, support our view that starburst activity dominates AGN activity in NGC 4945. ### 3.4 Emission and absorption features The infrared spectrum of the central region of NGC 4945 obtained with ISO-PHT-S (see Fig. 2) presents a new view of the ISM in starburst galaxies. Even at the low spectral resolving power of R$``$90 the spectrum is dominated by a wealth of emission and absorption features. Especially prominent is the family of PAH emission features at 3.3, 6.2, 7.7, 8.6 and 11.3$`\mu `$m, which ISO confirmed to be common-place in most galactic and extragalactic ISM spectra (e.g. Acosta-Pulido et al. Acosta96 (1996); Rigopoulou et al. Rigopoulou99 (1999); Mattila et al. Mattila (1999); Clavel et al. Clavel (2000)). Nevertheless, the weakness of the 8.6 and 11.3$`\mu `$m PAH bands in NGC 4945 is unusual. Consistent with A<sub>V</sub>$``$36 and with the strength of the absorption features discussed below, we explain this weakness by heavy extinction, which will suppress these two features placed in the wings of the silicate absorption feature. Perhaps the most important result, however, is the rich absorption spectrum, indicating that we are observing the infrared sources in the central region of NGC 4945 through a medium containing molecular ices. Interstellar absorptions of 4.27$`\mu `$m (2343cm<sup>-1</sup>) solid CO<sub>2</sub> and 4.68–4.67$`\mu `$m ‘XCN’+CO are detected, the first time in an extragalactic source to our knowledge. At our resolving power and signal-to-noise we cannot determine the contribution of 4.62$`\mu `$m (2165cm<sup>-1</sup>) XCN, 4.67$`\mu `$m (2140cm<sup>-1</sup>) CO ice and of gaseous CO absorptions to the 4.58–4.67$`\mu `$m absorption complex. The strength of the XCN absorption appears to be remarkable, suggestive of ice grain processing in an energetic environment (Lacy et al. Lacy (1984)). We defer a more detailed analysis of the XCN/CO feature to a future paper, which will also include the results of follow-up observations with ISAAC at the VLT. A strong silicate feature is observed around 9.7$`\mu `$m (see also Moorwood & Glass Moorwood84 (1984)). A deep minimum is also detected around 3.0$`\mu `$m, which is suggestive of water ice (or more precise, the O-H stretch) absorption. Table 3 gives column densities for some of the absorption features discussed above. The presence and strength of these absorption features is consistent with the high starburst obscuration derived from the emission lines (but see also Chiar et al. (Chiar (2000)) for variations in the strength of features along lines of sight of similar A<sub>V</sub>). At the resolution of ISO-PHT-S the molecular absorption features in NGC 4945 show striking similarities with the features seen in the ISO-SWS spectrum of the line of sight towards the Galactic center (Lutz et al. Lutz96 (1996); see Fig. 2). Observations at our spectral resolution do however not permit a detailed comparison. Regarding the 4.26$`\mu `$m CO<sub>2</sub> feature it is likely that the feature can be attributed to solid state CO<sub>2</sub>, since high spectral resolution ISO-SWS observations of other sources indicate that the contribution of gaseous CO<sub>2</sub> to the observed feature is very small (see in particular van Dishoeck et al. vanDishoeck (1996)). In the 4.6-4.8$`\mu `$m region, the spectra of NGC 4945 and the Galactic center differ more strongly, and a more careful inspection is required to assess the contributions of gaseous and solid CO and XCN. ISO-SWS spectroscopy of the Galactic center (Lutz et al. Lutz96 (1996)) clearly shows that what we see at low resolution as a relatively shallow and broad feature is in fact dominated by individual lines of gaseous CO. Contributions of a potential underlying solid CO/XCN component are possible but difficult to separate until our high resolution follow-up observations have been executed. ### 3.5 Molecular hydrogen: physical conditions, excitation and mass Near infrared observations of molecular hydrogen emission in NGC 4945 have been reported by several authors over the last 15 years. The most complete set of observations is published by Koornneef & Israel (Koornneef96 (1996)), who observed 8 ro-vibrational transitions with IRSPEC at the ESO NTT. With ISO-SWS and ISO-PHT-S we have extended the number of observed lines from 8 to 14 by observing the pure rotational transitions S(0), S(1), S(2), S(3), S(5) and S(7) as well as the (1-0) Q(3) line. The latter was also observed with IRSPEC and can therefore be used to determine the proper aperture correction factor for the IRSPEC data set. An overview of the observed lines is presented in Table 4. Information on the spatial extent of the H<sub>2</sub> emitting region is only available for the 2.12$`\mu `$m (1-0) S(1) line (Koornneef & Israel Koornneef96 (1996); Moorwood et al. 1996a ; Quillen et al. Quillen (1999); Marconi et al. Marconi (2000)). Based on Fig. 3a of Moorwood et al. (1996a ) we estimate that more than 90% of the (1-0) S(1) emission fits within the smallest ISO-SWS aperture (14$`\mathrm{}\times `$20$`\mathrm{}`$). It is not unreasonable to expect the H<sub>2</sub> emitting area to increase with decreasing H<sub>2</sub> temperature. The aperture sizes used to observe the respective H<sub>2</sub> transitions increase with increasing sensitivity to lower temperature H<sub>2</sub>. Based on this, we will assume in what follows that ISO-SWS and ISO-PHT-S have observed all available warm H<sub>2</sub>. Further to this, all three instruments were centered on the same nuclear position and viewed the nuclear region under more or less similar position angles (see Sect. 2). We will use the ratio of the 1-0 Q(3) line fluxes measured by ISO-SWS and IRSPEC to scale the other IRSPEC lines to the ISO-SWS aperture size. This ratio is 2.33. From the 14 transitions observed it is possible to compute the upper level populations for 12 levels. We assumed the H<sub>2</sub> levels to be optically thin. The excitation diagram in Fig. 3 shows a plot of the natural logarithm of the total number of H<sub>2</sub> molecules (N<sub>u</sub>), divided by the statistical weight (g<sub>u</sub>), in the upper level of each transition detected, versus the energy of that level (E<sub>u</sub>/k). The plot shows the situation after extinction correction (see below). The excitation temperature (T<sub>ex</sub>) of the gas is the reciprocal of the slope of the excitation diagram. If the warm H<sub>2</sub> is in LTE, the excitation temperature directly corresponds to the kinetic temperature. As is clearly visible from Fig. 3, (extinction corrected; see below) the excitation temperature increases monotonically with upper level energy, from 160K for the combination of (0-0) S(0) & S(1) to 2200K for the ro-vibrational lines. In a highly obscured galaxy like NGC 4945, extinction corrections to the H<sub>2</sub> emission will be important. The extinction can be estimated from the H<sub>2</sub> data themselves taking into account that any known excitation mechanism should produce a “smooth” excitation diagram for the pure rotational lines, and that transitions originating in a common upper level should give consistent results. More specifically, we use three criteria: * The excitation temperature should increase monotonically from the lowest to the highest energy levels. This sets limits on the extinction correction for the (0-0) S(3) line in the center of the 9.7$`\mu `$m silicate absorption feature. * The ratio of the (1-0) Q(3) & (1-0) S(1) lines at 2.42 & 2.12$`\mu `$m should be its intrinsic ratio determined by molecular constants only. The same applies to the (1-0) Q(2) & (1-0) S(0) lines at 2.41 & 2.22$`\mu `$m, that originate from identical upper levels too. * In LTE, the upper level populations normalized by the statistical weights should be similar for the 0-0 S(7) & 1-0 Q(3) lines at 5.51 & 2.42$`\mu `$m, which differ by only 4% in upper level energy. We have varied the extinction and tried several extinction laws. We present the most applicable extinction laws here: * Law A: A($`\lambda `$)$`\lambda ^{1.75}`$ for $`\lambda <`$8$`\mu `$m. For $`\lambda >`$8$`\mu `$m we took the Galactic center law of Draine (Draine (1989)), with A(9.7$`\mu `$m)/E(J–K)=0.71 (Roche & Aitken Roche85 (1985)) and E(J–K)=5. * Law B: The same as law ‘A’, except for the range $`\lambda `$=$`[`$2.6,8.8$`]\mu `$m, where we took the extinction law as found towards the Galactic center (Lutz Lutz99 (1999)). In the 3-8$`\mu `$m range this reddening law constitutes a significantly higher extinction than usually assumed. From Fig. 3 and the criteria defined above, moderate extinctions of A<sub>V</sub>=17–23 are clearly preferred. Extinction law A provides a somewhat better fit than extinction law B. None of the 3 solutions gives a good fit to the (1-0) Q(4) data point. In the following analysis, we use the preferrred extinction correction of A<sub>V</sub>=20$`{}_{3}{}^{}{}_{}{}^{+3}`$ and extinction law A. We note that the extinction to the H<sub>2</sub> emitting region is slightly less than that to the starburst $`[H\text{ii}]`$ regions (Sect. 3.2). This plausibly matches the morphological results of Moorwood et al. (1996a ), who find the starburst in an obscured disk, but the H<sub>2</sub> emission extending into a less obscured wind blown cavity. A rough estimate of the amount of warm molecular hydrogen in the nucleus of NGC 4945 can be derived from the level populations of the pure rotational S(0) and S(1) transitions. The excitation temperature for the upper levels of these transitions (J=2 and J=3) is 160K. Assuming the same excitation conditions for the J=0 and J=1 levels, we compute a warm molecular hydrogen mass of 2.4$`\times `$10<sup>7</sup>M. This is 9% of the total H<sub>2</sub> gas mass estimated from CO and 0.7% of the dynamical mass interior to the molecular ring (Bergman et al. Bergman (1992); see below). As already noted, the excitation temperature changes significantly with level energy. This is the consequence of the natural fact that the emitting gas will consist of a mixture of temperatures. The rich NGC 4945 dataset allows us to address this in a more quantitative way. Experiments with fits assuming a number of discrete temperature components lead us to suggest that a power law might give a good representation of the mass distribution as a function of temperature. We obtain a good fit for the following power law: dM/dT=4.43$`\times `$10<sup>15</sup> T<sup>-4.793</sup> M/K. The quality of the fit is shown in Fig. 4. Table 5 gives warm molecular masses for several low temperature cut-offs. Since H<sub>2</sub> gas at temperatures below 70K does not contribute to the (0-0) S(0) flux, nor to any of the other pure rotational lines, we cannot verify whether our power law mass distribution continues down to the lowest temperatures. Nevertheless, we have included mass estimates down to a low temperature cut-off of 50K. This number is reasonable, since we don’t expect the giant molecular clouds (GMCs) to be as cold as in the Galactic disk (10–20K). Rather we expect conditions as found near the Galactic center, where the GMCs are believed to have temperatures exceeding 50K (Armstrong & Barrett Armstrong (1985)). It is interesting to compare our warm molecular hydrogen gas mass estimate with values found in the literature (see Moorwood & Oliva Moorwood94 (1994) for a review). Bergman et al. (Bergman (1992)), using the inner molecular rotation curve of Whiteoak et al. (Whiteoak (1990)), compute a dynamical mass interior to the molecular ring (R$``$280pc=15.6$`\mathrm{}`$) of 3.3$`\times `$10<sup>9</sup>M. The same authors use CO to derive a total H<sub>2</sub> gas mass of 2.7$`\times `$10<sup>8</sup>M for the ring, assuming the rather high kinetic gas temperature of 100K. Note that a low temperature cut-off of the order 50–60K in our H<sub>2</sub> temperature distribution would bring our estimate of the total H<sub>2</sub> gas mass in agreement with that derived from the low level CO observations. The total H<sub>2</sub> gas mass of 2.7$`\times `$10<sup>8</sup>M agrees well with a starburst-like position of NGC 4945 in the L<sub>IR</sub>–M(H<sub>2</sub>) diagram (Young & Scoville Young (1991)). In Table 6 we list for a number of external galaxies and Galactic template sources temperatures and masses of the warm molecular hydrogen gas. With a value of 9%, NGC 4945 has a warm H<sub>2</sub> gas fraction similar to that found for most of the other galaxies listed. Note however that the results for NGC 3256, NGC 4038/39 and Arp 220 are less well constrained than for NGC 6946 and NGC 4945: only for the latter two can the temperature of the warm H<sub>2</sub> gas be determined from the (0-0) S(0) and S(1) transitions directly. For the same reason a comparison of the H<sub>2</sub> gas temperatures is of limited value unless they are derived from identical line combinations. Fairly low temperatures can be derived from the (0-0) S(0) and S(1) lines, 160K and 179K for NGC 4945 and NGC 6946, respectively. Limits for other galaxies listed in Table 6 are consistent with a similarly low temperature. This temperature is well below that observed for an Orion type shock ($`>`$500K). It is closer to what is observed for the same line combination in PDRs (e.g. Orion Bar: 155K, D. Rosenthal priv. comm.; S140: 159K, Draine & Bertoldi, DraineBertoldi (1999)). While a variety of regions may contribute to the galaxy-integrated temperature distribution, this comparison clearly shows Orion-like shocks to be not representative for the entire emission, and fairly normal PDRs (or less energetic shocks) to be perhaps more typical. If excited by shocks (as suggested by the morphology, Moorwood et al. 1996a ), then the near-infrared H<sub>2</sub> emission in NGC 4945 must trace a small subcomponent of faster shocks. ## 4 Conclusions The main results of this paper can be summarized as follows: * The nuclear starburst is heavily obscured by 36$`{}_{11}{}^{}{}_{}{}^{+18}`$ mag. of visual extinction, as infered from the $`[S\text{iii}]`$ 18.7$`\mu `$m/33.5$`\mu `$m ratio. * The excitation of the nuclear starburst is very low, as deduced from excitation indicators $`[Ne\text{iii}]`$15.56$`\mu `$m/$`[Ne\text{ii}]`$12.81$`\mu `$m and L<sub>bol</sub>/L<sub>lyc</sub>, consistent with an age of at least 5$`\times `$10<sup>6</sup>yrs. Comparison with starburst models implies that at least 50% of the bolometric luminosity is powered by the starburst. * The very low inferred black hole mass, the very cold mid-infrared to far-infrared colors, and the absence of any free line of sight to the NLR supports the conclusion that the starburst dominates the bolometric luminosity. * Our mid-infrared ISO spectroscopy does not provide any evidence for the existence of an AGN in the nucleus of NGC 4945. The only high excitation line detected, the 25.9$`\mu `$m $`[O\text{iv}]`$ line, is most likely produced in shocks associated with the nuclear starburst. * The AGN, detected in hard X-rays, is unusual in not revealing itself at optical, near-infrared and mid-infrared wavelengths. Hence, either the NLR is extremely obscured (A$`{}_{V}{}^{}>`$160), or UV photons from the AGN are absorbed close to the nucleus along all lines of sight, or the AGN is deficient in UV relative to its X-ray flux. * Many ISM solid state and molecular features have been observed with ISO-PHT-S in the 2.4–11.7$`\mu `$m range. Most prominent in emission are the PAH features at 3.3, 6.2, 7.7 and 11.2$`\mu `$m. The strongest absorption features are those of water ice, CO<sub>2</sub> and CO, seen against the nuclear spectrum. These features show striking similarities to the absorption features seen towards the Galactic center. * We have studied the physical conditions, excitation and mass of warm H<sub>2</sub>, combining IRSPEC and ISO observations of 14 transitions. We derive a visual extinction of 20$`{}_{3}{}^{}{}_{}{}^{+3}`$ mag. to the H<sub>2</sub> emitting region. From the (0-0) S(0)& S(1) lines, we compute a warm (160K) H<sub>2</sub> gas mass of 2.4$`\times `$10<sup>7</sup>M, 9% of the total gas mass inferred from CO. The excitation diagram is best fitted by a power law of the form dM/dT=4.43$`\times `$10<sup>15</sup> T<sup>-4.793</sup> M/K. The low excitation temperature of 160K shows Orion-like shocks not to be representative for the entire emission, and fairly normal PDRs to be perhaps more typical. ###### Acknowledgements. The authors wish to thank Dietmar Kunze and Fred Lahuis for help in the SWS datareduction and Matt Lehnert, Steve Lord and Eckhard Sturm for stimulating discussions.
warning/0003/hep-ph0003234.html	ar5iv	text	# Revisiting the UA(1) problems ## I Introduction Various statements of “U<sub>A</sub>(1) problems” related to the eta mesons and their possible resolutions have appeared in the literature now for almost three decades. Considerations of the eta U<sub>A</sub>(1) vacuum Ward identity were discussed by Glashow , Weinberg , Crewther and collaborators. The U<sub>A</sub>(1) axial current and its anomalous addition were studied by Kogut and Susskind . Semiclassical instantons with topological winding number were used by ’t Hooft . Lastly, the large $`N_c`$ limit together with the $`\theta `$ vacuum were explored by Witten . All of the above notions were invoked to resolve the U<sub>A</sub>(1) problem. These above U<sub>A</sub>(1) problems have recently been rekindled by ’t Hooft in his text . We prefer to focus on two U<sub>A</sub>(1)-type problems that have empirical resolutions and which also have a theoretical basis: 1. Goldstone boson structure of the observed $`\eta (547)`$ and $`\eta ^{}(958)`$ mesons via $`\eta `$$`\eta ^{}`$ mixing in the context of QCD. 2. Observed eta hadronic decay rates: 1. $`\mathrm{\Gamma }(\eta 3\pi ^0)=380\pm 36`$ eV appears large since it should vanish by the Sutherland theorem , or be a factor of two smaller in the context of chiral perturbation theory . 2. $`\mathrm{\Gamma }(\eta ^{}3\pi ^0)=313\pm 58`$ eV appears relatively suppressed because $`\eta ^{}3\pi ^0`$ phase space is six times larger than for $`\eta 3\pi ^0`$. 3. $`\mathrm{\Gamma }(\eta ^{}\eta \pi \pi )=131\pm 8`$ keV is a strong decay, whereas the smaller $`3\pi `$ decays in 2a, 2b above change isospin by one unit and are non-strong decays proceeding through the quark mass difference $`m_dm_u`$. 4. We invoke the $`\mathrm{\Delta }I=1`$ $`u_3=\overline{q}\lambda _3q`$ Coleman-Glashow (CG) quark tadpole to support the current-current Sutherland suppression of the $`\eta 3\pi `$ decay rates. The CG tadpole also explains all 13 hadron ($`P`$, $`V`$, $`B`$, $`D`$) SU(2) mass splittings . Then we use PCAC Consistency to compute the $`\eta `$, $`\eta ^{}3\pi ^0`$ decay rates in 2a, 2b above. The above problems are analyzed in Secs. II and III primarily on the basis of the input from meson phenomenology. However, the underlying notions of the quark model are also crucial in this analysis. Therefore, in Sec. IV we show the consistency of some of the results of Secs. II and III with a sophisticated quark model which has strong and clear connections with the fundamental theory – QCD. It is based on the so-called coupled Schwinger-Dyson (SD) and Bethe-Salpeter (BS) approach in which one, by solving the SD equation for dressed quark propagators of various flavors, explicitly constructs constituent quarks. They in turn build $`q\overline{q}`$ meson bound states which are solutions of the BS equation employing the dressed quark propagator obtained as the solution of the SD equation. If the SD and BS equations are so coupled in a consistent approximation, the light pseudoscalar mesons are simultaneously the $`q\overline{q}`$ bound states and the (quasi) Goldstone bosons of dynamical chiral symmetry breaking (D$`\chi `$SB). The resulting relativistically covariant constituent quark model (such as the variant of Ref. ) is consistent with current algebra because it incorporates the correct chiral symmetry behavior thanks to D$`\chi `$SB obtained in an, essentially, Nambu–Jona-Lasinio fashion, but the former model interaction is less schematic. In Refs. for example, it is combined nonperturbative and perturbative gluon exchange; the effective propagator function is the sum of the known perturbative QCD contribution and the modeled nonperturbative component. For details, we refer to Refs. , while here we just note that the momentum-dependent dynamically generated quark mass functions $`_f(q^2)`$ (i.e., the quark propagator SD solutions for quark flavors $`f`$) illustrate well how the coupled SD-BS approach provides a modern constituent model which is consistent with perturbative and nonperturbative QCD. For example, the perturbative QCD part of the gluon propagator leads to the deep Euclidean behaviors of quark propagators (for all flavors) consistent with the asymptotic freedom of QCD . However, what is important in the present paper, is the behavior of the same mass functions $`_f(q^2)`$ for low momenta \[$`q^2=0`$ to $`q^2(400\mathrm{MeV})^2`$\], where $`_f(q^2)`$ (due to D$`\chi `$SB) have values consistent with typical values of the constituent mass parameter in constituent quark models. For the (isosymmetric) $`u`$\- and $`d`$-quarks, our concrete model choice gives us $`_{u,d}(0)=356`$ MeV in the chiral limit (i.e., with vanishing $`\stackrel{~}{m}_{u,d}`$, the explicit chiral symmetry breaking bare mass term in the quark propagator SD equation, resulting in vanishing pion mass eigenvalue, $`m_\pi =0`$, in the BS equation), and $`_{u,d}(0)=375`$ MeV \[just 5% above $`_{u,d}(0)`$ in the chiral limit\] with the explicit chiral symmetry breaking bare mass $`\stackrel{~}{m}_{u,d}=3.1`$ MeV, leading to a realistically light pion, $`m_\pi =140.4`$ MeV. Similarly, for the $`s`$ quark, $`_s(0)=610`$ MeV. The simple-minded constituent mass parameters, denoted below by $`\widehat{m}`$ in the case of the isosymmetric $`u`$ and $`d`$ quarks, and by $`m_s`$ in the case of the $`s`$ quarks, have therefore close analogues in the coupled SD-BS approach which explicitly incorporates some crucial features of QCD, notably D$`\chi `$SB. ## II Goldstone structure of eta mesons To resolve U<sub>A</sub>(1) problem one, we invoke the U(3) pseudoscalar nonet structure $`(\stackrel{}{\pi },K,\eta `$, $`\eta ^{})`$ along with the Gell-Mann-Okubo mass formula $`m_\pi ^2+3m_{\eta _8}^2=4m_K^2`$, requiring an octet eta mass $`m_{\eta _8}567`$ MeV. While this $`\eta _8`$ mass is presumed to vanish in the SU(3) $`\times `$ SU(3) chiral limit (CL), the companion singlet $`\eta _0`$ mass is not expected to vanish in the CL. Using the standard relation mixing $`\eta ,\eta ^{}`$ to $`\eta _8,\eta _0`$ away from the CL one knows $$m_{\eta _8}^2+m_{\eta _0}^2=m_\eta ^2+m_\eta ^{}^21.22\text{GeV}^2,\text{or}m_{\eta _0}947\text{MeV}$$ (1) for masses $`\eta (547)`$, $`\eta ^{}(958)`$, $`\eta _8(567)`$. Here we assumed the standard, most traditional representation of the physical isoscalar pseudoscalars $`\eta `$ and $`\eta ^{}`$ as the orthogonal mixture $`\|\eta `$ $`=`$ $`\mathrm{cos}\theta \|\eta _8\mathrm{sin}\theta \|\eta _0,`$ (3) $`\|\eta ^{}`$ $`=`$ $`\mathrm{sin}\theta \|\eta _8+\mathrm{cos}\theta \|\eta _0,`$ (4) of the respective octet and singlet isospin zero states, $`\eta _8`$ and $`\eta _0`$. In the flavor SU(3) quark model, they are defined through quark–antiquark ($`q\overline{q}`$) basis states $`\|f\overline{f}`$ ($`f=u,d,s`$) as $`\|\eta _8`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}(\|u\overline{u}+\|d\overline{d}2\|s\overline{s}),`$ (6) $`\|\eta _0`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}(\|u\overline{u}+\|d\overline{d}+\|s\overline{s}).`$ (7) The exact SU(3) flavor symmetry ($`m_u=m_d=m_s`$) is nevertheless badly broken. It is an excellent approximation to assume the exact isospin symmetry ($`m_u=m_d`$), and a good approximation to take even the chiral symmetry limit ($`m_u=m_d=0`$) for $`u`$ and $`d`$-quark, but for a realistic description, the strange quark mass $`m_s`$ must be significantly heavier than $`m_u`$ and $`m_d`$. \[In particular, the CL is obviously phenomenologically unrealistic in the strange sector, although it is qualitatively meaningful, and in fact useful as a theoretical limit in the discussions of the U<sub>A</sub>(1) problem.\] Thus, with $`\|u\overline{u}`$ and $`\|d\overline{d}`$ being practically chiral states as opposed to a significantly heavier $`\|s\overline{s}`$, Eqs. (II) do not define the octet and singlet states of the exact SU(3) flavor symmetry, but the effective octet and singlet states. Hence, as in Ref. for example, only in the sense that the same $`q\overline{q}`$ states $`\|f\overline{f}`$ ($`f=u,d,s`$) appear in both Eq. (3) and Eq. (4) do these equations implicitly assume nonet symmetry (as pointed out by Gilman and Kauffman , following Chanowitz, their Ref. ). However, in order to avoid the U<sub>A</sub>(1) problem, this symmetry must ultimately be broken at least at the level of the masses. In particular, it must be broken in such a way that $`\eta \eta _8`$ becomes massless but $`\eta ^{}\eta _0`$ remains massive (as in Ref. ) when CL is taken for all three flavors, $`m_u,m_d,m_s0`$. Alternatively, one can work in a nonstrange (NS)–strange (S) basis $`\|\eta _{\text{NS}}`$ and $`\|\eta _\text{S}`$, where $`\|\eta _{\text{NS}}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\|u\overline{u}+\|d\overline{d})={\displaystyle \frac{1}{\sqrt{3}}}\|\eta _8+\sqrt{{\displaystyle \frac{2}{3}}}\|\eta _0,`$ (9) $`\|\eta _\text{S}`$ $`=`$ $`\|s\overline{s}=\sqrt{{\displaystyle \frac{2}{3}}}\|\eta _8+{\displaystyle \frac{1}{\sqrt{3}}}\|\eta _0.`$ (10) In analogy with Eq. (1), in this basis one finds $$m_{\eta _{NS}}^2+m_{\eta _S}^2=m_\eta ^2+m_\eta ^{}^21.22\text{GeV}^2,$$ (11) since $`\eta ^{}\|\eta =0`$ and since the NS–S mixing relations are $`\|\eta `$ $`=`$ $`\mathrm{cos}\varphi \|\eta _{\text{NS}}\mathrm{sin}\varphi \|\eta _S,`$ (13) $`\|\eta ^{}`$ $`=`$ $`\mathrm{sin}\varphi \|\eta _{\text{NS}}+\mathrm{cos}\varphi \|\eta _S.`$ (14) The more familiar singlet-octet state mixing angle $`\theta `$, defined by Eqs. (II), is geometrically related to the NS–S state mixing angle $`\varphi `$ above as $$\theta =\varphi \mathrm{arctan}\sqrt{2}=\varphi 54.74^{}.$$ (15) Although mathematically equivalent to the $`\eta _8`$$`\eta _0`$ basis, the NS–S mixing basis is more suitable for some quark model considerations, being more natural in practice when the symmetry between the NS and S sectors is broken as described in the preceding passage. There is also another important reason to keep in mind the $`\|\eta _{\text{NS}}`$-$`\|\eta _S`$ state mixing angle $`\varphi `$. This is because it offers the quickest way to show the consistency of our procedures and the corresponding results obtained using just one ($`\theta `$ or $`\varphi `$) state mixing angle, with the two-mixing-angle scheme considered in Refs. , which is defined with respect to the mixing of the decay constants. Namely, in the NS–S basis, the two decay-constant-mixing angles in the two-angle scheme ($`\varphi _q`$ and $`\varphi _s`$) are very close to each other. This basis thus has the advantage that even in the case of the decay-constant-mixing, the reduction to just one mixing angle occurs in a good approximation and one can use just one mixing angle around $`\varphi _q\varphi _s`$. Moreover, phenomenology seems to justify the central assumption of Feldmann, Kroll and Stech (FKS) that in the NS–S basis, the decay constants follow (in a good approximation) the pattern of particle state mixing, so that the NS–S state mixing angle $`\varphi \varphi _q\varphi _s`$. On the other hand, when we express our results in terms of $`\varphi `$ through the relation (15), which always holds since we use a state mixing angle, we find they are consistent with the FKS scheme. The application of the two-mixing-angle scheme is relegated to the Appendix, where we give our predictions for the $`\eta `$$`\eta ^{}`$ decay constants, as well as $`\theta _8`$ and $`\theta _0`$, the two decay-constant-mixing angles (in the octet-singlet basis), in that scheme. Here, we simply note that our considerations will ultimately lead us to $`\varphi 42^{}`$, practically the same as the result of FKS scheme and theory (e.g., in Table 2 of Feldmann’s review ), and in agreement with data. World $`\eta `$$`\eta ^{}`$ mixing angle data in 1989 led to $$\varphi =41^{}\pm 2^{}\text{or}\theta =14^{}\pm 2^{}.$$ (16) A more recent detailed analysis based on 1996 data for decays tensor to pseudoscalar $`TPP`$, radiative vector to pseudoscalar (or vice versa) $`VP\gamma `$, $`PV\gamma `$, double radiative decays $`\eta \gamma \gamma `$, $`\eta ^{}\gamma \gamma `$, and $`J/\psi VP`$ decays (14 such decays) leads to the empirical $`\eta `$$`\eta ^{}`$ mixing angles $$\varphi =43^{}\pm 5^{}\text{or}\varphi =42^{}\pm 3^{}$$ (17) found respectively from observed branching ratios, $`B(a_2\eta \pi /K\overline{K})=2.96\pm 0.53`$, $`B(a_2\eta ^{}\pi /\eta \pi )=0.037\pm 0.007`$, in complete agreement with (16). The $`\eta `$$`\eta ^{}`$ mixing angles in (16) or (17) (for 4 of 14 determinations) depend on the constituent quark mass ratio $`m_s/\widehat{m}1.45`$, as already found from baryon magnetic moments , meson charge radii and $`K^{}K\gamma `$ decays . ($`\widehat{m}`$ denotes the isosymmetric average mass $`m_{u,d}`$.) As for a theoretical determination of the $`\eta `$$`\eta ^{}`$ mixing angle $`\varphi `$ or $`\theta =\varphi 54.74^{}`$, we follow the path of Refs. . The contribution of the gluon axial anomaly to the singlet $`\eta _0`$ mass is essentially just parameterized and not really calculated, but some useful information can be obtained from the isoscalar $`q\overline{q}`$ annihilation graphs of which the “diamond” one in Fig. 1 is just the simplest example. That is, we can take Fig. 1 in the nonperturbative sense, where the two-gluon intermediate “states” represent any even number of gluons when forming a C<sup>+</sup> pseudoscalar $`\overline{q}q`$ meson , and where quarks, gluons and vertices can be dressed nonperturbatively, and possibly include gluon configurations such as instantons. Factorization of the quark propagators in Fig. 1 characterized by the ratio $`X\widehat{m}/m_s`$ leads to the pseudoscalar mass matrix in the NS–S basis $$\left(\begin{array}{cc}m_\pi ^2+2\beta \hfill & \sqrt{2}\beta X\hfill \\ \sqrt{2}\beta X\hfill & 2m_K^2m_\pi ^2+\beta X^2\hfill \end{array}\right)\left(\begin{array}{cc}m_\eta ^2\hfill & 0\hfill \\ 0\hfill & m_\eta ^{}^2\hfill \end{array}\right),$$ (18) where $`\beta `$ denotes the total annihilation strength of the pseudoscalar $`q\overline{q}`$ for the light flavors $`f=u,d`$, whereas it is assumed attenuated by a factor $`X`$ when a $`s\overline{s}`$ pseudoscalar appears. (The mass matrix in the $`\eta _8`$-$`\eta _0`$ basis reveals that in the $`X1`$ limit, the CL-nonvanishing singlet $`\eta _0`$ mass is given by $`3\beta `$.) The two parameters on the left-hand-side (LHS) of (18), $`\beta `$ and $`X`$, are determined by the two diagonalized $`\eta `$ and $`\eta ^{}`$ masses on the RHS of (18). The trace and determinant of the matrices in (18) then fix $`\beta `$ and $`X`$ to be $$\beta =\frac{(m_\eta ^{}^2m_\pi ^2)(m_\eta ^2m_\pi ^2)}{4(m_K^2m_\pi ^2)}0.28\text{GeV}^2,X0.78,$$ (19) with the latter value suggesting a constituent quark mass ratio $`X^1m_s/\widehat{m}1.3`$ , near the values in Refs. , $`m_s/\widehat{m}1.45`$. This fitted nonperturbative scale of $`\beta `$ in (19) depends only on the gross features of QCD. If instead one treats the QCD graph of Fig. 1 in the perturbative sense of literally two gluons exchanged, then one obtains only $`\beta _{2g}0.09\text{GeV}^2`$, which is about $`1/3`$ of the needed scale of $`\beta `$ found in (19). (This indicates that just the perturbative “diamond” graph can hardly represent even the roughest approximation to the effect of the gluon axial anomaly operator $`ϵ^{\alpha \beta \mu \nu }G_{\alpha \beta }^aG_{\mu \nu }^a`$.) The above fitted quark annihilation (nonperturbative) scale $`\beta `$ in (19) can be converted to the NS–S $`\eta `$$`\eta ^{}`$ mixing angle $`\varphi `$ in (11) from the alternative mixing relation $`\mathrm{tan}2\varphi =2\sqrt{2}\beta X(m_{\eta _S}^2m_{\eta _{NS}}^2)^19.2`$ to $$\varphi =\mathrm{arctan}\left[\frac{(m_\eta ^{}^22m_K^2+m_\pi ^2)(m_\eta ^2m_\pi ^2)}{(2m_K^2m_\pi ^2m_\eta ^2)(m_\eta ^{}^2m_\pi ^2)}\right]^{1/2}41.9^{}.$$ (20) This kinematical QCD mixing angle (20) or $`\theta =\varphi 54.74^{}12.8^{}`$ has dynamical analogs , namely the coupled SD-BS approach mentioned in the Introduction and used in Sec. IV below. Since this predicted $`\eta `$$`\eta ^{}`$ mixing angle in (20) is compatible with the empirical values in (16) and (17), we use (20) in the mixing angle relations (11) to infer the nonstrange and strange $`\eta `$ masses, $`m_{\eta _{NS}}^2=\mathrm{cos}^2\varphi m_\eta ^2+\mathrm{sin}^2\varphi m_\eta ^{}^2(758\text{MeV})^2`$ (22) $`m_{\eta _S}^2=\mathrm{sin}^2\varphi m_\eta ^2+\mathrm{cos}^2\varphi m_\eta ^{}^2(801\text{MeV})^2.`$ (23) Thus it is clear that the true physical masses $`\eta (547)`$ and $`\eta ^{}(958)`$ are respectively much closer to the Nambu-Goldstone (NG) octet $`\eta _8(567)`$ and the non-NG singlet $`\eta _0(947)`$ configurations than to the nonstrange $`\eta _{\text{NS}}(758)`$ and strange $`\eta _S(801)`$ configurations inferred in Eqs. (20). However, the mean $`\eta `$$`\eta ^{}`$ mass $`(548+958)/2753\text{MeV}`$ is quite near the nonstrange $`\eta _{\text{NS}}(758)`$. But since $`\eta _8(567)`$ appears far from the NG massless limit we must ask: how close is $`\eta _0(947)`$ to the chiral-limiting nonvanishing singlet $`\eta `$ mass? To answer this latter question, return to Fig. 1 and the quark annihilation strength $`\beta 0.28`$ GeV<sup>2</sup> in Eq. (19). These $`\overline{q}q`$ states presumably hadronize into the U<sub>A</sub>(1) singlet state $`\|\eta _0=\|\overline{u}u+\overline{d}d+\overline{s}s/\sqrt{3}`$, for effective squared mass in the SU(3) limit with $`\beta `$ remaining unchanged : $$m_{\eta _0}^2=3\beta (917\text{MeV})^2.$$ (24) This latter CL $`\eta _0`$ mass in (24) is only 3% shy of the exact chiral-broken $`\eta _0(947)`$ mass found in Eq. (1). (Such a 3% CL reduction also holds for the pion decay constant $`f_\pi 93`$ MeV $`90`$ MeV and for $`f_+(0)=10.97`$ , the $`K`$$`\pi `$ $`K_{l3}`$ form factor.) Thus this $`\eta `$$`\eta ^{}`$ mixing resolution of the first U<sub>A</sub>(1) problem is that the physical $`\eta (547)`$ is 97% of the chiral-broken NG boson $`\eta _8(567)`$. Also the mixing-induced CL singlet mass of 917 MeV in (24) is 97% of the chiral-broken singlet $`\eta _0(947)`$ in (1), which in turn is 99% of the physical $`\eta ^{}`$ mass $`\eta ^{}(958)`$. This speaks to Weinberg’s question as to why there is no isoscalar, pseudoscalar Goldstone boson (with mass less than about $`\sqrt{3}m_\pi 240\text{MeV}`$), associated with the spontaneous breakdown of the axial U<sub>A</sub>(1) symmetry. ## III Hadronic eta decays and the U<sub>A</sub>(1) problem As for the second U<sub>A</sub>(1) problem, Weinberg in correctly identified the rapidly varying $`\eta `$ and $`\pi ^0`$ poles for $`\eta 3\pi ^0`$ decay. However, one must also fold in the PCAC consistency approach of Refs. leading to the $`\eta 3\pi ^0`$ amplitude magnitude with $`f_\pi 93`$ MeV, $$\left\|3\pi ^0\left\|H_{\text{em}}\right\|\eta \right\|=(3/2f_\pi ^2)\left\|\pi ^0\left\|H_{\text{em}}\right\|\eta \right\|+𝒪(m_\pi ^2/m_\eta ^2).$$ (26) Here the factor of $`3/f_\pi ^2`$ on the RHS of (26) corresponds to the three successive double pion PCAC reductions, while the factor of 1/2 characterizes Weinberg’s rapidly varying $`\eta `$ and $`\pi ^0`$ pole terms. Also this $`\mathrm{\Delta }I=1`$ $`\eta \pi `$ non-strong transition in (26) reduces to $$\pi ^0\left\|H_{\text{em}}\right\|\eta =\mathrm{cos}\varphi \pi ^0\left\|u_3\right\|\eta _{\text{NS}}=\mathrm{cos}42^{}(\mathrm{\Delta }m_K^2\mathrm{\Delta }m_\pi ^2)3900\mathrm{MeV}^2.$$ (27) In (27) we have invoked the CG $`u_3=\overline{q}\lambda _3q`$ quark tadpole (which is known to explain all $`P`$, $`V`$, $`B`$, $`D`$ hadron SU(2) electromagnetic (em) mass splittings) using the SU(3) form $`\pi ^0\left\|u_3\right\|\eta _{\text{NS}}=\mathrm{\Delta }m_K^2\mathrm{\Delta }m_\pi ^20.0052`$ GeV<sup>2</sup>, where $`\mathrm{\Delta }m_K^2=m_{K^+}^2m_{K^0}^2`$, etc. Also in (27) we have again invoked the $`\eta `$$`\eta ^{}`$ mixing relations (13) with mixing angle predicted by (20). Substituting (27) into (26), one obtains the $`\eta 3\pi ^0`$ amplitude $$\left\|3\pi ^0\left\|H_{\text{em}}\right\|\eta \right\|=(3/2f_\pi ^2)\left\|\pi ^0\left\|H_{\text{em}}\right\|\eta \right\|0.68.$$ (29) As for the experimental $`\eta _{3\pi ^0}`$ decay amplitude, taking a constant matrix element (29) integrated over the Dalitz plot, one predicts an $`\eta 3\pi ^0`$ decay rate $$\mathrm{\Gamma }(\eta _{3\pi ^0})=(816\text{eV})\left\|3\pi ^0\left\|H_{\text{em}}\right\|\eta \right\|^2377\text{eV}.$$ (30) The latter almost perfectly matches the 1998 PDG rate of $`380\pm 36\text{eV}`$ at the central value. Alternatively we can extract the effective constant 3-body matrix elements $`A_a,A_b,A_c`$ from data $`\mathrm{\Gamma }(\eta 3\pi ^0)0.82\left\|A_a\right\|^2\text{keV}0.38\text{keV},`$ (32) $`\mathrm{\Gamma }(\eta ^{}3\pi ^0)5.58\left\|A_b\right\|^2\text{keV}0.31\text{keV},`$ (33) $`\mathrm{\Gamma }(\eta ^{}\eta \pi ^0\pi ^0)1.06\left\|A_c\right\|^2\text{keV}42\text{keV},`$ (34) leading to the dimensionless 3-body amplitudes $$\left\|A_a\right\|0.68,\left\|A_b\right\|0.24,\left\|A_c\right\|6.3.$$ (35) Note that the PCAC amplitude for $`3\pi ^0\left\|H_{em}\right\|\eta `$ in (29) recovers the observed $`\eta 3\pi ^0`$ rate in (30) or equivalently the constant Dalitz plot amplitude forms in (III) give $`\left\|A_a\right\|0.68`$ which was earlier used to predict the $`\eta 3\pi ^0`$ rate in Eqs. (III). This consistency pattern can also be applied to $`\eta ^{}3\pi ^0`$ decay, presumably dominated by $`\eta ^{}\eta \pi ^0\pi ^0`$ followed by an em transition $`\pi ^0\left\|H_{em}\right\|\eta `$: $`\left\|3\pi ^0\left\|H_{\mathrm{em}}\right\|\eta ^{}\right\|=3\left\|\pi ^0\left\|H_{\mathrm{em}}\right\|\eta \pi ^0\pi ^0\eta \|\eta ^{}\right\|(m_\eta ^2m_\pi ^2)^1`$ (37) $`3(3900\text{MeV}^2)(6.3)(281000\text{MeV}^2)^10.26.`$ (38) In (38) we have again used the em scale (27) (three times), the $`\eta `$ propagator on the $`\pi ^0`$ mass shell and the constant amplitude $`\left\|A_c\right\|6.3`$ in (35). The result 0.26 is near the constant amplitude $`\left\|A_b\right\|0.24`$ in (35), or equivalently the $`\eta _{3\pi }^{}`$ decay rate is predicted to be $$\mathrm{\Gamma }(\eta ^{}3\pi ^0)5.58\left\|3\pi ^0\left\|H_{\mathrm{em}}\right\|\eta ^{}\right\|^2\text{keV}377\text{eV},$$ (39) near data $`313\pm 58\text{eV}`$. Finally we consider the strong decays $`\eta ^{}\eta \pi \pi `$, with the charged to neutral pion branching ratio being about 2, as expected via SU(2) symmetry. At first these decays were thought to be controlling the $`\eta `$$`\eta ^{}`$ mixing angle. Now, however, one begins by assuming an $`\eta `$$`\eta ^{}`$ mixing angle \[such as $`\varphi 42^{}`$ or $`\theta 13^{}`$ found earlier in Eqs. (16,17,20)\] , and then attempts to explain the observed $`\eta ^{}\eta \pi \pi `$ rate given in Sec. I. To this end Singh and Pasupathy in Ref. studied the $`\delta =a_0(983)`$ scalar meson pole amplitude in $`\eta ^{}\delta \pi `$, $`\delta \eta \pi `$. Later Deshpande and Truong in also included a scalar meson $`\sigma `$ pole in this analysis with $`\eta ^{}\eta \sigma `$, $`\sigma \pi \pi `$. These second authors in justified introducing this latter $`\sigma `$ in order to mask a soft-pion Adler zero which would drastically alter the $`\pi \pi `$ phase space. In fact the $`\eta ^{}\eta \pi ^0\pi ^0`$ data shows only a small deviation from phase space, with linear amplitude $`A(1+\alpha y)`$ now requiring $`\alpha =0.058\pm 0.013`$, and $`\alpha =0.08\pm 0.03`$ for $`\eta ^{}\eta \pi ^+\pi ^{}`$ decay. Keeping only these two $`\delta `$ and $`\sigma `$ pole terms, we slightly modify Refs. and write this combined $`\eta ^{}\eta \pi ^+\pi ^{}`$ amplitude magnitude as $$A=\left\|A(\eta ^{}\eta \pi ^+\pi ^{})\right\|\left\|\frac{g_{\delta \eta \pi }g_{\eta ^{}\delta \pi }}{m_\delta ^2uim_\delta \mathrm{\Gamma }_\delta }+\frac{g_{\sigma \pi \pi }g_{\eta ^{}\eta \sigma }}{m_\sigma ^2sim_\sigma \mathrm{\Gamma }_\sigma }\right\|.$$ (40) Here the combined $`\delta `$ and $`\sigma `$ pole amplitudes have the same structure as in Ref. except we always (rather than partially) keep the non-narrow widths $`\mathrm{\Gamma }_\delta 100`$ MeV and $`\mathrm{\Gamma }_\sigma 700`$ MeV . Also to estimate the pole denominators in (40), we follow Ref. and take $`m_\delta ^2u2m_\eta ^{}E_12m_\eta ^{}(m_\eta ^{}m_\eta )`$ in the $`\eta ^{}`$ rest frame with $`p_\pi p_\pi ^{}0`$ soft and $`s=\left[6.772.4y\right]m_\pi ^2`$. Finally we choose the nonstrange $`\sigma `$ mass from the recent data analysis of Ref. : $$m_\sigma =400\text{to}900\text{MeV},\text{mean mass}m_\sigma 650\text{MeV}.$$ (41) This is near $`\epsilon `$ (700) used in and is supported by the 1998 PDG tables . Moreover a $`\sigma (650)`$ is generated from linear $`\sigma `$ model (L$`\sigma `$M) dynamics with L$`\sigma `$M coupling constants using the mixing relations (11): $`g_{\delta \eta \pi }=\mathrm{cos}\varphi g_{\delta \eta _{NS\pi }}`$ $`=`$ $`\mathrm{cos}\varphi \left({\displaystyle \frac{m_\delta ^2m_{\eta _{NS}}^2}{2f_\pi }}\right)1.56\text{GeV},`$ (43) $`g_{\eta ^{}\delta \pi }=\mathrm{sin}\varphi g_{\delta \eta _{NS}\pi }`$ $`=`$ $`\mathrm{sin}\varphi \left({\displaystyle \frac{m_\delta ^2m_{\eta _{NS}}^2}{2f_\pi }}\right)1.40\text{GeV},`$ (44) $`g_{\sigma \pi \pi }`$ $`=`$ $`m_{\sigma _{NS}}^2/2f_\pi 2.27\text{GeV},`$ (45) $`g_{\eta ^{}\eta \sigma }`$ $`=`$ $`\mathrm{cos}\varphi \mathrm{sin}\varphi g_{\sigma \pi \pi }1.13\text{GeV}.`$ (46) Note that $`g_{\delta \eta _{NS\pi }}=g_{\sigma \pi \pi }`$ in the chiral limit and also that the $`\eta `$$`\eta ^{}`$ mixing angle used ($`\varphi =41.9^{}`$) is as found from Eq. (20). Substituting the above numerical values back into (40) leads to the $`\eta ^{}\eta \pi ^+\pi ^{}`$ amplitude magnitude $$\left\|A\right\|\left\|\frac{2.20}{0.79i0.10}+\frac{2.57}{0.29i0.46}\right\|6.64.$$ (47) This L$`\sigma `$M prediction in (47) should be compared with the original estimates in of $`\left\|A\right\|8.5`$, $`\alpha 0.012`$. Also, $`\left\|A\right\|6.64`$ in (47) is near $`\left\|A_c\right\|6.3`$ in (35) assuming a constant matrix element and isospin invariance. Lastly accounting for the $`\eta ^{}\eta \pi ^0\pi ^0`$ as well as the $`\eta ^{}\eta \pi ^+\pi ^{}`$ amplitude, and folding in the slight Dalitz plot slope we predict the total decay rate (for the average slope $`\alpha 0.07`$): $`\mathrm{\Gamma }(\eta ^{}\eta \pi \pi )`$ $`=`$ $`3\left\|A\right\|^2(1+0.24\alpha +0.27\alpha ^2)\text{keV}`$ (49) $``$ $`130\text{keV}.`$ (50) This prediction (50) is in very good agreement with present data ($`131\pm 8`$ keV) as given in Sec. I. We differ from Ref. primarily in that we use the L$`\sigma `$M meson-meson couplings in Eqs. (41). An extraction of the $`\delta \eta \pi `$ coupling from the width of $`\mathrm{\Gamma }(\delta \eta \pi )100`$ MeV gives for $`q=321`$ MeV: $$\mathrm{\Gamma }(\delta \eta \pi )=\frac{q\left\|2g_{\delta \eta \pi }\right\|^2}{8\pi m_\delta ^2},\text{or}\left\|g_{\delta \eta \pi }\right\|1.38\text{GeV}.$$ (51) The latter coupling in (51) is reasonably near the L$`\sigma `$M coupling 1.56 GeV in (43). ## IV Consistency with dynamical calculations As pointed out in the Introduction and Sec. II, there is a dynamical approach to the question of the Goldstone boson structure of the mixed $`\eta (547)`$ and $`\eta ^{}(958)`$ mesons , namely the coupled SD-BS approach incorporating some crucial features of QCD, which leads to the similar conclusions on the mixing angle and masses as the analysis in Sec. II. Before addressing its mass matrix, let us see what this approach tells us about the mixing angle that can be inferred from $`\gamma \gamma `$ decays. Since the SD-BS approach incorporates the correct chiral symmetry behavior thanks to D$`\chi `$SB and is consistent with current algebra, it reproduces (when care is taken to preserve the vector Ward-Takahashi identity of QED) the Abelian axial anomaly results, which are otherwise notoriously difficult to reproduce in bound-state approaches, as discussed in Ref. . This gives particular weight to the constraints placed on the mixing angle $`\theta `$ by the SD-BS results on $`\gamma \gamma `$ decays of pseudoscalars. ### A $`\gamma \gamma `$ decays of the bound-state $`\pi ^0,\eta ,\eta ^{}`$ We express the broken–SU(3) pseudoscalar states $`\pi ^0,\eta _8`$ and $`\eta _0`$ through the quark basis states $`\|f\overline{f}`$ by $$\|P=\underset{f}{}\left(\frac{\lambda ^P}{\sqrt{2}}\right)_{ff}\|f\overline{f},(f=u,d,s),$$ (52) where $`P=\pi ^0,\eta _8,\eta _0`$ simultaneously have the meaning of the respective indices $`j=3,8,0`$ on the SU(3) Gell-Mann matrices $`\lambda ^j(j=1,\mathrm{},8)`$ and on $`\lambda ^0(\sqrt{2/3})\mathrm{𝟏}_3`$. This picks out the diagonal $`\lambda ^3,\lambda ^8,\lambda ^0`$ in Eq. (52). For future convenience we write the $`P(p)\gamma (k)\gamma (k^{})`$ amplitudes as $$T_P(k^2,k^2)=\underset{f}{}\left(\frac{\lambda ^P}{\sqrt{2}}\right)_{ff}Q_f^2\stackrel{~}{T}_{f\overline{f}}(k^2,k^2),$$ (53) where $`\stackrel{~}{T}_{f\overline{f}}(k^2,k^2)T_{f\overline{f}}(k^2,k^2)/Q_f^2`$ are the “reduced” two-photon amplitudes obtained by removing the squared charge factors $`Q_f^2`$ from $`T_{f\overline{f}}`$, the $`\gamma \gamma `$ amplitude of the pseudoscalar quark-antiquark bound state of the hidden flavor $`f\overline{f}`$. The decay amplitudes (into real photons, $`k^2=k^2=0`$) of the physical states $`\eta `$ and $`\eta ^{}`$, are given in terms of the predicted $`\gamma \gamma `$ decay amplitudes of the SU(3) states $`\eta _8`$ and $`\eta _0`$ as $`T_\eta (0,0)`$ $`=`$ $`\mathrm{cos}\theta T_{\eta _8}(0,0)\mathrm{sin}\theta T_{\eta _0}(0,0),`$ (54) $`T_\eta ^{}(0,0)`$ $`=`$ $`\mathrm{sin}\theta T_{\eta _8}(0,0)+\mathrm{cos}\theta T_{\eta _0}(0,0).`$ (55) The best fit to the experimental $`\gamma \gamma `$ decay amplitudes was found in Ref. for $`\theta =12^{}`$ for the concrete SD-BS model and parameters adopted there. In order to show that in the SD-BS approach $`\gamma \gamma `$ decays imply $`\theta `$ somewhere in that ballpark (i.e., less negative than values favored by $`\chi `$PT) regardless of any model choice, and to be able to compare with other theoretical approaches which usually try to express $`P\gamma \gamma `$ amplitudes in terms of the leptonic (axial-current) decay constants $`f_P`$, let us start with the light $`u,d`$ sector in the chiral (and soft) limit. There, the SD-BS approach yields analytically and exactly<sup>§</sup><sup>§</sup>§The same holds for the related process $`\gamma \pi ^+\pi ^0\pi ^{}`$., and independently of the internal bound-state pion structure, $$\stackrel{~}{T}_{\pi ^0}(0,0)\stackrel{~}{T}_{u\overline{u}}(0,0)=\stackrel{~}{T}_{d\overline{d}}(0,0)=\frac{N_c}{2\sqrt{2}\pi ^2f_\pi },$$ (56) $$T_{\pi ^0}(0,0)=\frac{N_c}{2\sqrt{2}\pi ^2f_\pi }\underset{f}{}\left(\frac{\lambda ^3}{\sqrt{2}}\right)_{ff}Q_f^2=\frac{1}{4\pi ^2f_\pi }.$$ (57) Of course, the calculated value of $`f_\pi `$ does depend on the (modeling of the) internal pion structure, but the empirically successful axial-anomaly chiral-limit relation (57) does not. The $`\pi ^0\gamma \gamma `$ decay amplitude for a possibly nonvanishing pion mass, can be used as a definition of pionic $`\gamma \gamma `$-decay constant $`\overline{f}_\pi `$ by demanding that this amplitude be written in the form of the massless, CL amplitude (57), but with $`\overline{f}_\pi `$ in place of $`f_\pi `$: $`T_{\pi ^0}(0,0)=1/4\pi ^2\overline{f}_\pi `$. Obviously, $`\overline{f}_\pi =f_\pi `$ in the CL, and $`\overline{f}_\pi `$ is a convenient way to re-express the $`\gamma \gamma `$ amplitude in the case of a nonvanishing pion mass, because the Veltman-Sutherland theorem, PCAC, and the empirical success of the chiral-limit anomaly result (57), guarantee that $`\overline{f}_\pi f_\pi `$ always holds for any realistic description of the light $`u,d`$ sector. For simplicity of discussion, we therefore use $`\overline{f}_\pi =f_\pi `$ in this subsection, as the Veltman-Sutherland theorem guarantees that this can be wrong only by several percent. Although the chiral limit formula (57) can be applied without reservations only to pions, it is customary to write the amplitudes for $`\eta _8,\eta _0\gamma \gamma `$ in the same form as (57), defining thereby the $`\gamma \gamma `$-decay constants $`\overline{f}_{\eta _8}`$ and $`\overline{f}_{\eta _0}`$: $$T_{\eta _8}(0,0)\frac{N_c}{2\sqrt{2}\pi ^2\overline{f}_{\eta _8}}\underset{f}{}\left(\frac{\lambda ^8}{\sqrt{2}}\right)_{ff}Q_f^2=\frac{f_\pi }{\overline{f}_{\eta _8}}\frac{T_{\pi ^0}(0,0)}{\sqrt{3}},$$ (58) $$T_{\eta _0}(0,0)\frac{N_c}{2\sqrt{2}\pi ^2\overline{f}_{\eta _0}}\underset{f}{}\left(\frac{\lambda ^0}{\sqrt{2}}\right)_{ff}Q_f^2=\frac{f_\pi }{\overline{f}_{\eta _0}}\frac{\sqrt{8}T_{\pi ^0}(0,0)}{\sqrt{3}}.$$ (59) As pointed out by , $`\overline{f}_{\eta _8}`$ and $`\overline{f}_{\eta _0}`$ are not a priori simply connected with the usual axial-current decay constants $`f_{\eta _8}`$ and $`f_{\eta _0}`$, in contrast to $`f_\pi \overline{f}_\pi `$. Expressing $`T_{\eta _8}(0,0)`$ and $`T_{\eta _0}(0,0)`$ through the $`\gamma \gamma `$-decay constants $`\overline{f}_{\eta _8}`$ and $`\overline{f}_{\eta _0}`$, yields the customary (see, e.g. ) forms for the $`\eta `$ and $`\eta ^{}`$ decay widths: $`\mathrm{\Gamma }(\eta \gamma \gamma )`$ $`=`$ $`{\displaystyle \frac{\alpha _{\mathrm{em}}^2}{64\pi ^3}}{\displaystyle \frac{m_\eta ^3}{3f_\pi ^2}}\left[{\displaystyle \frac{f_\pi }{\overline{f}_{\eta _8}}}\mathrm{cos}\theta \sqrt{8}{\displaystyle \frac{f_\pi }{\overline{f}_{\eta _0}}}\mathrm{sin}\theta \right]^2,`$ (60) $`\mathrm{\Gamma }(\eta ^{}\gamma \gamma )`$ $`=`$ $`{\displaystyle \frac{\alpha _{\mathrm{em}}^2}{64\pi ^3}}{\displaystyle \frac{m_\eta ^{}^3}{3f_\pi ^2}}\left[{\displaystyle \frac{f_\pi }{\overline{f}_{\eta _8}}}\mathrm{sin}\theta +\sqrt{8}{\displaystyle \frac{f_\pi }{\overline{f}_{\eta _0}}}\mathrm{cos}\theta \right]^2.`$ (61) The even more customary version of (60) and (61) in which the axial-current decay constants $`f_{\eta _8}`$ and $`f_{\eta _0}`$ appear in place of $`\overline{f}_{\eta _8}`$ and $`\overline{f}_{\eta _0}`$ requires a derivation where PCAC and soft meson technique are applied to the $`\eta `$$`\eta ^{}`$ complex . For the indeed light pion, these assumptions are impeccable (leading to $`f_\pi =\overline{f}_\pi `$), but not for the $`\eta `$$`\eta ^{}`$ complex. For such a heavy particle as $`\eta ^{}`$ they are quite dubious. However, we do not need and do not use these assumptions since we directly calculated the $`\eta _8`$ and $`\eta _0`$ decay amplitudes, i.e., $`\overline{f}_{\eta _8}`$ and $`\overline{f}_{\eta _0}`$, just as the axial-current pseudoscalar decay constants $`f_{\eta _8}`$ and $`f_{\eta _0}`$ were calculated independently of the $`\gamma \gamma `$ processes. In contrast to $`f_\pi =\overline{f}_\pi `$, $`f_{\eta _8}`$ and $`\overline{f}_{\eta _8}`$ cannot be equated, as the difference between them was found to be quite important . The precise values of $`\overline{f}_{\eta _8}`$ and $`\overline{f}_{\eta _0}`$ are model dependent, but $`\overline{f}_{\eta _8}<\overline{f}_\pi f_\pi `$ holds in this approachThis is different from chiral perturbation theory . Nevertheless, for the axial-current decay constants our approach gives $`f_{\eta _8}>f_\pi `$ (see Appendix or Ref. ). Even the numerical value obtained in our concrete SD-BS calculation , $`f_{\eta _8}=1.31f_\pi `$, is rather close to $`f_{\eta _8}=1.25f_\pi `$ obtained in chiral perturbation theory . generally, i.e., independently of chosen model details, as long as the $`s`$-quark mass is realistically heavier than the $`u,d`$-quark masses. To see this, let us start by noting that $`\overline{f}_{\eta _8}<f_\pi `$ is equivalent to $`T_{\eta _8}(0,0)>T_{\pi ^0}(0,0)/\sqrt{3}`$, and since we can re-write Eq. (53) for $`\eta _8`$ as $$T_{\eta _8}(0,0)=\frac{T_{\pi ^0}(0,0)}{\sqrt{3}}+\frac{1}{9}\frac{2}{\sqrt{6}}\left[\stackrel{~}{T}_{d\overline{d}}(0,0)\stackrel{~}{T}_{s\overline{s}}(0,0)\right],$$ (62) the inequality $`\overline{f}_{\eta _8}<f_\pi `$ is in our approach simply the consequence of the fact that the (“reduced”) $`\gamma \gamma `$-amplitude of the $`s\overline{s}`$-pseudoscalar bound state, $`\stackrel{~}{T}_{s\overline{s}}`$, is smaller than the corresponding non-strange $`\gamma \gamma `$-amplitude $`\stackrel{~}{T}_{d\overline{d}}`$ ($`=\stackrel{~}{T}_{u\overline{u}}=\stackrel{~}{T}_{\pi ^0}`$ in the isosymmetric limit), for any realistic relationship between the non-strange and much larger strange quark masses. This is the reason why in this approach one cannot fit well the experimental $`\eta ,\eta ^{}\gamma \gamma `$ widths with the mixing angle as negative as in chiral perturbation theory descriptions ($`\theta 20^{}`$), but rather with $`\theta 12^{}`$. This is easily understood, for example, with the help of Fig. 1. of Ball et al. , where the values of $`\overline{f}_{\eta _{8(0)}}/f_\pi `$ consistent with experiment are given as a function of the mixing angle $`\theta `$. Their curve shows that values $`\overline{f}_{\eta _8}/f_\pi <1`$ permit accurate reproduction of $`\eta ,\eta ^{}\gamma \gamma `$ widths only for $`\theta `$-values less negative than $`15^{}`$. \[It does not matter that they in fact plotted $`f_{\eta _{8(0)}}/f_\pi `$ and not $`\overline{f}_{\eta _{8(0)}}/f_\pi `$. Namely, they used Eqs. (60)-(61) above for comparison with the experimental $`\gamma \gamma `$-widths, just with $`f_{\eta _{8(0)}}/f_\pi `$ instead of $`\overline{f}_{\eta _{8(0)}}/f_\pi `$, so that the experimental constraints displayed in their Fig. 1 apply to whatever ratios are used in these expressions. One should also note that since in our approach $`\overline{f}_{\eta _8},\overline{f}_{\eta _0}`$ and $`f_\pi `$ are not free parameters but predicted quantities, the two widths $`\eta ,\eta ^{}\gamma \gamma `$ cannot be fitted exactly by adjusting just one parameter, $`\theta `$. Rather, we fix $`\theta `$ by performing a $`\chi ^2`$ fit to the widths.\] On the other hand, the more negative values $`\theta \text{}20^{}`$ give good $`\eta ,\eta ^{}\gamma \gamma `$ widths in conjunction with the ratio $`\overline{f}_{\eta _8}/f_\pi =f_{\eta _8}/f_\pi =1.25`$ obtained by in $`\chi `$PT. However, the coupled SD-BS approach belongs among constituent quark approachesHowever, at least one effective–meson–Lagrangian approach, that of Benayoun et al. , yields results quite close to ours: $`\theta =11.59^{}\pm 0.76^{}`$ and their Eq. (29), where their $`(f_8,f_1)`$ correspond to our $`(\overline{f}_{\eta _8},\overline{f}_{\eta _0})`$, with $`f_8/f_\pi =0.82\pm 0.02`$ and $`f_1/f_\pi =1.15\pm 0.02`$. and for them, considerably less negative angles, $`\theta 14^{}\pm 2`$ , are natural. Ref. showed that these bounds and estimates are very robust under SD-BS model variations and can be taken as model independent. For example, for chiral $`u,d`$ quarks, $$\overline{f}_{\eta _8}=\frac{3f_\pi }{5\frac{4\pi ^2\sqrt{2}f_\pi }{N_c}\stackrel{~}{T}_{s\overline{s}}(0,0)},\overline{f}_{\eta _0}=\frac{6f_\pi }{5+\frac{2\pi ^2\sqrt{2}f_\pi }{N_c}\stackrel{~}{T}_{s\overline{s}}(0,0)},$$ (63) leading to the bounds $`\frac{3}{5}f_\pi <\overline{f}_{\eta _8}<f_\pi `$ and $`f_\pi <\overline{f}_{\eta _0}<\frac{6}{5}f_\pi `$. Also, considerations based on the Goldberger–Treiman relation showed that $`\stackrel{~}{T}_{s\overline{s}}(0,0)<\stackrel{~}{T}_{u\overline{u}}(0,0)`$ is simply due to $`f_{s\overline{s}}f_\pi +2(f_{K^+}f_\pi )>f_\pi `$ (where $`f_{s\overline{s}}`$ is the axial-current decay constant of the unphysical $`s\overline{s}`$ pseudoscalar bound state), and that a good estimate of the $`\gamma \gamma `$-amplitude ratio is the inverse ratio of the pertinent constituent quark masses: $`\stackrel{~}{T}_{s\overline{s}}(0,0)/\stackrel{~}{T}_{u\overline{u}}(0,0)\widehat{m}/m_s`$. Equations (63) then give the relations \[reducing to $`\overline{f}_{\eta _8}=f_\pi `$ and $`\overline{f}_{\eta _0}=f_\pi `$ in the U(3) limit, just like Eqs. (63) themselves\] $$\overline{f}_{\eta _8}\frac{3f_\pi }{52\widehat{m}/m_s},\overline{f}_{\eta _0}\frac{6f_\pi }{5+\widehat{m}/m_s},$$ (64) obtained also by Ref. using the simple quark loop model with constant constituent masses. These estimates are (for reasonable $`\widehat{m}/m_s`$) close to what Ref. calculated with a concrete SD-BS model choice , namely $`\overline{f}_{\eta _0}/f_\pi =1.067`$ and $`\overline{f}_{\eta _8}/f_\pi =0.797`$. For these concrete model values, $`\eta ,\eta ^{}\gamma \gamma `$ widths (60)-(61) fit the data best for $`\theta =12.0^{}`$. ### B Introducing $`X`$ into the SD-BS mass matrix For the very predictive SD-BS approach to be consistent, the above mixing angle extracted from $`\eta ,\eta ^{}\gamma \gamma `$ widths, should be close to the angle $`\theta `$ predicted by diagonalizing the $`\eta `$$`\eta ^{}`$ mass matrix. In this subsection, it is given in the quark $`f\overline{f}`$ basis: $$M^2=\text{diag}(M_{u\overline{u}}^2,M_{d\overline{d}}^2,M_{s\overline{s}}^2)+\beta \left[\begin{array}{ccc}1& 1& 1\hfill \\ 1& 1& 1\hfill \\ 1& 1& 1\hfill \end{array}\right].$$ (65) As in Sec. II, $`3\beta `$ (called $`\lambda _\eta `$ in Ref. ) is the contribution of the gluon axial anomaly to $`m_{\eta _0}^2`$, the squared mass of $`\eta _0`$. We denote by $`M_{f\overline{f}^{}}`$ the masses obtained as eigenvalues of the BS equations for $`q\overline{q}`$ pseudoscalars with the flavor content $`f\overline{f}^{}`$ ($`f,f^{}=u,d,s`$). However, since Ref. had to employ a rainbow-ladder approximation (albeit the improved one of Ref. ), it could not calculate the gluon axial anomaly contribution $`3\beta `$. It could only avoid the U<sub>A</sub>(1)-problem in the $`\eta `$$`\eta ^{}`$ complex by parameterizing $`3\beta `$, namely that part of the $`\eta _0`$ mass squared which remains nonvanishing in the CL. Because of the rainbow-ladder approximation (which does not contain even the simplest annihilation graph – Fig. 1), the $`q\overline{q}`$ pseudoscalar masses $`M_{f\overline{f}^{}}`$ do not contain any contribution from $`3\beta `$, unlike the nonstrange and strange $`\eta `$ masses $`m_{\eta _{NS}}`$ \[in Eq. (22)\] and $`m_{\eta _S}`$ \[in Eq. (23)\], which do, and which must not be confused with $`M_{u\overline{u}}=M_{d\overline{d}}`$ and $`M_{s\overline{s}}`$. Since the flavor singlet gluon anomaly contribution $`3\beta `$ does not influence the masses $`m_\pi `$ and $`m_K`$ of the non-singlet pion and kaon, the realistic rainbow-ladder modeling aims directly at reproducing the empirical values of these masses: $`M_{u\overline{u}}=M_{d\overline{d}}=m_\pi `$ and $`M_{s\overline{d}}=m_K`$. In contrast, the masses of the physical etas, $`m_\eta `$ and $`m_\eta ^{}`$, must be obtained by diagonalizing the $`\eta _8`$-$`\eta _0`$ sub-matrix containing both $`M_{f\overline{f}}`$ and the gluon anomaly contribution to $`m_{\eta _0}^2`$. Since the gluon anomaly contribution $`3\beta `$ vanishes in the large $`N_c`$ limit as $`1/N_c`$, while all $`M_{f\overline{f}^{}}`$ vanish in CL, our $`q\overline{q}`$ bound-state pseudoscalar mesons behave in the $`N_c\mathrm{}`$ and chiral limits in agreement with QCD and $`\chi `$PT (e.g., see ): as the strict CL is approached for all three flavors, the SU(3) octet pseudoscalars including $`\eta `$ become massless Goldstone bosons, whereas the chiral-limit-nonvanishing $`\eta ^{}`$-mass $`3\beta `$ is of order $`1/N_c`$ since it is purely due to the gluon anomaly. If one lets $`3\beta 0`$ (as the gluon anomaly contribution behaves for $`N_c\mathrm{}`$), then for any quark masses and resulting $`M_{f\overline{f}}`$ masses, the “ideal” mixing ($`\theta =54.74^{}`$) takes place so that $`\eta `$ consists of $`u,d`$ quarks only and becomes degenerate with $`\pi `$, whereas $`\eta ^{}`$ is the pure $`s\overline{s}`$ pseudoscalar bound state with the mass $`M_{s\overline{s}}`$. In Ref. , numerical calculations of the mass matrix were performed for the realistic chiral and SU(3) symmetry breaking, with the finite quark masses (and thus also the finite BS $`q\overline{q}`$ bound-state pseudoscalar masses $`M_{f\overline{f}}`$) fixed by the fit to static properties of many mesons but excluding the $`\eta `$$`\eta ^{}`$ complex. The mixing angle which diagonalizes the $`\eta _8`$-$`\eta _0`$ mass matrix thus depended in Ref. only on the value of the additionally introduced “gluon anomaly parameter” $`3\beta `$. Its preferred value turned out to be $`3\beta =1.165`$ GeV<sup>2</sup>=(1079 MeV)<sup>2</sup>, leading to the mixing angle $`\theta =12.7^{}`$ \[compatible with $`\varphi =41.9^{}`$ in Eq. (20)\] and acceptable $`\eta \gamma \gamma `$ and $`\eta ^{}\gamma \gamma `$ decay amplitudes. Also, the $`\eta `$ mass was then fitted to its experimental value, but such a high value of $`3\beta `$ inevitably resulted in a too high $`\eta ^{}`$ mass, above 1 GeV. (Conversely, lowering $`3\beta `$ aimed to reduce $`m_\eta ^{}`$, would push $`\theta `$ close to $`20^{}`$, making predictions for $`\eta ,\eta ^{}\gamma \gamma `$ intolerably bad.) However, unlike Eq. (18) in the present paper, it should be noted that Ref. did not introduce into the mass matrix the “strangeness attenuation parameter” $`X`$ which should suppress the nonperturbative quark $`f\overline{f}f^{}\overline{f}^{}`$ annihilation amplitude (illustrated by the “diamond” graph in Fig. 1) when $`f`$ or $`f^{}`$ are strange. On the other hand, the influence of this suppression should be substantial, since $`X\widehat{m}/m_s`$ should be a reasonable estimate of it, and this nonstrange-to-strange constituent mass ratio in the considered variant of the SD-BS approach is not far from $`X`$ in Eq. (19) and from the mass ratios in Refs. , and is even closer to the mass ratios in the Refs. . Namely, two of us found it to be around $`_u(0)/_s(0)=0.615`$ if the constituent mass was defined at the vanishing argument $`q^2`$ of the momentum-dependent SD mass function $`_f(q^2)`$. We therefore introduce the suppression parameter $`X`$ the same way as in the NS–S mass matrix (18), whereby the mass matrix in the $`f\overline{f}`$ basis becomes $$M^2=\text{diag}(M_{u\overline{u}}^2,M_{d\overline{d}}^2,M_{s\overline{s}}^2)+\beta \left[\begin{array}{ccc}1& 1& X\hfill \\ 1& 1& X\hfill \\ X& X& X^2\hfill \end{array}\right].$$ (66) In a very good approximation, Eq. (66) recovers (in the $`\pi ^0`$NS–S basis) Eq. (18) for the $`2\times 2`$ $`\eta `$$`\eta ^{}`$ subspace. This is because $`M_{s\overline{s}}^2`$ differs from $`2m_K^2m_\pi ^2`$ only by a couple of percent, thanks to the good chiral behavior of the masses $`M_{f\overline{f}^{}}`$ calculated in SD-BS approach. (These $`M_{f\overline{f}^{}}^2`$ and the CL model values of $`f_\pi `$ and quark condensate, satisfy Gell-Mann-Oakes-Renner relation to the first order in the explicit chiral symmetry breaking .) The SD-BS–predicted octet (quasi-)Goldstone masses $`M_{f\overline{f}^{}}`$ are known to be empirically successful in our concrete model choice , but the question is whether the SD-BS approach can also give some information on the $`X`$-parameter. If we treat both $`3\beta `$ and $`X`$ as free parameters, we can of course fit both the $`\eta `$ mass and the $`\eta ^{}`$ mass to their experimental values. For the model parameters as in Ref. (for these parameters our independent calculation gives $`m_\pi =M_{u\overline{u}}=140.4`$ MeV and $`M_{s\overline{s}}=721.4`$ MeV), this happens at $`3\beta =0.753`$ GeV<sup>2</sup> =(868 MeV)<sup>2</sup> and $`X=0.835`$. However, the mixing angle then comes out as $`\theta =17.9^{}`$, which is too negative to allow consistency of the empirically found two-photon decay amplitudes of $`\eta `$ and $`\eta ^{}`$, with predictions of our SD-BS approach for the two-photon decay amplitudes of $`\eta _8`$ and $`\eta _0`$ . Therefore, and also to avoid introducing another free parameter in addition to $`3\beta `$, we take the path where the dynamical information from our SD-BS approach is used to estimate $`X`$. Namely, our $`\gamma \gamma `$ decay amplitudes $`T_{f\overline{f}}`$ can be taken as a serious guide for estimating the $`X`$-parameter instead of allowing it to be free. We did point out in Sec. II that the attempted treatment of the gluon anomaly contribution through just the “diamond diagram” contribution to $`3\beta `$, indicated that just this partial contribution is quite insufficient. This limits us to keeping $`3\beta `$ as a free parameter, but we can still suppose that this diagram can help us get the prediction of the strange-nonstrange ratio of the complete pertinent amplitudes $`f\overline{f}f^{}\overline{f}^{}`$ as follows. Our SD-BS modeling in Ref. employs an infrared-enhanced gluon propagator weighting the integrand strongly for low gluon momenta squared. Therefore, in analogy with Eq. (4.12) of Kogut and Susskind (see also Refs. ), we can approximate the Fig. 1 amplitudes $`f\overline{f}2\mathrm{g}\mathrm{l}\mathrm{u}\mathrm{o}\mathrm{n}\mathrm{s}f^{}\overline{f}^{}`$, i.e., the contribution of the quark-gluon diamond graph to the element $`ff^{}`$ of the $`3\times 3`$ mass matrix, by the factorized form $$\stackrel{~}{T}_{f\overline{f}}(0,0)𝒞\stackrel{~}{T}_{f^{}\overline{f}^{}}(0,0).$$ (67) In Eq. (67), the quantity $`𝒞`$ is given by the integral over two gluon propagators remaining after factoring out $`\stackrel{~}{T}_{f\overline{f}}(0,0)`$ and $`\stackrel{~}{T}_{f^{}\overline{f}^{}}(0,0)`$, the respective amplitudes for the transition of the $`q\overline{q}`$ pseudoscalar bound state for the quark flavor $`f`$ and $`f^{}`$ into two vector bosons, in this case into two gluons. The contribution of Fig. 1 is thereby expressed with the help of the (reduced) amplitudes $`\stackrel{~}{T}_{f\overline{f}}(0,0)`$ we calculated for the transition of $`q\overline{q}`$ pseudoscalars to two real photons ($`k^2=k_{}^{}{}_{}{}^{2}=0`$). Although $`𝒞`$ is in principle computable, all this unfortunately does not amount to determining $`\beta ,\beta X`$ and $`\beta X^2`$ in Eq. (66) since the higher (four-gluon, six-gluon, … , etc.) contributions are clearly lacking. We therefore must keep the total (light-)quark annihilation strength $`\beta `$ as a free parameter. However, if we assume that the suppression of the diagrams with the strange quark in a loop is similar for all of them, Eq. (67) and the “diamond” diagram in Fig. 1 help us to at least estimate the parameter $`X`$ as $`X\stackrel{~}{T}_{s\overline{s}}(0,0)/\stackrel{~}{T}_{u\overline{u}}(0,0)`$. This is a natural way to build in the effects of the SU(3) flavor symmetry breaking in the $`q\overline{q}`$ annihilation graphs. We get $`X=0.663`$ from the two-photon amplitudes we obtained in the chosen SD-BS model . This value of $`X`$ agrees well with the other way of estimating $`X`$, namely the nonstrange-to-strange constituent mass ratio of Refs. . With $`X=0.663`$, requiring that the $`2\times 2`$ matrix trace, $`m_\eta ^2+m_\eta ^{}^2`$, be fitted to its empirical value, fixes the chiral-limiting nonvanishing singlet mass squared to $`3\beta =0.832`$ GeV<sup>2</sup>=(912 MeV)<sup>2</sup>, just 0.5% below Eq. (24). The resulting mixing angle and $`\eta `$, $`\eta ^{}`$ masses are $$\theta =13.4^{},m_\eta =588\text{MeV},m_\eta ^{}=933\text{MeV}.$$ (68) The above results of the SD-BS approach are very satisfactory since they agree well with what was found in Sec. II by different methods. Let us close this section by exploring the stability of these results on model variations. Except the introduction of $`3\beta (=\lambda _\eta )`$, these SD-BS results were obtained without any other parameter fitting, with the model parameters resulting from the very broad previous fit , but actually giving us, in our independent calculation, a few percent too high results for $`m_\pi `$ and $`m_K`$. To possibly improve, and in any case check the robustness of the consistency with Sec. II (and subsection IV.A) on variations of our model description, we therefore perform a refitting in the sector of $`u,d`$ and $`s`$ quarks, to reproduce exactly the average isotriplet pion mass $`m_\pi =M_{u\overline{u}}=137.3`$ MeV and isodoublet kaon mass $`m_K=495.7`$ MeV. As Table I shows, the changes are small, and lead to $`_u(0)/_s(0)=0.622`$ and $`X=\stackrel{~}{T}_{s\overline{s}}(0,0)/\stackrel{~}{T}_{u\overline{u}}(0,0)=0.673`$. Using this $`X`$ to fit the sum of the squared $`\eta `$ and $`\eta ^{}`$ masses to the empirical value, yields the column B in Table II, where we see a slight improvement in the $`\eta `$ and $`\eta ^{}`$ masses with respect to the results (68), while the mixing angle is still acceptable, being less than $`2^{}`$ away from the angle favored in Sec. II. If we treat $`X`$ as the second free parameter (this procedure yields the column C of Table II) so that we are able to fit $`m_\eta `$ and $`m_\eta ^{}`$ precisely to their experimental values, we get $`X=0.805`$, along with the mixing angle $`\theta =14.9^{}`$ and the chiral-limit-nonvanishing singlet mass $`3\beta =0.801`$ GeV<sup>2</sup>=(895 MeV)<sup>2</sup>. This is noticeably closer to $`\theta `$ and $`3\beta `$ resulting from other procedures (where $`X`$ is not a free parameter) than before the aforementioned $`\pi ^0K`$ refitting to $`m_\pi =137.3`$ MeV and $`m_K=495.7`$ MeV. Next, we note in the column D of Table II that the slightly improved fit to the masses also led to somewhat improved $`\eta ,\eta ^{}\gamma \gamma `$ widths when we extract from them $`\theta =12.8^{}`$, practically the same as Ref. and the Sec. II result (20). All the three possibilities B, C, and D, do not differ too much from each other, and agree reasonably with the experimental masses and $`\gamma \gamma `$ widths given in column E as well as with the corresponding results of Sec. II. This contrasts with column A, which also contains the results of the new fit but with $`X=1`$. Column A shows that when $`X=1`$, a good description of the masses requires a $`\theta `$ value too negative for a good description of the $`\gamma \gamma `$ widths in the SD-BS approach. Column A thus convinces us that it was precisely the lack of the strangeness attenuation factor $`X`$ that prevented Ref. from satisfactorily reproducing the $`\eta ^{}`$ mass when it successfully did so with the $`\eta `$ mass and $`\gamma \gamma `$ widths. ## V Conclusion In Sec. II we studied the first U<sub>A</sub>(1) problem associated with the Goldstone structure of $`\eta (547)`$ and $`\eta ^{}(958)`$ mesons. Following a QCD gluon-mediated approach to $`\eta `$$`\eta ^{}`$ particle mixing, we began by extracting an $`\eta `$$`\eta ^{}`$ mixing angle $`\varphi 42^{}`$ in the NS–S basis or $`\theta 13^{}`$ in the singlet-octet basis. This led to eta masses $`\eta _8(567)`$, $`\eta _0(947)`$ with chiral-limiting (CL) $`\eta _0(917)`$. Then the physical eta mass $`\eta (547)`$ is 97% of $`\eta _8(567)`$, while $`\eta ^{}(958)`$ is 104% of the CL $`\eta _0(917)`$. Such a 3–4% CL suppression is likewise found for the pion decay constant $`f_\pi 93`$ MeV $`90`$ MeV and for the $`K_{l3}`$ form factor $`f_+(0)=10.96`$$`0.97`$. Then in Sec. III we studied the second U<sub>A</sub>(1) problem associated with eta meson hadronic decay rates. The $`\eta ,\eta ^{}3\pi ^0`$ ($`\mathrm{\Delta }I=1`$) decay rates of 377 eV followed from PCAC Consistency. Also a (strong) decay rate of 130 keV for $`\eta ^{}\eta \pi \pi `$ was obtained from $`\delta `$ and $`\sigma `$ scalar meson poles combined with linear $`\sigma `$ model couplings. These three rates are compatible with data finding $`\mathrm{\Gamma }(\eta 3\pi ^0)=380\pm 36`$ eV, $`\mathrm{\Gamma }(\eta ^{}3\pi ^0)=313\pm 58`$ eV and $`\mathrm{\Gamma }(\eta ^{}\eta \pi \pi )=131\pm 8`$ keV. Finally, in Sec. IV we showed the consistency of the above results with those obtained in a chirally well-behaved quark model which was explicitly constructed through D$`\chi `$SB, SD and BS equations. For example, described variations of our SD-BS approach lead to $`\theta 13^{}\pm 2^{}`$ and to the corresponding CL $`\eta _0`$ mass $`\sqrt{3\beta }=912\pm 18`$ MeV. Successful reproduction of the Abelian axial anomaly amplitudes in the CL in this bound-state approach, gives particular weight to our conclusion that so far away from the CL as in the case of the $`\eta `$$`\eta ^{}`$ complex, $`\gamma \gamma `$-decay constants ($`\overline{f}_{\eta _8},\overline{f}_{\eta _0}`$) differ significantly from the usual axial-current decay constants ($`f_{\eta _8},f_{\eta _0}`$). By allowing for the effects of the SU(3) flavor symmetry breaking also in $`q\overline{q}`$ annihilation graphs, we have improved the $`\eta `$$`\eta ^{}`$ mass matrix with respect to the mass matrix in Ref. (via the strangeness attenuation factor $`X=0.663`$). The consistency of our approach and results with the two-mixing-angle scheme is shown in detail in the Appendix, where we also compute the mixing angles and axial decay constants in that scheme. Acknowledgments: D. Kl. and D. Ke. acknowledge the support of the Croatian Ministry of Science and Technology under the respective contract numbers 1–19–222 and 009802. M. D. S. appreciates discussions with H. F. Jones, R. Delbourgo and V. Elias, and thanks the Univ. of Western Ontario for hospitality. ## Connection with the two–mixing–angles scheme In this appendix, we clarify the relationship of our approach with the two-mixing-angle scheme considered by Leutwyler and Kaiser as well as FKS , and reviewed by Feldmann . The two-mixing-angle scheme is defined with respect to the mixing of the decay constants. As such, it is very suitable in studies where manipulating decay constants is crucial, e.g., when one expresses amplitudes through them with the help of PCAC. In other situations, the mixing of the states may be crucial. The SD–BS variant of our approach, where one explicitly solves for quark-antiquark bound states, and then uses them for direct calculation of amplitudes, is an especially clear example of that. (See Ref. for another recent example.) As FKS themselves state several lines below their Eq. (1.3), the appearance of the four parameters in the two-mixing-angle scheme, namely $`f_8,f_0,\theta _8`$ and $`\theta _0`$, raises anew the problem of their mutual relations and their connection with the mixing angle of the particle states, which is necessarily a single one. In our approach, it is convenient to utilize a mixing scheme defined with respect to a state basis corresponding to the broken SU(3) flavor symmetry, $`\|\eta _{\text{NS}}`$ and $`\|\eta _\text{S}`$ or, equivalently, the effective<sup>\</sup><sup>\</sup>\*Note that in spite of the differences in notation, the effective octet and singlet states (II) of the broken SU(3) flavor, correspond to the effective octet and singlet states $`\eta _8`$ and $`\eta _0`$ given, e.g., by Eq. (85) in Feldmann’s review . SU(3)–broken $`\|\eta _8`$ and $`\|\eta _0`$ \[Eqs. (II)\], i.e., the state mixing angle $`\varphi `$ of Eqs. (11) or (mathematically completely equivalently) the state mixing angle $`\theta `$ of Eqs. (II). Nevertheless, we show below that a) we can calculate quantities utilized in the two-mixing-angle scheme defined with respect to the mixing of the decay constants, and b) what we find for these quantities is close to what is quoted in Refs. . Independently of any specific approach, one can always define quite generally the axial-current decay constants $`f_\eta ^8`$, $`f_\eta ^{}^8`$, $`f_\eta ^0`$, and $`f_\eta ^{}^0`$ as the matrix elements $$0\|A^{a\mu }(x)\|P(p)=if_P^ap^\mu e^{ipx},a=8,0;P=\eta ,\eta ^{}.$$ (69) These definitions are to be contrasted with the somewhat arbitrary definitions of the just two individual axial–current decay constants $`f_\eta `$ and $`f_\eta ^{}`$ often used in this context, where $`f_\eta \mathrm{cos}^2\varphi f_{\text{NS}}+\mathrm{sin}^2\varphi f_\text{S}`$ and $`f_\eta ^{}\mathrm{sin}^2\varphi f_{\text{NS}}+\mathrm{cos}^2\varphi f_\text{S}`$, where $`f_{\text{NS}}`$ and $`f_\text{S}`$ are defined below in Eqs. (Connection with the two–mixing–angles scheme). These $`f_\eta `$ and $`f_\eta ^{}`$ stem from Eqs. (11) in conjunction with the rather arbitrary definitions $`A_\eta ^\mu (x)\mathrm{cos}\varphi A_{\text{NS}}^\mu (x)\mathrm{sin}\varphi A_\text{S}^\mu (x)`$ and $`A_\eta ^{}^\mu (x)\mathrm{sin}\varphi A_{\text{NS}}^\mu (x)+\mathrm{cos}\varphi A_\text{S}^\mu (x)`$, with Eqs. (Connection with the two–mixing–angles scheme) below defining the nonstrange and strange axial currents, $`A_{\text{NS}}^\mu (x)`$ and $`A_\text{S}^\mu (x)`$. In contrast to this, the four decay constants $`f_\eta ^8`$, $`f_\eta ^{}^8`$, $`f_\eta ^0`$, and $`f_\eta ^{}^0`$ are precisely defined by Eqs. (69) and therefore can have unambiguous and process-independent meaning . Following the convention of Leutwyler and Kaiser , the four decay constants $`f_\eta ^8`$, $`f_\eta ^{}^8`$, $`f_\eta ^0`$, and $`f_\eta ^{}^0`$ are parametrized in terms of two decay constants $`f_0`$, $`f_8`$, and two angles $`\theta _0`$, $`\theta _8`$: $`f_\eta ^8`$ $`=`$ $`\mathrm{cos}\theta _8f_8,`$ (71) $`f_\eta ^{}^8`$ $`=`$ $`\mathrm{sin}\theta _8f_8,`$ (72) $`f_\eta ^0`$ $`=`$ $`\mathrm{sin}\theta _0f_0,`$ (73) $`f_\eta ^{}^0`$ $`=`$ $`\mathrm{cos}\theta _0f_0.`$ (74) We define the currents: $`A_{\text{NS}}^\mu (x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}A^{8\mu }(x)+\sqrt{{\displaystyle \frac{2}{3}}}A^{0\mu }(x)={\displaystyle \frac{1}{2}}\left(\overline{u}(x)\gamma ^\mu \gamma _5u(x)+\overline{d}(x)\gamma ^\mu \gamma _5d(x)\right),`$ (76) $`A_\text{S}^\mu (x)`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}A^{8\mu }(x)+{\displaystyle \frac{1}{\sqrt{3}}}A^{0\mu }(x)={\displaystyle \frac{1}{\sqrt{2}}}\overline{s}(x)\gamma ^\mu \gamma _5s(x).`$ (77) The corresponding NS–S decay constants (analogous to the constants $`f_\eta ^8`$, $`f_\eta ^{}^8`$, $`f_\eta ^0`$, and $`f_\eta ^{}^0`$ defined above) are defined as $$0\|A_F^\mu (x)\|P(p)=if_P^Fp^\mu e^{ipx},F=\text{NS},S;P=\eta ,\eta ^{}.$$ (78) The relations (Connection with the two–mixing–angles scheme) between the currents dictate that the relations between these two sets of decay constants are given, exactly and model independently, by $$\left[\begin{array}{cc}f_\eta ^{\text{NS}}& f_\eta ^\text{S}\\ f_\eta ^{}^{\text{NS}}& f_\eta ^{}^\text{S}\end{array}\right]=\left[\begin{array}{cc}f_\eta ^8& f_\eta ^0\\ f_\eta ^{}^8& f_\eta ^{}^0\end{array}\right]\left[\begin{array}{cc}\frac{1}{\sqrt{3}}& \sqrt{\frac{2}{3}}\\ \sqrt{\frac{2}{3}}& \frac{1}{\sqrt{3}}\end{array}\right],$$ (79) where we used matrix notation for compactness. If we have well–defined nonstrange–strange states $`\|\eta _{\text{NS}}`$ and $`\|\eta _\text{S}`$ (II) \[as in our SD–BS approach\] we can define the decay constants $`f_{\text{NS}}`$ and $`f_\text{S}`$ through $`0\|A_{\text{NS}}^\mu (x)\|\eta _{\text{NS}}(p)`$ $`=`$ $`if_{\text{NS}}p^\mu e^{ipx},`$ (81) $`0\|A_\text{S}^\mu (x)\|\eta _\text{S}(p)`$ $`=`$ $`if_\text{S}p^\mu e^{ipx},`$ (82) $`0\|A_{\text{NS}}^\mu (x)\|\eta _\text{S}(p)`$ $`=`$ $`0,`$ (83) $`0\|A_\text{S}^\mu (x)\|\eta _{\text{NS}}(p)`$ $`=`$ $`0.`$ (84) Since the states $`\|\eta `$ and $`\|\eta ^{}`$ are given by Eqs. (11) as the linear combinations of $`\|\eta _{\text{NS}}`$ and $`\|\eta _\text{S}`$, we can relate the constants $`\{f_\eta ^{\text{NS}},f_\eta ^{}^{\text{NS}},f_\eta ^\text{S},f_\eta ^{}^\text{S}\}`$ with $`\{f_{\text{NS}},f_\text{S}\}`$: $$\left[\begin{array}{cc}f_\eta ^{\text{NS}}& f_\eta ^\text{S}\\ f_\eta ^{}^{\text{NS}}& f_\eta ^{}^\text{S}\end{array}\right]=\left[\begin{array}{cc}\mathrm{cos}\varphi & \hfill \mathrm{sin}\varphi \\ \mathrm{sin}\varphi & \hfill \mathrm{cos}\varphi \end{array}\right]\left[\begin{array}{cc}f_{\text{NS}}& \hfill 0\\ 0& \hfill f_\text{S}\end{array}\right].$$ (85) Using Eq. (79), we can relate the decay constants $`\{f_\eta ^8,f_\eta ^{}^8,f_\eta ^0,f_\eta ^{}^0\}`$ with $`\{f_{\text{NS}},f_\text{S}\}`$: $$\left[\begin{array}{cc}f_\eta ^8& f_\eta ^0\\ f_\eta ^{}^8& f_\eta ^{}^0\end{array}\right]=\left[\begin{array}{cc}\mathrm{cos}\varphi & \hfill \mathrm{sin}\varphi \\ \mathrm{sin}\varphi & \hfill \mathrm{cos}\varphi \end{array}\right]\left[\begin{array}{cc}f_{\text{NS}}& \hfill 0\\ 0& \hfill f_\text{S}\end{array}\right]\left[\begin{array}{cc}\frac{1}{\sqrt{3}}& \sqrt{\frac{2}{3}}\\ \sqrt{\frac{2}{3}}& \frac{1}{\sqrt{3}}\end{array}\right],$$ (86) which is the same<sup>††</sup><sup>††</sup>††Note that the last matrix in our Eq. (86) is just $`U^{}(\theta _{ideal})`$ in the notation of Ref. , where $`\theta _{ideal}\mathrm{arctan}\sqrt{2}`$. as in Feldmann et al. , who use the notation $`f_q=f_{\text{NS}},f_s=f_\text{S}`$. The above equations together with the definitions (Connection with the two–mixing–angles scheme) give the following solutions for $`f_8`$, $`f_0`$, $`\theta _8`$, and $`\theta _0`$: $`f_8`$ $`=`$ $`\sqrt{{\displaystyle \frac{1}{3}}f_{\text{NS}}^2+{\displaystyle \frac{2}{3}}f_\text{S}^2},`$ (88) $`\theta _8`$ $`=`$ $`\varphi \text{arctan}\left({\displaystyle \frac{\sqrt{2}f_\text{S}}{f_{\text{NS}}}}\right),`$ (89) $`f_0`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}f_{\text{NS}}^2+{\displaystyle \frac{1}{3}}f_\text{S}^2},`$ (90) $`\theta _0`$ $`=`$ $`\varphi \text{arctan}\left({\displaystyle \frac{\sqrt{2}f_{\text{NS}}}{f_\text{S}}}\right).`$ (91) These relations (the same as already found by FKS ) are obtained in a general way, independently of specifics of any given approach; therefore, we can also apply them within our framework. In Sec. II, we have already pointed out the agreement of our results for our preferred state mixing angle $`\varphi 42^{}`$ with the FKS results quoted in Refs. , but the values we find for $`f_0`$, $`f_8`$, $`\theta _0`$, and $`\theta _8`$ are also similar to theirs. The important quantity here is $`y=f_{\text{NS}}/f_\text{S}`$, giving the extent of the SU(3) breaking. In our approach, this quantity is also present and plays a crucial role. The information it carries enables us to calculate what the decay-constant-mixing angles $`\theta _8`$ and $`\theta _0`$ would be in our approach (although we stress again that because of the way we calculate, their “rough average” $`\theta `$ must retain the central role because it has the meaning of the state-mixing angle). This is all because the $`y`$-ratio is essentially (as easily seen through the Goldberger–Treiman relation for constituent quarks) our parameter $`X`$, which can be (at least approximately) expressed as the ratio of the nonstrange-to-strange constituent quark mass: $`X\widehat{m}/m_s`$. In addition, in our very predictive coupled SD-BS approach, we can directly calculate all decay constants including $`f_q=f_{\text{NS}}`$ and $`f_s=f_S`$, and this again gives (in a good approximation) the same value for $`X=y=f_{\text{NS}}/f_S`$. While we have $`f_{\text{NS}}=f_\pi `$ (also assumed by FKS , in their theoretical analysis), our calculation yields $`f_S=1.4505f_\pi `$, less than 3% more than the theoretical FKS prediction . (Note that we used the symbol $`f_{s\overline{s}}`$ for $`f_s=f_S`$ .) Our chosen model therefore gives $`y=f_{\text{NS}}/f_S=0.6894`$, leading to $`f_8=1.318f_\pi `$ and $`f_0=1.170f_\pi `$. (Interestingly, this is practically equal to our SD-BS model values $`f_{\eta _8}=1.31f_\pi `$ and $`f_{\eta _0}=1.16f_\pi `$ for the octet and singlet axial-current decay constants $`f_{\eta _8}`$ and $`f_{\eta _0}`$, mentioned in Sec. IV. They are defined in the standard way through the matrix elements $`0\|A^{a\mu }\|\eta _a`$, $`(a=8,0)`$, so that the definitions (II) and (Connection with the two–mixing–angles scheme) imply that $`f_{\eta _8}`$ and $`f_{\eta _0}`$ are straightforwardly expressed through $`f_{\text{NS}}`$ and $`f_\text{S}`$ by the relations $`f_{\eta _8}=\frac{1}{3}f_{\text{NS}}+\frac{2}{3}f_\text{S}`$ and $`f_{\eta _0}=\frac{2}{3}f_{\text{NS}}+\frac{1}{3}f_\text{S}`$. We thus note that the quadratic relations (88) and (90) for differently defined octet and singlet constants $`f_8`$ and $`f_0`$, lead to similar values as the linear relations for $`f_{\eta _8}`$ and $`f_{\eta _0}`$.) Using in Eqs. (Connection with the two–mixing–angles scheme) our preferred state mixing angle $`\varphi =42^{}`$, our model value $`y=f_{\text{NS}}/f_S=0.6896`$ also leads to the following decay-constant-mixing angles in the $`\eta _8\eta _0`$ basis: $`\theta _8=22^{}`$ and $`\theta _0=2.3^{}`$, close to the theoretical FKS results . See also Table 1 in Ref. , line “FKS scheme & theory”, giving the $`\theta _8=21.0^{}`$ and $`\theta _0=2.7^{}`$, while the line “FKS scheme & phenomenology” in the same table has only somewhat more negative $`\theta _0`$ but larger $`f_0/f_\pi `$. The “FKS scheme & theory” then implies the state-mixing angle $`\theta 12^{}`$, in agreement with our results. This is as it should be, as we note that the mass matrix in Ref. , its Eq. (73), coincides with ours when the anomaly contribution $`a^2`$ is identified with the “Veneziano term” $`\lambda _\eta ^2/3(=\beta )`$, as we do. Recalling that we have already pointed out the agreement of the NS–S mixing angles obtained by us and by FKS , one can see that everything tallies. ## Figure captions Nonperturbative QCD quark annihilation illustrated by the diagram with two-gluon exchange. It shows the transition of the $`f\overline{f}`$ pseudoscalar $`P`$ into the pseudoscalar $`P^{}`$ having the flavor content $`f^{}\overline{f}^{}`$. The dashed lines and full circles depict the $`q\overline{q}`$ bound-state pseudoscalars and vertices, respectively.
warning/0003/math0003188.html	ar5iv	text	# Krichever correspondence for algebraic varieties ## 1 Introduction In 70’s years I. M. Krichever suggested a construction how to attach to some algebraic-geometric data, connected with algebraic curves and vector bundles on them, an infinite-dimensional (Fredholm) subspace in the space $`k((z))`$ of Laurent power series ( ). This construction was successfully used in the theory of integrable systems, in paricular, in the theory of KP and KdV equations (, , ). There were also found applications of this construction to the theory of modules of algebraic curves (, ). Besides, this construction turned out to be connected with description of commutative subrings in the rings of pseudo-differential operators (, ). Now this construction is called the Krichever correspondence or the Krichever map (, , , ). But in these works it is essentially that algebraic-geometric data are connected with $`1`$-dimensional varieties and $`1`$-dimensional local field $`k((z))`$. Recently, in works , it were pointed out some connections between the theory of the KP-equations and $`n`$-dimensional local fields; also it was suggested a variant of the Krichever map for algebraic-geometric data which is connected with algebraic surfaces, vector bundles on them and $`2`$-dimensional local fields. One of the typical examples of multidimensional local field is the field of Laurent iterated series $`k((z_1))\mathrm{}((z_n))`$. Such fields serve for natural generalization of local objects of $`1`$-dimensional varieties to the case of multidimensional varieties. Let us consider an $`n`$-dimensional algebraic scheme $`X`$. Let $`Y_0\mathrm{}Y_n`$ be a flag of closed subschemes on $`X`$ such that $`Y_0=X`$, $`Y_i`$ is of codimension $`1`$ in $`Y_{i1}`$, and $`Y_n=x`$ is a closed point. Then there exists a construction (, , ), attaching in the canonical way to such flag some ring, which is an $`n`$-dimensional local field provided that $`x`$ is a smooth point on all $`Y_i`$. Moreover, if $`X`$ is an algebraic variety over a field $`k`$, $`x`$ is a $`k`$-rational point, and we fix local parameters $`z_1,z_2,\mathrm{},z_n\widehat{𝒪}_{x,X}`$ such that $`z_{ni+1}=0`$ is a local equation of variety $`Y_i`$ in the formal neighbourhood of the point $`x`$ on the variety $`Y_{i1}`$ ($`1in`$), then the obtained $`n`$-dimensional local field it is possible to identify with $`k((z_1))((z_2))\mathrm{}((z_n))`$. A concept of multidimensional local field has appeared in the middle of 70’s years, and, originally, such fields were used for the development of generalization of class field theory to the schemes of higher dimension. Later there were also found applications of multidimensional local fields to many problems of algebraic geometry, where it makes sence to speak about local components of geometric objects (see ). In this work, using multidimensional local fields, we construct the Krichever correspondence for varieties of arbitrary dimension $`n`$: that is some injective map from algebraic-geometric data, connected with projective algebraic varieties, full flags of ample divisors and their local parameters in the formal neighbourhood of the last point of the flag, vector bundles and their trivialisations in the formal neighbourhood of the last point of the flag, to some $`k`$-subspaces of finite dimensional vector space over the $`n`$-dimensional local field $`k((z_1))\mathrm{}((z_n))`$. If $`n=1`$, then our constructed map is a variant of the Krichever map for curves. If $`n=2`$, then our constructed map coincides with the map constructed in . The work is organized as follows. In §2 we give various technical lemmas about cohomology of coherent sheaves, projective and injective limits, which will be useful further in the work. In §3 we give a construction of family of functors, which is connected with quasicoherent sheaves and a fixed flag of subvarieties, and which can be interpreted as a cohomology system of coeficients on the standart symplex. In §4, using the construction of §3, we construct complexes of sheaves of abelian groups, which is acyclic resolutions of arbitrary quasicoherent sheaves on schemes. In §5 we prove some theorems about intersections among components of resolutions constructed in §4. In some cases the whole resolution can be reconstructed from one $`k`$-subspace of finite dimensional vector space over $`k((z_1))\mathrm{}((z_n))`$. Using this, we construct the Krichever map in higher dimensions. Note that connected with $`n`$-dimensional local fileds resolutions of quasicoherent sheaves on schemes were in works , . But in contrast to these works, our resolutions depend only on a single flag of subvarieties and are not resolutions of adelic type. Note also that, in contrast to , all the constructions and proofs in this work are internal ones, i. e., they are not reduced to multidimensional adelic complexes. During all the work we shall keep the following notations and agreements. For any finite set $`I`$ let $`\mathrm{}I`$ be the number of elements of the set $`I`$. If $`X`$ is a scheme, then $`Sh(X)`$ is the category of sheaves of abelian groups on $`X`$, $`CS(X)`$ is the category of coherent sheaves on $`X`$, $`QS(X)`$ is the catehory of quasicoherent sheaves on $`X`$, $`Ab`$ is the catehory of abelian groups. If $`f:YX`$ is a morphism of two schemes, then always $`f^{}`$ is the pull-back functor in the category of sheaves of abelian groups, $`f_{}`$ is the direct image functor in the category of sheaves of abelian groups. If $`𝒰`$ is an open covering of $`X`$, $``$ is a sheaf of abelian groups on $`X`$, then $`\stackrel{ˇ}{H}^{}(𝒰,)`$ are the Čech cohomologies groups with respect to the covering $`𝒰`$. Let $`YX`$ be a closed subscheme of a scheme $`X`$, which is defined by the ideal sheaf $`J`$. Then by $`(Y,𝒪_X/J^k)`$ denote the scheme whose topological space coincides with the topological space of the scheme $`Y`$ and the structure sheaf is $`𝒪_X/J^k`$. ($`𝒪_X`$ is the structure sheaf of the scheme $`X`$.) The author would like to express the deep gratitude to his scientific adviser A. N. Parshin for the constant attention to the work. ## 2 Technical lemmas ###### Lemma 1 On a noetherian scheme $`X`$ any short exact sequence of quasicoherent sheaves is direct limit of short exact sequences of coherent sheaves. For $`\varphi :𝒢\mathrm{Mor}(QS(X))`$ there are $`\varphi _i:_i𝒢_i\mathrm{Mor}(CS(X))`$ with $`\underset{}{lim}\varphi _i=\varphi `$, $`_i`$, $`𝒢_i𝒢`$. Proof. See \[11, lemma 1.2.2\] and \[10, lemma 2.1.5\]. ###### Lemma 2 Let $`X`$ be a noetherian scheme. Let $`\psi :CS(X)Sh(X)`$ be an exact additive functor. Then $`\psi `$ commutes with direct limits. Proof. (By analogy with \[11, lemma 1.2.3\] or \[10, lemma 2.2.2\].) First let us prove that if we have a direct system of sheaves $$\{_i:iI,\varphi _{ij}:_i_j(ij)\}$$ with $`\underset{}{lim}_i=0`$, then $`\underset{}{lim}\psi (_i)=0`$. For this one we prove that for any open $`UX`$ $`\underset{}{lim}H^0(U,\psi (_i))=0`$. Let $`x\underset{}{lim}H^0(U,\psi (_i))`$. Let this $`x`$ be represented by $`x_iH^0(U,\psi (_i))`$. Then from coherent property of the sheaf $`_i`$ and noetherian property of the scheme $`X`$ there is some $`jI`$ such that $`\varphi _{ij}=0`$. Since $`\psi `$ is an additive functor, we have that $`\psi (\varphi _{ij})=0`$. Therefore $$\begin{array}{ccc}\hfill {}_{}{}^{0}(U,\psi (_i))& & H^0(U,\psi (\varphi _{ij})(\psi (_i)))\hfill \\ \hfill x_i& & 0\hfill \end{array}$$ Now consider the general case: let $`\underset{}{lim}_i=`$, $`\varphi _i:_i`$ be the canonical morphisms. Consider the following exact sequence of coherent sheaves: $$0\mathrm{Ker}\varphi _i_i\mathrm{Coker}\varphi _i0$$ The functor $`\psi `$ is an exact functor, therefore we have the following exact sequence: $$0\psi (\mathrm{Ker}\varphi _i)\psi (_i)\psi ()\psi (\mathrm{Coker}\varphi _i)0$$ From $`\underset{}{lim}\mathrm{Ker}\varphi _i=0`$ and $`\underset{}{lim}\mathrm{Coker}\varphi _i=0`$ it follows by arguments above that $`\underset{}{lim}\psi (\mathrm{Ker}\varphi _i)=0`$ and $`\underset{}{lim}\psi (\mathrm{Coker}\varphi _i)=0`$. Direct limit maps exact sequences to exact sequence. Therefore $`\underset{}{lim}\psi (_i)=\psi ()`$. Lemma 2 is proved. ###### Lemma 3 Let $`X`$ be a noetherian scheme. Then an exact additive functor $`\psi :CS(X)Sh(X)`$ can be uniquely extended to a functor $`\psi ^{}:QS(X)Sh(X)`$ which commutes with direct limits. This new functor is exact as well. Proof. (By analogy with \[11, lemma 1.2.4\].) Let $`\mathrm{Ob}(QS(X))`$. By lemma 1 $`=\underset{}{lim}_i`$, where $`_i\mathrm{Ob}(CS(X))`$. Define $$\psi ^{}()=\underset{}{lim}\psi (_i)\text{.}$$ We have $`\psi ^{}()=\psi ()`$ for $`Ob(CS(X))`$ by lemma 2. By lemma 1, for any $`\varphi \mathrm{Mor}(QS(X))`$ we have $`\varphi =\underset{}{lim}\varphi _i`$, where $`\varphi _i\mathrm{Mor}(CS(X))`$. Define $$\psi ^{}(\varphi )=\underset{}{lim}\psi (\varphi _i)\text{.}$$ By lemma 2, we have that $`\psi ^{}(\varphi )=\psi (\varphi )`$ for $`\varphi \mathrm{Mor}(CS(X))`$. It is clear that this definition is the only one possible. And by lemma 2, it is well defined. Lemma 3 is proved. ###### Lemma 4 Let $`X`$ be a noetherian scheme, $`i:YX`$ be a closed subscheme, which is defined by the ideal sheaf $`J`$ on $`X`$. Let $`j:UY`$ be an open subscheme of $`Y`$ such that for any point $`xX`$ there exists an affine neighborhood $`Vx`$ such that $`VU`$ is an affine subscheme. Let the supports of sheaves $`_iSh(X)`$ ($`i=1,\mathrm{},3`$) are in $`Y`$, and the sheaf $`_1`$ is a quasicoherent sheaf with respect to the subscheme $`(Y,𝒪_X/J^k)`$ for some $`kN`$. Then from exactness of the sequence of sheaves $$0_1_2_30$$ (1) it follows exactness of the following sequence $$0i_{}j_{}j^{}_1i_{}j_{}j^{}_2i_{}j_{}j^{}_30\text{.}$$ Proof. First note that for any affine open subscheme $`WU`$ and for any quasicoherent sheaf $`𝒢`$ on the scheme $`(U,𝒪_X/J^k_U)`$ we have $$H^1(W,𝒢)=0\text{.}$$ (2) In fact, if the sheaf $`𝒢`$ is a quasicoherent sheaf with respect to the subscheme $`U=(U,𝒪_X/J_U)`$, then equality (2) follows from affineness of the scheme $`W`$. Now if $``$ is a quasicoherent sheaf with respect to the subscheme $`(U,𝒪_X/J^k_U)`$, $`k𝐍`$, $`k1`$, then consider the following exact sequence: $$0J/J0\text{.}$$ (3) But the sheaves $`J`$ and $`/J`$ are quasicoherent sheaves with respect to the subscheme $`(U,𝒪_X/J^{k1}_U)`$. Therefore we can do induction, from which it follows that $$H^1(W,J)=0\text{and}H^1(W,/J)=0\text{.}$$ Hence and from the long cohomological sequence associated with sequence (3) we obtain equality (2). Return to sequence (1). We have exactness of the following sequence: $$0j^{}_1j^{}_2j^{}_30\text{.}$$ Appliing the functor $`j_{}`$, we obtain $$0j_{}j^{}_1j_{}j^{}_2j_{}j^{}_3R^1j_{}(j^{}_1)$$ Let us show that the sheaf $$R^1j_{}(j^{}_1)=0\text{.}$$ The sheaf $`j^{}_1`$ is a quasicoherent sheaf with respect to the subscheme $`(U,𝒪_X/J^k_U)`$ for some $`k𝐍`$, therefore the sheaf $`R^1j_{}(j^{}_1)`$ is a quasicoherent sheaf on the scheme $`(Y,𝒪_X/J^k)`$ with respect to the same $`k𝐍`$. Therefore it suffices to show that for affine open $`V`$ from the lemma’s conditions $$H^0(VY,R^1j_{}(j^{}_1))=0\text{.}$$ (4) But $`H^0(VY,R^1j_{}(j^{}_1))=H^1(VU,j^{}_1)`$. And equality (4) follows from equality (2). Therefore we have exactness of the following sequence: $$0j_{}j^{}_1j_{}j^{}_2j_{}j^{}_30\text{.}$$ From $`i:YX`$ is a closed imbedding it follows that $`i_{}`$ is an exact functor. Therefore the following sequence is exact: $$0i_{}j_{}j^{}_1i_{}j_{}j^{}_2i_{}j_{}j^{}_30\text{.}$$ Lemma 4 is proved. Let $`X`$ be a noetherian scheme. Suppose that we have an exact and additive functor $`\mathrm{\Phi }:QS(X)Sh(X)`$. Let $`i:YX`$ be a closed subscheme of the scheme $`X`$, which is defined by the ideal sheaf $`J`$. Let $`j:UY`$ be an open subscheme of $`Y`$. Then define a functor $$C_U\mathrm{\Phi }:CS(X)Sh(X)\text{as following:}$$ for any sheaf $`CS(X)`$: $$C_U\mathrm{\Phi }()\stackrel{\mathrm{def}}{=}\underset{k𝐍}{\underset{}{lim}}\mathrm{\Phi }(i_{}j_{}j^{}(/J^k))\text{.}$$ ###### Remark 1 It is not difficult to understand that if the sheaf $``$ is a coherent sheaf on $`X`$, then for any $`k𝐍`$ the sheaf $`i_{}j_{}j^{}(/J^k)`$ is a quasicoherent sheaf on $`X`$. In fact, the sheaf $`/J^k`$ is a coherent sheaf on the scheme $`X`$. Moreover, the sheaf $`/J^k`$ is a coherent sheaf on the closed subscheme $`(Y,𝒪_X/J^k)`$ of the scheme $`X`$. Then $`j^{}(/J^k)`$ is a coherent sheaf on the scheme $`(U,𝒪_X/J^k_U)`$, $`j_{}j^{}(/J^k)`$ is a quasicoherent sheaf on the scheme $`(Y,𝒪_X/J^k)`$. And since $`i_{}`$ coincides with the direct image functor from the subscheme $`(Y,𝒪_X/J^k)`$, we see that $`i_{}j_{}j^{}(/J^k)`$ is a quasicoherent sheaf on $`X`$. ###### Lemma 5 Let $`X`$ be a noetherian scheme, $`\mathrm{\Phi }:QS(X)Sh(X)`$ be an exact additive functor, $`i:YX`$ be a closed subscheme of the scheme $`X`$, which is defined by the ideal sheaf $`J`$ on $`X`$. Let $`j:UY`$ be an open subscheme of $`Y`$. In addition, suppose the following: for any point $`xX`$, for any open $`WX`$, $`xW`$, there exists an affine open subscheme $`VW`$, $`xV`$ such that: 1. $`VU`$ is an affine subscheme; 2. for any quasicoherent sheaf $``$ on $`X`$ $$H^1(V,\mathrm{\Phi }())=0$$ (5) Then $`C_U\mathrm{\Phi }:CS(X)Sh(X)`$ is an exact and additive functor. ###### Remark 2 For example, condition 1 of lemma 5 is satisfied in the following cases: * $`X`$ is a separated scheme, $`U`$ is an affine subscheme (as on a separated scheme the intersection of two affine open subschemes is an affine subscheme); * $`X`$ is a separated scheme, and $`U`$ is a complement to some Cartier divisor in $`Y`$. Proof (of lemma 5). Additivity of the functor $`C_U\mathrm{\Phi }`$ is obvious from the construction. Let us show exactness. Let $$0_1_2_30$$ be an exact sequence of coherent sheaves on $`X`$. For any $`xX`$, for any open $`WX`$ consider an open affine $`VW`$, $`Vx`$ satisfying conditions 1 2 of lemma 5. Then for the proof of exactness of the functor $`C_U\mathrm{\Phi }`$ it suffices to show exactness of the following sequence: $$0H^0(V,C_U\mathrm{\Phi }(_1))H^0(V,C_U\mathrm{\Phi }(_2))H^0(V,C_U\mathrm{\Phi }(_3))0\text{.}$$ (6) On the other hand, the following sequence is exact: $$0_1/J^k_2_1_2/J^k_2_3/J^k_30\text{.}$$ Since the sheaves in the last sequence are coherent sheaves on the scheme $`(Y,𝒪_X/J^k)`$, by lemma 4 the following sequence is exact: $$0i_{}j_{}j^{}(_1/J^k_2_1)i_{}j_{}j^{}(_2/J^k_2)i_{}j_{}j^{}(_3/J^k_3)0\text{.}$$ Since we apply the exact functor $`\mathrm{\Phi }`$ to the last sequence, we obtain exactness of the following sequence from $`Sh(X)`$: $$0\mathrm{\Phi }(i_{}j_{}j^{}(_1/J^k_2_1))\mathrm{\Phi }(i_{}j_{}j^{}(_2/J^k_2))\mathrm{\Phi }(i_{}j_{}j^{}(_3/J^k_3))0\text{.}$$ Now from (5) we obtain exactness of the following sequence: $`0H^0(V,\mathrm{\Phi }(i_{}j_{}j^{}(_1/J^k_2_1)))H^0(V,\mathrm{\Phi }(i_{}j_{}j^{}(_2/J^k_2)))`$ $$H^0(V,\mathrm{\Phi }(i_{}j_{}j^{}(_3/J^k_3)))0\text{.}$$ (7) Also for any natural numbers $`k_1k_2`$ we have the following exact sequence of sheaves: $$0J^{k_1}_2_1/J^{k_2}_2_1_1/J^{k_2}_2_1_1/J^{k_1}_2_10\text{.}$$ Hence, as above, the sheaf $`J^{k_1}_2_1/J^{k_2}_2_1`$ is a coherent sheaf on the scheme $`(Y,𝒪_X/J^{k_2})`$. Therefore by lemma 4 the following sequence is exact: $`0i_{}j_{}j^{}(J^{k_1}_2_1/J^{k_2}_2_1)i_{}j_{}j^{}(_1/J^{k_2}_2_1)`$ $`i_{}j_{}j^{}(_1/J^{k_1}_2_1)0\text{.}`$ Since the functor $`\mathrm{\Phi }`$ is exact, we have exactness of the sequence: $`0\mathrm{\Phi }(i_{}j_{}j^{}(J^{k_1}_2_1/J^{k_2}_2_1))\mathrm{\Phi }(i_{}j_{}j^{}(_1/J^{k_2}_2_1))`$ $`\mathrm{\Phi }(i_{}j_{}j^{}(_1/J^{k_1}_2_1))0\text{.}`$ And from (5) we obtain that the following map is a surjective map: $$\mathrm{\Phi }(i_{}j_{}j^{}(_1/J^{k_2}_2_1))\mathrm{\Phi }(i_{}j_{}j^{}(_1/J^{k_1}_2_1))$$ (8) Now taking the projective limit with respect to all $`k𝐍`$ and using (8), from which it follows the Mittag-Leffler condition (see \[9, ch.II, §9\]), we obtain exactness of the following sequence: $`0\underset{k𝐍}{\underset{}{lim}}H^0(V,\mathrm{\Phi }(i_{}j_{}j^{}(_1/J^k_2_1)))\underset{k𝐍}{\underset{}{lim}}H^0(V,\mathrm{\Phi }(i_{}j_{}j^{}(_2/J^k_2)))`$ $`\underset{k𝐍}{\underset{}{lim}}H^0(V,\mathrm{\Phi }(i_{}j_{}j^{}(_3/J^k_3)))0`$ or $`0\underset{k𝐍}{\underset{}{lim}}H^0(V,\mathrm{\Phi }(i_{}j_{}j^{}(_1/J^k_2_1)))`$ $`H^0(V,C_U\mathrm{\Phi }(_2))H^0(V,C_U\mathrm{\Phi }(_3))0\text{.}`$ Thus for the proof of exactness of sequence (6) we have to show that $$H^0(V,C_U\mathrm{\Phi }(_1))=\underset{k𝐍}{\underset{}{lim}}H^0(V,\mathrm{\Phi }(i_{}j_{}j^{}(_1/J^k_2_1)))\text{.}$$ (9) Since the scheme $`X`$ is a noetherian scheme, by Artin-Rees lemma (see \[2, cor. 10.10\]) there exists $`l𝐍`$ such that for all $`kl`$: $$J^k_2_1=J^{kl}(J^l_2_1)\text{.}$$ (10) From (10) we obtain that the maps $$_1/J^k_2_1_1/J^{kl}_1$$ are well defined and surjective. The same is for $$_1/J^k_1_1/J^k_2_1\text{.}$$ Further, appliing successively the functors $`j^{}`$, $`j_{}`$, $`i_{}`$, $`\mathrm{\Phi }`$ and $`H^0(V,)`$ and using lemma 4 and condition (5), we obtain cofinality of the projective systems: $$H^0(V,\mathrm{\Phi }(i_{}j_{}j^{}(_1/J^k_2_1)))\text{and}H^0(V,\mathrm{\Phi }(i_{}j_{}j^{}(_1/J^k_1)))\text{.}$$ Therefore (9) is satisfied, and lemma 5 is proved. ###### Lemma 6 Let $`X`$ be a noetherian scheme, $`\mathrm{\Phi }:QS(X)Sh(X)`$ be an exact additive functor, $`i:YX`$ be a closed subscheme of the scheme $`X`$, which is defined by the ideal sheaf $`J`$ on $`X`$. Let $`j:UY`$ be an open subscheme of $`Y`$ such that for any point $`xX`$ there exists an affine neighbourhood $`Vx`$ such that $`VU`$ is an affine subscheme. Let $`𝒰=\{U_i\}_{iI}`$ be an affine open covering of the scheme $`X`$. Suppose that for any $`k1`$, for any coherent sheaf $`𝒢`$ on the scheme $`X`$ we have $`\stackrel{ˇ}{H}^m(𝒰,\mathrm{\Phi }(i_{}j_{}j^{}(𝒢/J^k𝒢)))`$ $`=`$ $`0\text{for any}m1`$ (11) $`H^1({\displaystyle \underset{iI_0}{}}U_i,\mathrm{\Phi }(i_{}j_{}j^{}(𝒢/J^k𝒢)))`$ $`=`$ $`0\text{for any subset}I_0I`$ (12) $`H^1(X,\mathrm{\Phi }(i_{}j_{}j^{}(𝒢/J^k𝒢)))`$ $`=`$ $`0\text{.}`$ (13) Then $`\stackrel{ˇ}{H}^m(𝒰,C_U\mathrm{\Phi }())=0`$ for any $`m1`$ and for any coherent sheaf $``$ on $`X`$. Proof. Let $``$ be any coherent sheaf on $`X`$. Let $$p_n:\underset{\mathrm{}I_0=n+1}{\underset{I_0I}{}}H^0(\underset{iI_0}{}U_i,\mathrm{\Phi }(i_{}j_{}j^{}(/J^k)))\underset{\mathrm{}I_0=n+2}{\underset{I_0I}{}}H^0(\underset{iI_0}{}U_i,\mathrm{\Phi }(i_{}j_{}j^{}(/J^k)))$$ be the map which is arised from the Čech complex with respect to the covering $`𝒰`$. Define $`H_k^n\stackrel{\mathrm{def}}{=}\mathrm{Ker}p_n`$. In addition, $$H_k^0=H^0(X,\mathrm{\Phi }(i_{}j_{}j^{}(/J^k)))\text{.}$$ From (11) we obtain at once that for any $`n1`$ $$H_k^n=\mathrm{Im}p_{n1}\text{.}$$ Therefore for any $`m0`$ the following sequence is exact: $$0H_k^n\underset{\mathrm{}I_0=n+1}{\underset{I_0I}{}}H^0(\underset{iI_0}{}U_i,\mathrm{\Phi }(i_{}j_{}j^{}(/J^k)))\stackrel{p_n}{}H_k^{n+1}0$$ (14) For any natural numbers $`k_1k_2`$ we have the exact sequence $$0J^{k_1}/J^{k_2}/J^{k_2}/J^{k_1}0\text{.}$$ Since the sheaves of this sequence are coherent sheaves on the scheme $`(Y,𝒪_X/J^{k_2})`$, by lemma 4 the following sequence is exact: $$0i_{}j_{}j^{}(J^{k_1}/J^{k_2})i_{}j_{}j^{}(/J^{k_2})i_{}j_{}j^{}(/J^{k_1})0\text{.}$$ Further, from exactness of the functor $`\mathrm{\Phi }`$ and condition (13) we obtain surjectivity of the following maps for any natural numbers $`k_1k_2`$: $$H^0(X,\mathrm{\Phi }(i_{}j_{}j^{}(/J^{k_2})))H^0(X,\mathrm{\Phi }(i_{}j_{}j^{}(/J^{k_1})))\text{.}$$ (15) By condition (12), we obtain as well that for any $`k_1k_2`$ and any $`I_0I`$ the maps $$H^0(\underset{iI_0}{}U_i,\mathrm{\Phi }(i_{}j_{}j^{}(/J^{k_2})))H^0(\underset{iI_0}{}U_i,\mathrm{\Phi }(i_{}j_{}j^{}(/J^{k_1})))\text{.}$$ (16) are surjective maps. Let us prove that for any $`n0`$ the map $$H_{k_2}^nH_{k_1}^n$$ (17) is a surjective map for any $`k_1k_2`$. In the case $`n=0`$ it is statement (15). For arbitrary $`n`$ it follows from surjectivity of $`p_{n1}`$ in exact sequence (14) and, also, surjectivity of map (16). Now taking the projective limit in (14) and using surjectivity of (17), from which it follows the Mittag-Leffler condition for the projective systems $`{}_{k}{}^{}{}_{}{}^{n}`$, we obtain exactness of the following sequence for any $`n0`$: $$0\underset{k}{\underset{}{lim}}H_k^n\underset{\mathrm{}I_0=n+1}{\underset{I_0I}{}}H^0(\underset{iI_0}{}U_i,C_U\mathrm{\Phi }())\underset{k}{\underset{}{lim}}H_k^{n+1}0\text{.}$$ Hence it follows at once that for any $`m1`$ $`\stackrel{ˇ}{H}^m(𝒰,C_U\mathrm{\Phi }())=0`$. Lemma 6 is proved. In the sequel we’ll use the following variant of A. Kartan lemma, which connects the Čech cohomologies groups with the usual cohomologies groups of sheaves. For a sheaf $`𝒜`$ on a topological space $`V`$ by $`\stackrel{ˇ}{H}^q(V,𝒜)`$ denote direct limit of Čech cohomologies with respect to all coverings of the space $`V`$. ###### Lemma 7 Let $`X`$ be a topological space and $`𝒜`$ be a sheaf on . Suppoce that it is possible to cover $`X`$ by family $`𝐔`$ of open sets such that this family has the following properties: 1. If $`𝐔`$ contains $`U^{}`$ and $`U^{\prime \prime }`$, then it contains $`U^{}U^{\prime \prime }`$; 2. $`𝐔`$ contains arbitrarily small open sets; 3. $`\stackrel{ˇ}{H}^q(U,𝒜)=0`$ for any $`q1`$ and $`U𝐔`$. Under these conditions we have isomorphism: $$\stackrel{ˇ}{H}^q(X,𝒜)H^q(X,𝒜)\text{.}$$ Proof. See theorem 5.9.2 of \[8, ch.2\]. ###### Lemma 8 Let $`X`$ be a noetherian scheme, $`j:UX`$ be an open affine subscheme. Then for any quasicoherent sheaf $``$ on $`U`$: $`H^0(X,j_{})`$ $`=`$ $`H^0(U,)`$ (18) $`H^i(X,j_{})`$ $`=`$ $`0\text{if}i>0\text{.}`$ (19) Proof. Equality (18) follows from the construction of the functor $`j_{}`$. Let us prove (19). Embed the quasicoherent sheaf $``$ in a flasque quasicoherent sheaf $`𝒢`$ on $`U`$. (It always can do, see \[9, ch. III, §3\].) $$0𝒢𝒢/0$$ Now the following sequence is exact: $$0j_{}j_{}𝒢j_{}(𝒢/)0\text{.}$$ (20) (Indeed, $`R^1j_{}=0`$. The last follows from quasicoherentness of the sheaf $`R^1j_{}`$ and for any affine open $`VX`$: $`H^0(V,R^1j_{})=H^1(VU,)=0`$, as from separateness of $`X`$ it follows that $`VU`$ is an affine scheme.) Besides, it is not difficult to see that the sheaf $`j_{}𝒢`$ is an flasque sheaf. Therefore $$H^i(X,j_{}𝒢)=0\text{for any}i>0\text{.}$$ (21) Besides, the map $$H^0(X,j_{}𝒢)H^0(X,j_{}(𝒢/))$$ is surjective, as $`H^0(X,j_{}𝒢)=H^0(U,𝒢)`$, $`H^0(X,j_{}(𝒢/))=H^0(U,𝒢/)\text{.}`$ And from afinneness of $`U`$ it follows that $`H^1(U,)=0`$, therefore the following map is surjective: $$H^0(U,𝒢)H^0(U,𝒢/)\text{.}$$ Hence and from (21) we obtain that $$H^1(X,j_{})=0\text{.}$$ Further, if $`i>1`$, then from (21) and from the long cohomological sequence associated with (20) it follows that $$H^i(X,j_{})=H^{i1}(X,j_{}(𝒢/))\text{.}$$ But the sheaf $`𝒢/`$ is a quasicoherent sheaf on $`U`$. Hence, by induction, it is possible to assume that $$H^{i1}(X,j_{}(𝒢/))=0\text{.}$$ Therefore $`H^i(X,j_{})=0`$. Lemma 8 is proved. ###### Lemma 9 Let $`X`$ be a noetherian scheme. Let $`YX`$ be a closed subscheme, which is defined by the ideal sheaf $`J`$, and $`j:UX`$ be the open subscheme which is complement to the subscheme $`Y`$. Let a sheaf $``$ be a quasicoherent sheaf on $`X`$, and consider the following exact sequence of quasicoherent sheaves on $`X`$: $$0j_{}j^{}𝒢0$$ which is induced by the natural map $`j_{}j^{}`$. Let $`=\underset{i}{\underset{}{lim}}_i`$ and $`𝒢=\underset{i}{\underset{}{lim}}𝒢_i`$, $`iI`$, where the sheaves $`_i`$ $`𝒢_i`$ are coherent sheaves on the scheme $`X`$ for any $`iI`$. Then for any $`iI`$ there exists $`l(i)𝐍`$ such that $$J^{l(i)}_i=0\text{and}J^{l(i)}𝒢_i=0\text{.}$$ (22) Proof. From the exact sequence $$0J^m𝒪_X𝒪_X/J^m0$$ it follows the following sequence of quasicoherent sheaves on $`X`$ $$0om_X(𝒪_X/J^m,)om_X(𝒪_X,)om_X(J^m,)xt_X^1(𝒪_X/J^m,)\text{.}$$ Taking direct limit with respect to $`m`$, we obtain $$0\underset{}{lim}_mom_X(𝒪_X/J^m,)\underset{}{lim}_mom_X(J^m,)\underset{}{lim}_mxt_X^1(𝒪_X/J^m,)\text{.}$$ By \[9, ch. III, ex. 3.7(a)\] we have $$\underset{}{lim}_mom_X(J^m,)=j_{}j^{}\text{.}$$ Now (22) follows from $$=\underset{}{lim}_mom_X(𝒪_X/J^m,)\text{and}𝒢\underset{}{lim}_mxt_X^1(𝒪_X/J^m,)\text{.}$$ Lemma 9 is proved. ## 3 Construction and its original properties Let $`X`$ be a noetherian separated scheme. Consider a flag of closed subschemes $$XY_0Y_1\mathrm{}Y_n$$ in the scheme $`X`$. Let $`J_j`$ be the ideal sheaf of the subscheme $`Y_j`$ in $`X`$ ($`0jn`$). Let $`i_j`$ be the embedding of the subscheme $`Y_jX`$. Let $`U_i`$ be an open subscheme of $`Y_i`$ which is complement to the closed subscheme $`Y_{i+1}`$ ($`0in1`$). Let $`j_i:U_iY_i`$ be the open embedding of the subscheme $`U_i`$ to the scheme $`Y_i`$ ($`0in1`$). By definition, let $`U_n=Y_n`$ and $`j_n`$ be the identity morphism from $`U_n`$ to $`Y_n`$. Assume that for any point $`xX`$ there exists an open affine neithbourhood $`Ux`$ such that $`UU_i`$ is an affine scheme for any $`0in`$. In the sequel we’ll say that a flag of subschemes $`\{Y_i,\mathrm{\hspace{0.33em}0}in\}`$ with such condition is the flag with locally affine complements. ###### Remark 3 For example, the last condition of locally affineness of complements is satisfied in the following cases * $`Y_{i+1}`$ is the Cartier divisor on the scheme $`Y_i`$ ($`0in1`$), or * $`U_i`$ is an affine scheme for any $`0in1`$. (On a separated scheme the intersection of two open affine subschemes is an affine subscheme.) Consider the $`n`$-dimensional simplex and its standard simplicial set (without degenerations). To be precise, consider the set: $$(\{0\},\{1\},\mathrm{},\{n\})\text{.}$$ (Here are all the integers between $`0`$ and $`n`$.) Then the simplicial set $`S=\{S_k\}`$ : * $`S_0\stackrel{\mathrm{def}}{=}\{\eta \{0\},\{1\},\mathrm{},\{n\}\}`$. * $`S_k\stackrel{\mathrm{def}}{=}\{(\eta _0,\mathrm{},\eta _k),\text{where}\eta _lS_0\text{and}\eta _{l1}<\eta _l\}`$. The boundary map $`_i`$ ($`0<i<k`$) is given by eliminating the $`i`$-th component of the vector $`(\eta _0,\mathrm{},\eta _k)`$. (It is the $`i`$-th face of $`(\eta _0,\mathrm{},\eta _k)`$.) Definition. For any $`(\eta _0,\mathrm{},\eta _k)S_k`$ define the functor $$V_{(\eta _0,\mathrm{},\eta _k)}:QS(X)Sh(X)$$ uniquely determined by the following inductive conditions: 1. $`V_{(\eta _0,\mathrm{},\eta _k)}`$ commutes with direct limits. 2. If $``$ is a coherent sheaf, and $`\eta S_0`$, then $$V_\eta ()\stackrel{\mathrm{def}}{=}\underset{m𝐍}{\underset{}{lim}}(i_\eta )_{}(j_\eta )_{}(j_\eta )^{}(/J_\eta ^m)\text{.}$$ 3. If $``$ is a coherent sheaf, and $`(\eta _0,\mathrm{},\eta _k)S_k`$ ($`k1`$), then $$V_{(\eta _0,\eta _1,\mathrm{},\eta _k)}()\stackrel{\mathrm{def}}{=}\underset{m𝐍}{\underset{}{lim}}V_{(\eta _1,\mathrm{},\eta _k)}\left((i_{\eta _0})_{}(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m)\right)\text{.}$$ In the sequel, to avoid the confusion of notations in the case of a lot of schemes and flags of closed subschemes we’ll use sometimes the equivalent notation for $`V_{(\eta _0,\mathrm{},\eta _k)}()`$, in which the closed subschemes is written explicitly: $$V_{(\eta _0,\mathrm{},\eta _k)}()=V_{(Y_{\eta _0},\mathrm{},Y_{\eta _k})}(X,)\text{.}$$ ###### Proposition 1 Let $`\sigma =(\eta _0,\mathrm{},\eta _k)S_k`$. Then 1. The functor $`V_\sigma :QS(X)Sh(X)`$ is well defined. 2. The functor $`V_\sigma `$ is exact and additive. 3. The functor $`V_\sigma `$ is local on $`X`$, i. e., for any open $`UX`$ for any quasicoherent sheaf $``$ on $`X`$: $$V_{(Y_{\eta _0},\mathrm{},Y_{\eta _k})}(X,)_U=V_{(Y_{\eta _0}U,\mathrm{},Y_{\eta _k}U)}(U,_U)\text{.}$$ (Here if $`Y_jU=ø`$, then $`Y_iU`$ is an empty subscheme of $`U`$ which is defined by the ideal sheaf $`𝒪_U`$.) 4. For any quasicoherent sheaf $``$ on the scheme $`X`$ the sheaf $`V_{(\eta _0,\mathrm{},\eta _k)}()`$ is a sheaf of $`𝒪_X`$-modules with the support on the subscheme $`Y_{\eta _k}`$. (Usually, this sheaf is not quasicoherent.) 5. For any quasicoherent sheaf $``$ on $`X`$: $$V_\sigma ()=V_\sigma (𝒪_X)_{𝒪_X}\text{.}$$ 6. If all $`U_i`$ is affine ($`0in`$), then for any affine covering $`𝒰`$ of the scheme $`X`$, for any quasicoherent sheaf $``$ on $`X`$, for any $`m1`$: $$\stackrel{ˇ}{H}^m(𝒰,V_\sigma ()))=0\text{.}$$ 7. If all $`U_i`$ is affine ($`0in`$), then for any quasicoherent sheaf $``$ on $`X`$, for any $`m1`$ $$H^m(X,V_\sigma ())=0\text{.}$$ Proof. 1. Well-posedness of the definition of $`V_\sigma `$ is proved by induction by means of using of lemma 1, lemma 2, lemma 3, lemma 5, lemma 6 and lemma 7. Let us check the base of induction for lemmas 5 and 6 (when the functor $`\mathrm{\Phi }=id`$). Namely * for any affine scheme $`𝒱`$ and any quasicoherent sheaf $``$ on $`V`$ $$H^1(V,)=0\text{.}$$ * Let us show that if $`i:YX`$ is a closed subscheme with the ideal sheaf $`J`$, $`j:UY`$ is an open imbedding of the affine scheme $`U`$ in $`Y`$. Then for any quasicoherent sheaf $``$ on $`X`$, for any $`k1`$, for any affine open covering $`𝒰`$ of the scheme $`X`$, for any $`m1`$: $$\stackrel{ˇ}{H}^m(𝒰,i_{}j_{}j^{}(/J^k))=0\text{.}$$ From affineness of the covering $`𝒰`$ it follows that it is acyclic for quasicoherent sheaves. Consequently the Čech cohomologies groups with respect to this covering coincide with the usual cohomologies groups of quasicoherent sheaves. (See \[9, ch.3, theorem 4.5\]) Therefore it suffices to prove that for any integer $`k1`$ $$H^m(X,i_{}j_{}j^{}(/J^k))=0\text{.}$$ (23) If $`k=1`$, then the sheaf $`/J`$ is quasicoherent with respect to the subscheme $`Y`$, and $$H^m(X,i_{}j_{}j^{}(/J))=H^m(Y,j_{}(j^{}(/J)))=0\text{.}$$ Where the last equality follows from lemma 8. If $`k>1`$, then by lemma 4 the following sequence is exact $$0i_{}j_{}j^{}(J^{k1}/J^k)i_{}j_{}j^{}(/J^k)i_{}j_{}j^{}(/J^{k1})0\text{.}$$ (24) In addition, the sheaves $`J^{k1}/J^k`$ and $`/J^{k1}`$ are quasicoherent with respect to the subscheme $`(Y,𝒪_X/J^{k1})`$. Therefore, by induction, we obtain $$H^m(X,i_{}j_{}j^{}(J^{k1}/J^k))=0\text{and}H^m(X,i_{}j_{}j^{}(/J^{k1}))=0\text{.}$$ Hence, from (24) we have (23). Item 1 of proposition 1 is proved. 2. The proof of this item is analogous to the proof of item 1 by means of the same lemmas. 3. Localness follows by induction from the construction of the functor $`V_{(\eta _0,\mathrm{},\eta _k)}`$. 4. This item follows by induction from the construction. 5. We have the natural map: $$V_\sigma ()\text{,}$$ which induces the following map: $$V_\sigma (𝒪_X)_{𝒪_X}V_\sigma (𝒪_X)_{𝒪_X}V_\sigma ()V_\sigma (𝒪_X)_{V_\sigma (𝒪_X)}V_\sigma ()=V_\sigma ()\text{.}$$ (25) Let us show that (25) gives us an isomorphism between $`V_\sigma (𝒪_X)_{𝒪_X}`$ and $`V_\sigma ()`$. Since the functor $`V_\sigma `$ and tensor products commute with direct limits, we can assume that $``$ is a coherent sheaf. In view of item 3 of this proposition, we can restrict ourself to the local situation. That is, we suppose $`X=\mathrm{Spec}A`$, $`=\stackrel{~}{M}`$ for some finitely generated $`A`$-module $`M`$. Then for some $`r𝐍`$ there exists an exact sequence of sheaves as: $$0\stackrel{~}{N}𝒪_X^r\stackrel{~}{M}0\text{,}$$ where $`N`$ is some finitely generated $`A`$-module. Hence we obtain the commutative diagram: $$\begin{array}{ccccccc}& V_\sigma (𝒪_X)_{𝒪_X}\stackrel{~}{N}& & V_\sigma (𝒪_X)_{𝒪_X}𝒪_X^r& & V_\sigma (𝒪_X)_{𝒪_X}\stackrel{~}{M}& 0\\ & \text{}& & \text{}& & \text{}& \\ 0& V_\sigma (\stackrel{~}{N})& & V_\sigma (𝒪_X^r)& \stackrel{\delta }{}& V_\sigma (\stackrel{~}{M})& 0\text{,}\end{array}$$ where the lower row is exact by virtue of item 2. Besides, it is clear that $`\beta `$ is an isomorphism. Therefore from surjectivity of $`\delta `$ it follows that $`\alpha `$ is surjective. Since $`\stackrel{~}{N}`$ is a coherent sheaf, we have that the map $`\gamma `$ is surjective as well. Hence, from exactness of the lower row and non complicated diagram search it follows that the map $`\alpha `$ is injective. 6. The proof is similar to the proof of item 1 by means of the same lemmas. 7. This item follows from the previous item of this proposition and lemma 7. (Since every point has arbitrarily small affine neithbourhood with affine intersection to all $`U_i`$, we obtain that this affine neithbourhood satisfies item 6 of proposition 1.) ###### Proposition 2 1. Let $`X`$ be a noetherian separated scheme. Let $$Y_0Y_1\mathrm{}Y_n\text{and}Y_0^{}Y_1^{}\mathrm{}Y_n^{}$$ be two flags of closed subschemes in $`X`$ with the corresponding ideal sheaves $`J_i`$ and $`J_i^{}`$ ($`0in`$) such that for any $`0in`$ there exist integers $`l_i1`$ and $`l_i^{}1`$ with the following properties: $$J_i^{l_i}J_i^{}\text{and}(J_i^{})^{l_i^{}}J_i\text{.}$$ (26) Then the functors $$V_{(Y_{\eta _0},\mathrm{},Y_{\eta _k})}(X,)\text{and}V_{(Y_{\eta _0}^{},\mathrm{},Y_{\eta _k}^{})}(X,)$$ coincides for any $`(\eta _0,\mathrm{},\eta _k)S_k`$. 2. Consider the flag of closed subschemes: $$XZY_0\mathrm{}Y_n$$ on a noetherian separated scheme $`X`$. Let $`i:ZX`$ be a closed imbedding. Then for any quasicoherent sheaf $``$ on the scheme $`Z`$ we have $$i_{}\left(V_{(Y_0,\mathrm{},Y_n)}(Z,)\right)=V_{(Y_0,\mathrm{},Y_n)}(X,i_{})\text{.}$$ ###### Remark 4 Condition (26) is equivalent to the statement that the topological spaces of subschemes $`Y_i`$ and $`Y_i^{}`$ are the same for all $`0in`$. Proof (of proposition 2). 1. It suffices to prove for the case when $`J_i=J_i^{}`$ for all $`ij`$, $`0in`$, where some fixing $`0jn`$. From inductance of the definition of the functor $`V_{(Y_{\eta _0},\mathrm{},Y_{\eta _k})}(X,)`$ (and $`V_{(Y_{\eta _0}^{},\mathrm{},Y_{\eta _k}^{})}(X,)`$) we can restrict ourself to the case $`j=0`$. This case follows from cofinality of the projective systems $`/J_j^k`$ and $`/(J_j^{})^k`$ (from condition 26). 2. This item follows at once from the construction and the fact that supports of all appearing from induction sheaves are on the subscheme $`Z`$. ## 4 Complexes and their exactness Consider again the usual $`n`$-simplex without degenerations $`S=\{S_k,\mathrm{\hspace{0.33em}0}kn\}`$. If $`\sigma =(\eta _0,\mathrm{},\eta _k)S_k`$, then $`_i(\sigma )`$ is the $`i`$-th face of $`\sigma `$ ($`0ik`$). Then define the morphism of functors $$d_i(\sigma ):V_{_i(\sigma )}V_\sigma \text{,}\text{as}$$ commuting with direct limits, and on coherent sheaves it is a map $$V_{_i(\sigma )}()V_\sigma ()$$ (27) which is defined by the following rules: * if $`i=0`$, then (27) is obtained from application of the functor $`V_{_0(\sigma )}`$ to the map $$(i_{\eta _0})_{}(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m)$$ and passage to the projective limit on $`m`$; * if $`i=1`$, $`k=1`$, then we have the natural map $$(i_{\eta _0})_{}(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m)V_{(\eta _1)}((i_{\eta _0})_{}(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))\text{.}$$ Now after passage to the projective limit on $`m`$ we obtain the map (27) in this case. * $`i0`$, $`k>1`$, then from induction on $`k`$ we can suppose that we have the map $$V_{_{i1}(_0(\sigma ))}((i_{\eta _0})_{}(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))V_{_0(\sigma )}((i_{\eta _0})_{}(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))\text{.}$$ And passage to the projective limit on $`m`$ gives us the map (27) in this case. ###### Proposition 3 For any $`1kn`$, $`0ik`$ define $$d_i^k\stackrel{\mathrm{def}}{=}\underset{\sigma S_k}{}d_i(\sigma ):\underset{\sigma S_{k1}}{}V_\sigma \underset{\sigma S_k}{}V_\sigma \text{.}$$ Also define $$d_0^0:id\underset{\sigma S_0}{}V_\sigma $$ as the direct sum of the natural maps $`V_\sigma ()`$. (Here $`id`$ is the functor of the natural imbedding of $`QS(X)`$ into $`Sh(X)`$, $``$ is a quasicoherent sheaf on $`X`$, $`\sigma S_0`$.) Then for all $`0i<jkn1`$ we have $$d_j^{k+1}d_i^k=d_i^{k+1}d_{j1}^k\text{.}$$ (28) Proof. Using the inductance of the definition, the proof is done by induction from non complicated consideration of some cases. It suffices to consider the small $`i`$ and $`k`$ only. (For example, see similar cases in \[11, §2.4\] or .) Define $$d^m\stackrel{\mathrm{def}}{=}\underset{0im}{}(1)^id_i^m$$ Then proposition 3 makes possible to construct the complex of sheaves $`V()`$ from any quasicoherent sheaf $``$ on $`X`$ in the standard way: $$\mathrm{}\underset{\sigma S_{m1}}{}V_\sigma ()\stackrel{d^m}{}\underset{\sigma S_m}{}V_\sigma ()\mathrm{}$$ Where $`d^{m+1}d^m=0`$ follows from (28) by means of non complicated direct calculations. ###### Theorem 1 Let $`X`$ be a noetherian separated scheme. Let $`Y_0Y_1\mathrm{}Y_n`$ be a flag of closed subschemes with locally affine complements. Assumee that $`Y_0=X`$. Then the following complex is exact: $$0\stackrel{d^0}{}V()0\text{.}$$ (29) Proof. It suffices to consider only the case when the sheaf $``$ is coherent. Consider the exact sequence of sheaves $$0(j_0)_{}(j_0)^{}𝒢0$$ (30) Here $``$ and $`𝒢`$ is the kernel and the cokernel of the natural map of sheaves $`(j_0)_{}(j^0)^{}`$. From exactness of functors $`V_\sigma `$ (for any $`\sigma `$) we obtain the following exact sequence of complexes of sheaves: $$0V()V()V((j_0)_{}(j_0)^{})V(𝒢)0$$ (31) By lemma 9 the supports of sheaves $``$ and $`𝒢`$ are on $`Y_1`$, therefore in the case $`\eta _0=0`$ we have $`V_{(Y_{\eta _0},\mathrm{},Y_{\eta _k})}(X,)=0`$ and $`V_{(Y_{\eta _0},\mathrm{},Y_{\eta _k})}(X,𝒢)=0`$. Therefore, using it, lemma 9 (which decompose the sheaves $``$ and $`𝒢`$ in direct limits of sheaves which is coherent on subschemes with topological space $`Y_1`$), proposition 2, permutability of the functors $`V_\sigma `$ with direct limits, we can apply induction on the length of flag and suppose that the complexes $$0\stackrel{d^0}{}V()0\text{and}$$ (32) $$0𝒢\stackrel{d^0}{}V(𝒢)0$$ (33) are already exact. It is not difficult to understand that for any $`\sigma =(\eta _0,\mathrm{},\eta _k)S_k`$ if $`\eta _0=0`$, then $`V_\sigma ((j_0)_{}(j_0)^{})=V_\sigma ()`$; if $`\eta _00`$, then $`V_\sigma ((j_0)_{}(j_0)^{})=V_\sigma ^{}()`$, where $`\sigma ^{}=(0,\eta _0,\mathrm{},\eta _k)S_{k+1}`$. Hence, the complex $`V((j_0)_{}(j_0)^{})`$ has the same components $`V_\sigma ^{}()`$ in the degree $`k`$ and $`k+1`$. Therefore, successively from the highest degrees spliting off the trivial complexes $$0V_\sigma ((j_0)_{}(j_0)^{})V_\sigma ^{}((j_0)_{}(j_0)^{})0\text{,}$$ we obtain exactness of the complex $$0(j_0)_{}(j_0)^{}\stackrel{d^0}{}V((j_0)_{}(j_0)^{})0\text{.}$$ (34) Now, since complexes (32), (33), (34) are exact, we obtain exactness of complex (29) from exactness of (31) and (30). Theorem 1 is proved. For any $`\sigma S_k`$ define $$A_\sigma ()\stackrel{\mathrm{def}}{=}H^0(X,V_\sigma ())\text{.}$$ ###### Proposition 4 Let $`X`$ be a noetherian separated scheme. $`Y_0Y_1\mathrm{}Y_n`$ be a flag of closed subschemes such that all $`U_i`$ are affine ($`0in`$). Let $`\sigma S_k`$ be arbitrary. Then 1. $`A_\sigma `$ is an exact and additive functor: $`QS(X)Ab`$. 2. If $`X=\mathrm{Spec}A`$, $`M`$ is some $`A`$-module, then $$A_\sigma (\stackrel{~}{M})=A_\sigma (𝒪_X)_AM\text{.}$$ Proof. 1. This item follows at once from items 2 and 7 of proposition 1. 2. Similarly to the proof of item 5 of proposition 1 we can suppose that the module $`M`$ is finitely generated over $`A`$. Now consider the exact sequence of $`A`$-modules: $$0NA^^rM0\text{.}$$ Hence we obtain the commutative diagramm: $$\begin{array}{ccccccc}& A_\sigma (𝒪_X)_A\stackrel{~}{N}& & A_\sigma (𝒪_X)^^r& & A_\sigma (𝒪_X)_AM& 0\\ & \text{}& & \text{}& & \text{}& \\ 0& A_\sigma (\stackrel{~}{N})& & A_\sigma (𝒪_X^r)& \stackrel{\delta }{}& A_\sigma (\stackrel{~}{M})& 0\text{,}\end{array}$$ where the lower row is exact by virtue of item 1 of this proposition. It is clear that $`\beta `$ is an isomorphism. Therefore, arguing as in item 5, we obtain at first surjectivity of the map $`\alpha `$, and afterwards we obtain injectivity of $`\alpha `$. Proposition 4 is proved. Let $``$ be any quasicoherent sheaf on $`X`$. Apply the functor $`H^0(X,)`$ to the complex $`V()`$. We obtain the complex of abelian groups $`A()`$: $$\mathrm{}\underset{\sigma S_{m1}}{}A_\sigma ()\underset{\sigma S_m}{}A_\sigma ()\mathrm{}\text{.}$$ ###### Theorem 2 Let $`X`$ be a noetherian separated scheme. Let $`Y_0Y_1\mathrm{}Y_n`$ be a flag of closed subschemes such that $`Y_0=X`$ and all $`U_i`$ are affine ($`0in`$). Then cohomology of the complex $`A()`$ coincide with cohomology of the sheaf $``$ on $`X`$, i. e., for any $`i`$ $$H^i(X,)=H^i(A())\text{.}$$ Proof. From theorem 1 and item 7 of proposition 1 it follows that $`V()`$ is an acyclic resolution for the sheaf $``$. Therefore it is possible to calculate cohomology of the sheaf $``$ by means of global sections of this resolution. Theorem 2 is proved. From the last theorem we obtain at once the following geometrical corollary. ###### Theorem 3 Let $`X`$ be a projective algebraic scheme of dimension $`n`$ over a field. Let $`Y_0Y_1\mathrm{}Y_n`$ be a flag of closed subschemes such that $`Y_0=X`$ and $`Y_i`$ is an ample divisor on the scheme $`Y_{i1}`$ for any $`1in`$. Then for any quasicoherent sheaf $``$ on $`X`$, for any $`i`$ we have $$H^i(X,)=H^i(A())\text{.}$$ Proof. In fact, since $`Y_i`$ is an ample divisor on $`Y_{i1}`$ for all $`1in`$,we have that $`U_i`$ is an affine scheme for all $`0in1`$. Since $`dimY_n=0`$, we have that $`U_n=Y_n`$ is affine as well. Now application of theorem 2 concludes the proof. ###### Remark 5 Let us remark that for any quasicoherent sheaf $``$, for any $`\sigma =(\eta _0)S_0`$$`A_\sigma ()`$ is the group of section over $`U_{\eta _0}`$ of the sheaf $``$ lifted to the formal neighbourhood of the subscheme $`Y_{\eta _0}`$ in $`X`$. And the complex $`A()`$ can be interpreted as the Čech complex for the such ”covering ” of the scheme $`X`$. ## 5 Combinatorial properties and the Krichever map. ###### Lemma 10 Let $`X`$ be a noetherian separated scheme. Let $`Y_0Y_1\mathrm{}Y_n`$ be a flag of closed subschemes such that $`Y_0=X`$ and $`Y_i`$ is an ample Cartier divisor on the scheme $`Y_{i1}`$ ($`1in`$). Let $`J_i`$ be the ideal sheaves on $`X`$ defining the corresponding subschemes $`Y_i`$ in $`X`$. Let $`\sigma =(\eta _0,\mathrm{},\eta _k)S_k`$. Then for any $`i\eta _0`$, for any quasicoherent sheaf $``$ on $`X`$ we have $$A_\sigma ()=\underset{m}{\underset{}{lim}}A_\sigma (/J_i^m)\text{.}$$ (35) ###### Remark 6 We consider the sheaf $`/J_i^m`$ in (35) as the sheaf on the scheme $`X`$. The corresponding functor of direct image from the subscheme $`Y_i`$ is omitted for the sake of simplication of notations. Further we shall do the same in analogous situations. Proof. From the definition of the functor $`A_\sigma `$ we have $$\underset{m}{\underset{}{lim}}A_\sigma (/J_i^m)=\underset{m}{\underset{}{lim}}\underset{l}{\underset{}{lim}}A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}(/(J_i^m+J_{\eta _0}^l)))=$$ $$=\underset{(m,l)}{\underset{}{lim}}A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}(/(J_i^m+J_{\eta _0}^l)))=\underset{l}{\underset{}{lim}}A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^l))=A_\sigma ()\text{,}$$ where next to the last equality follows from cofinality of the projective systems $$A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}(/(J_i^m+J_{\eta _0}^l)))\text{and}A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^l))\text{.}$$ The last follows from cofinality of the systems $`/(J_i^m+J_{\eta _0}^l)`$ and $`/J_{\eta _0}^l`$. Besides, in our reasonings we meant that if $`\sigma S_0`$, then $`A_{_0(\sigma )}=H^0(X,)`$. The lemma is proved. ###### Lemma 11 Let $`X`$ be a Cohen-Macaulay noetherian scheme. Let $`Y_0Y_1\mathrm{}Y_n`$ be a flag of closed subschemes such that $`Y_0=X`$ and $`Y_i`$ is a Cartier divisor on the scheme $`Y_{i1}`$ ($`1in`$). Let $`J_i`$ be the ideal sheaves on $`X`$ defining the corresponding subschemes $`Y_i`$ in $`X`$. Let $`j_k`$ be an open imbedding of $`Y_k\backslash Y_{k+1}`$ into $`Y_k`$ ($`0k<n`$). Then for any $`0k<n`$, for any $`m1`$, for any locally free sheaf $``$ on $`X`$ the natural map $$/J_k^m(j_k)_{}(j_k)^{}(/J_k^m)$$ is an imbedding. Proof. Let us do induction on $`m`$. Let $`m=1`$. Then $`/J_k`$ is a locally free sheaf on $`Y_k`$. Therefore, applying lemma 9 to the pair $`Y_{k+1}Y_k`$, we obtain injectivity in this case, since Cartier divisor is locally generated by one element which is not divisor of zero in the structure sheaf. Let $`m>1`$. Then since $`Y_i`$ are Cartier divisors on $`Y_{i1}`$, we have that $`Y_k`$ is a local complete intersection in $`X`$ and locally defined by a regular sequence on $`X`$. Therefore the sheaf $`J_k^{m1}/J_k^m`$ is locally free on $`Y_k`$ (see \[9, ch. II, th. 8.21A\]). From the exact sequence $$0J_k^{m1}/J_k^m/J_k^m/J_k^{m1}0$$ (36) and lemma 4 we obtain exactness of the following sequence $$0(j_k)_{}(j_k)^{}(J_k^{m1}/J_k^m)(j_k)_{}(j_k)^{}(/J_k^m)(j_k)_{}(j_k)^{}(/J_k^{m1})0\text{.}$$ (37) The sheaf $`J_k^{m1}/J_k^m=_{𝒪_X}J_k^{m1}/J_k^m`$ is locally free $`Y_k`$. Therefore the map $$J_k^{m1}/J_k^m(j_k)_{}(j_k)^{}(J_k^{m1}/J_k^m)$$ is an imbedding. Also the map $$/J_k^{m1}(j_k)_{}(j_k)^{}(/J_k^{m1})$$ is an imbedding by induction hypothesis. From this, (36), (37) and non complicated diagram search we obtain that the map $$/J_k^m(j_k)_{}(j_k)^{}(/J_k^m)$$ is an imbedding. The lemma is proved. ###### Lemma 12 Let $`X`$ be a projective equidimensional Cohen-Macaulay algebraic scheme of dimension $`n`$ over a field. Let $`Y_0Y_1\mathrm{}Y_k`$ be a flag of closed subschemes such that $`Y_0=X`$ and $`Y_i`$ is an ample Cartier divisor on the scheme $`Y_{i1}`$ for any $`i`$ ($`1ik`$). Let $`J_i`$ be the ideal sheaves on $`X`$ defining the corresponding subschemes $`Y_i`$ in $`X`$. Then for any locally free sheaf $``$ on $`X`$ the natural map $$H^0(X,)\underset{m}{\underset{}{lim}}H^0(X,/J_k^m)$$ is an imbedding; and if $`k<n`$, then this one is an isomorphism. Proof. The proof will be done by induction on $`k`$. Let at first $`k=1`$. The sheaf $`J_1`$ is the dual of the ample invertible sheaf on $`X`$. And from conditions on $`X`$ (that is Cohen-Macaulayness, projectiveness , equidimensionality) there exist (see \[9, ch. III, th. 7.6\]) $`l>0`$ such that for any $`m>l`$ we have $`H^0(X,J_1^m)=0`$, and if $`n2`$, then we have $`H^1(X,J_1^m)=0`$ as well. Hence and from the exact sequence $$0J_1^m/J_1^m0$$ we obtain that the map $$H^0(X,)H^0(X,/J_1^m)$$ is an imbedding for $`m>l`$, and $$H^0(X,)=H^0(X,/J_1^m)$$ for $`m>l`$ and $`n2`$. And after passage to the projective limit on $`m`$ we obtain the imbedding $$H^0(X,)\underset{m}{\underset{}{lim}}H^0(X,/J_1^m)=H^0(X,\underset{m}{\underset{}{lim}}/J_1^m)\text{,}$$ and if $`n>1`$, then this one is the isomorphism $$H^0(X,)H^0(X,\underset{m}{\underset{}{lim}}/J_1^m)\text{.}$$ Now let $`k>1`$ be arbitrary. From lemmas conditions it follows that $`Y_k`$ is locally defined by a regular sequence on $`X`$. Therefore $`Y_k`$ is a Cohen-Macaulay scheme, and $`J_k^m/J_k^{m+1}`$ are locally free sheaves on $`Y_k`$ (see \[9, ch. II, th. 8.21A\]). By induction hypothesis we have $$\underset{l}{\underset{}{lim}}H^0(X,/J_{k1}^l)=H^0(X,)\text{.}$$ (38) From cofinality of projective systems $`/J_k^m`$ and $`/(J_k^m+J_{k1}^l)`$ we have that $$\underset{m}{\underset{}{lim}}H^0(X,/J_k^m)=\underset{(l,m)}{\underset{}{lim}}H^0(X,/(J_k^m+J_{k1}^l))=\underset{l}{\underset{}{lim}}\underset{m}{\underset{}{lim}}H^0(X,/(J_k^m+J_{k1}^l))\text{.}$$ Frim this one and equality (38) it suffices for the proof of the lemma to show that for any $`l1`$ the map $$H^0(X,/J_{k1}^l)\underset{m}{\underset{}{lim}}H^0(X,/(J_k^m+J_{k1}^l))$$ is an imbedding, and if $`k<n`$ then this map is an isomorphism. For this one let us consider the exact sequence $$0\frac{J_k^m+J_{k1}^l}{J_{k1}^l}/J_{k1}^l/(J_k^m+J_{k1}^l)0\text{.}$$ It suffices to show that for all sufficiently large $`m`$ $$H^0(X,\frac{J_k^m+J_{k1}^l}{J_{k1}^l})=0\text{,}$$ (39) and if $`k<n`$, then $$H^1(X,\frac{J_k^m+J_{k1}^l}{J_{k1}^l})=0\text{.}$$ (40) For this one let us do induction on $`l`$. If $`l=1`$, then (39) and (40) follows at once from \[9, ch. III, th. 7.6\] and the fact that the sheaf $`J_k/J_{k1}`$ is the dual of the ample invertible sheaf on $`Y_{k1}`$. If $`l>1`$, then from the identity $`\frac{A+B}{B}=\frac{A}{AB}`$ it follows the exact sequence $$0\frac{J_k^mJ_{k1}^{l1}}{J_k^mJ_{k1}^l}_{𝒪_X}\frac{J_k^m+J_{k1}^l}{J_{k1}^l}\frac{J_k^m+J_{k1}^{l1}}{J_{k1}^{l1}}0\text{.}$$ Restricting ourself to the local situation and using the fact that $`J_i`$ is generated by a regular sequence and \[9, ch. II, th. 8.21A\], it is not difficult to understand that $$\frac{J_k^mJ_{k1}^{l1}}{J_k^mJ_{k1}^l}_{𝒪_X}=\frac{J_k^{ml+1}+J_{k1}}{J_{k1}}\left(\frac{J_{k1}^{l1}}{J_{k1}^l}_{𝒪_X}\right)\text{.}$$ In addition, the sheaf $`\frac{J_{k1}^{l1}}{J_{k1}^l}_{𝒪_X}`$ is locally free on $`Y_{k1}`$. The sheaf $`J_k/J_{k1}`$ is the dual of the ample invertible sheaf on the Cohen-Macaulay scheme $`Y_{k1}`$. Therefore from \[9, ch. III, th. 7.6\] we have for sufficiently large $`m`$ that: $$H^0(X,\frac{J_k^mJ_{k1}^{l1}}{J_k^mJ_{k1}^l}_{𝒪_X})=0\text{, and}$$ $$\text{if}k<n\text{, then}H^1(X,\frac{J_k^mJ_{k1}^{l1}}{J_k^mJ_{k1}^l}_{𝒪_X})=0\text{.}$$ The lemma is proved. Corollary(from lemma 12) Under the conditions of lemma 12 for any $`\sigma S_0`$ the natural map $$H^0(X,)A_\sigma ()$$ is an imbedding. Proof. Let $`\sigma =(m)`$. By lemma 12 we have the imbedding $$0/J_k^m(j_k)_{}(j_k)^{}(/J_k^m)\text{.}$$ Hence we obtain the imbedding $$0H^0(X,/J_k^m)H^0(X,(j_k)_{}(j_k)^{}(/J_k^m))\text{.}$$ After passage to the projective limit on $`m`$ we obtain the imbedding $$0\underset{m}{\underset{}{lim}}H^0(X,/J_k^m)A_\sigma ()\text{.}$$ Now application of lemma 12 concludes the proof of the corollary. ###### Theorem 4 Let $`X`$ be a projective equidimensional Cohen-Macaluay scheme of dimension $`n`$ over a field. Let $`Y_0Y_1\mathrm{}Y_n`$ be a flag of closed subschemes such that $`Y_0=X`$ and $`Y_i`$ is an ample Cartier divisor on the scheme $`Y_{i1}`$ for any $`1in`$. Then for any locally free sheaf $``$ on $`X`$ we have that 1. for any $`\sigma S_k`$ ($`0kn`$) the natural map $$H^0(X,)A_\sigma ()$$ is an imbedding, 2. for any $`\sigma S_k`$ ($`1kn`$), for any $`i`$ ($`0ik`$) the natural map $$d_i(\sigma ):A_{_i(\sigma )}()A_\sigma ()$$ ia an imbedding. ###### Remark 7 Taking into account (28) from proposition 3, it is possible to reformulate item 2 of this theorem in the following way: for any locally free sheaf $``$ on $`X`$, for any $`\sigma _1,\sigma _2S`$, $`\sigma _1\sigma _2`$ the natural map $`A_{\sigma _1}()A_{\sigma _2}()`$ ia an imbedding. Proof. Let $`J_i`$ be the ideal sheaves on $`X`$ defining the corresponding subschemes $`Y_i`$ in $`X`$. Let us prove first item 2 of the theorem. Consider 3 cases. Case 1. Let $`\sigma =(\eta _0,\eta _1,\mathrm{},\eta _k)`$, and $`i=0`$. Then $`_0(\sigma )=(\eta _1,\mathrm{},\eta _k)`$. By lemma 11 for any $`m1`$ the map $$/J_{\eta _0}^m(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m)$$ is an imbedding. Apply to this sequence the exact functor $`A_{_0(\sigma )}`$. We obtain the imbedding: $$A_{_0(\sigma )}(/J_{\eta _0}^m)A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))\text{.}$$ After passage to the projective limit on $`m`$ we obtain the imbedding $$\underset{m}{\underset{}{lim}}A_{(\eta _1,\mathrm{},\eta _k)}(/J_{\eta _0}^m)A_\sigma ()\text{.}$$ In addition, from lemma 10 we have that $$\underset{m}{\underset{}{lim}}A_{(\eta _1,\mathrm{},\eta _k)}(/J_{\eta _0}^m)=A_{(\eta _1,\mathrm{},\eta _k)}()\text{.}$$ Thus we obtain that in this case the map $$A_{_0(\sigma )}()A_\sigma ()$$ is an imbedding. Case 2. Let $`\sigma =(\eta _0,\eta _1)`$, and $`i=1`$. In this case we have that $$A_{_1(\sigma )}()=\underset{m}{\underset{}{lim}}H^0(X,(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))$$ $$A_\sigma ()=\underset{m}{\underset{}{lim}}A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))\text{,}$$ and for the proof of this case it suffices to show that for any $`m1`$ the map $$H^0(X,(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))$$ (41) is an imbedding. Let us show this by induction. Let $`m=1`$. Then $`/J_{\eta _0}`$ is a locally free sheaf on $`Y_{\eta _0}`$. Besides, since $`Y_{\eta _0+1}`$ is a Cartier divisor on $`Y_{\eta _0}`$, we have from item 5 of proposition 1 that $$(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0})=((j_{\eta _0})_{}(j_{\eta _0})^{}𝒪_{Y_{\eta _0}})_{𝒪_{Y_{\eta _0}}}(/J_{\eta _0})=$$ $$=\underset{j}{\underset{}{lim}}𝒪_{Y_{\eta _0}}(jY_{\eta _0+1})_{𝒪_{Y_{\eta _0}}}(/J_{\eta _0})=\underset{j}{\underset{}{lim}}(/J_{\eta _0})(jY_{\eta _0+1})\text{.}$$ (42) The sheaves $`(/J_{\eta _0})(jY_{\eta _0+1})`$ are locally free on $`Y_{\eta _0}`$ as well. Therefore by corollary from lemma 12 the map $$H^0(X,(/J_{\eta _0})(jY_{\eta _0+1}))A_{_0(\sigma )}((/J_{\eta _0})(jY_{\eta _0+1}))$$ is an imbedding. After passage to the projective limit on $`j`$ we obtain injectivity of (41) in the case $`m=1`$. If $`m>1`$, then the statement will follow from induction hypothesis and consideration of the following two exact sequences: $$\begin{array}{c}0H^0(X,(j_{\eta _0})_{}(j_{\eta _0})^{}𝒢)H^0(X,(j_{\eta _0})_{}(j_{\eta _0})^{}/J_{\eta _0}^m)\hfill \end{array}$$ $$\begin{array}{c}\hfill H^0(X,(j_{\eta _0})_{}(j_{\eta _0})^{}/J_{\eta _0}^{m1})0\end{array}$$ $$\begin{array}{c}0A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}𝒢)A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}/J_{\eta _0}^m)\hfill \end{array}$$ $$\begin{array}{c}\hfill A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}/J_{\eta _0}^{m1})0\text{,}\end{array}$$ where the sheaf $`𝒢=J_{\eta _0}^{m1}/J_{\eta _0}^m`$ is a locally free on $`Y_{\eta _0}`$. Case 2 is analyzed. Case 3. We shall consider all that are not in cases 1 and 2. Let $`\sigma =(\eta _0,\eta _1,\mathrm{},\eta _k)`$ and $`i0`$. Let us do induction on $`i`$. The case $`i=0`$ is already analyzed (case 1). We have $$A_{_i(\sigma )}()=\underset{}{lim}_mA_{_{i1}(_0(\sigma ))}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))\text{and}$$ $$A_\sigma ()=\underset{}{lim}_mA_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))\text{.}$$ By induction hypothesis applied to the scheme $`Y_{\eta _0}`$ we can suppoce that for any locally free sheaf $``$ on $`Y_{\eta _0}`$ the map $$0A_{_{i1}(_0(\sigma ))}()A_{_0(\sigma )}()$$ (43) is an imbedding. Let us show that for any $`m`$ the map $$A_{_{i1}(_0(\sigma ))}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))$$ (44) is an imbedding. If $`m=1`$, then, using (42), as in case 2, we reduce all at once to the sequence (43). Afterwards we pass to the direct limit. If $`m>1`$, then as in case 2 we can do induction on $`m`$ by means of using of two following exact sequences: $$\begin{array}{c}0A_{_{i1}(_0(\sigma ))}(j_{\eta _0})_{}(j_{\eta _0})^{}𝒢)A_{_{i1}(_0(\sigma ))}((j_{\eta _0})_{}(j_{\eta _0})^{}/J_{\eta _0}^m)\hfill \end{array}$$ $$\begin{array}{c}\hfill A_{_{i1}(_0(\sigma ))}((j_{\eta _0})_{}(j_{\eta _0})^{}/J_{\eta _0}^{m1})0\end{array}$$ $$\begin{array}{c}0A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}𝒢)A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}/J_{\eta _0}^m)\hfill \end{array}$$ $$\begin{array}{c}\hfill A_{_0(\sigma )}((j_{\eta _0})_{}(j_{\eta _0})^{}/J_{\eta _0}^{m1})0\text{,}\end{array}$$ where the sheaf $`𝒢=J_{\eta _0}^{m1}/J_{\eta _0}^m`$. Now after passage in (44) to the limit we conclude the proof of case 3. Now consider item 1 of the theorem. If $`k=0`$, then this is corollary of lemma 12. If $`k>0`$, then consider the map $$H^0(X,)A_\sigma ()$$ as composition of the maps $$H^0(X,)A_{_0(\sigma )}()\text{and}A_{_0(\sigma )}()A_\sigma ()\text{,}$$ where we can suppose that by induction on $`k`$ the first map is injective, and by item 2 of this theorem the second map is injective as well. Theorem 4 is proved. ###### Lemma 13 Let all the conditions of theorem 4 be satisfied. By $`J_i`$ denote the ideal sheaves on $`X`$ defining the corresponding $`Y_i`$ in $`X`$. Then for any locally free sheaf $``$ on $`X`$, for any $`m>0`$ we have that 1. the map $$H^0(X,/J_i^m)A_\sigma (/J_i^m)$$ is an imbedding for any $`\sigma =(\zeta _0,\mathrm{},\zeta _k)`$, $`0i\zeta _0`$. 2. the map $$A_{\sigma _1}(/J_i^m)A_{\sigma _2}(/J_i^m)$$ is an imbedding for any $`\sigma _1,\sigma _2S`$, $`\sigma _1\sigma _2=(\eta _0,\mathrm{},\eta _k)`$ , $`0i\eta _0`$. Proof. Let us show item 1. If $`m=1`$, then this follows from theorem 4, which is applied to the scheme $`Y_i`$. If $`m>1`$, then apply induction. For this consider the exact sequence $$0J_i^{m1}/J_i^m/J_i^m/J_i^{m1}0\text{.}$$ Hence and from exactness of the functor $`A_\sigma `$ we have the following commutative diagram: $$\begin{array}{ccccccc}0& H^0(X,J_i^{m1}/J_i^m)& & H^0(X,/J_i^m)& & H^0(X,/J_i^{m1})& \\ & \text{}& & \text{}& & \text{}& \\ 0& A_\sigma (J_i^{m1}/J_i^m)& & A_\sigma (/J_i^m)& & A_\sigma (/J_i^{m1})& 0\text{.}\end{array}$$ The sheaf $`J_i^{m1}/J_i^m=\frac{J_i^{m1}}{J_i^m}_{𝒪_X}`$ is locally free on $`Y_i`$ (see \[9, ch. II, th. 8.21A\]). Therefore the map $`\alpha `$ is an imbedding. The map $`\gamma `$ is an imbedding by the induction hypothesis. Now from non complicated diagram search it follows that $`\beta `$ is an imbedding as well. Item 1 is proved. Item 2 follows at once from analogous to item 1 and inductive on $`m`$ reasonings applied to the following diagram: $$\begin{array}{ccccccc}0& A_{\sigma _1}(J_i^{m1}/J_i^m)& & A_{\sigma _1}(/J_i^m)& & A_{\sigma _1}(/J_i^{m1})& \\ & \text{}& & \text{}& & \text{}& \\ 0& A_{\sigma _2}(J_i^{m1}/J_i^m)& & A_{\sigma _2}(/J_i^m)& & A_{\sigma _2}(/J_i^{m1})& 0\text{.}\end{array}$$ The lemma is proved. ###### Theorem 5 Let all the conditions of theorem 4 be satisfied. Then for any locally free sheaf $``$, for any $`\sigma _1,\sigma _2S`$ we have that 1. if $`\sigma _1\sigma _2=ø`$, then $$A_{\sigma _1}()A_{\sigma _2}()=H^0(X,)\text{;}$$ 2. if $`\sigma _1\sigma _2ø`$, then $$A_{\sigma _1}()A_{\sigma _2}()=A_{\sigma _1\sigma _2}()\text{.}$$ ###### Remark 8 According to theorem 4, the intersections make sense, because we can always imbed $`A_{\sigma _1}()`$ and $`A_{\sigma _2}()`$ into $`A_\eta ()`$, where $`\eta `$ contains $`\sigma _1`$ and $`\sigma _2`$. For instance, $`\eta =\sigma _1\sigma _2`$. Proof. As usually, by $`J_i`$ denote the ideal sheaves on $`X`$ defining the corresponding subschemes $`Y_i`$ in $`X`$. Let us show item 1. Let $`\sigma _1=(\eta _0,\mathrm{})`$, $`\sigma _2=(\zeta _0,\mathrm{})`$, $`\sigma _1\sigma _2=ø`$. Without loss of generality it can be assumed that $`\zeta _0>\eta _0`$. Assume that $`\sigma _1S_0`$. By lemma 11 for any $`m>1`$ we have the exact sequence: $$0/J_{\eta _0}^m(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m)𝒢_m0\text{,}$$ where the sheaf $`𝒢_m=\frac{(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m)}{/J_{\eta _0}^m}`$. Let us show by induction on $`m`$ that the natural map $$A_{_0(\sigma _1)}(𝒢_m)A_{_0(\sigma _1)\sigma _2}(𝒢_m)$$ (45) is an imbedding. If $`m=1`$, then $$𝒢_1=\frac{(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0})}{/J_{\eta _0}}=\frac{(j_{\eta _0})_{}(j_{\eta _0})^{}𝒪_{Y_{\eta _0}}_{𝒪_{Y_{\eta _0}}}(/J_{\eta _0})}{/J_{\eta _0}}=$$ $$=\frac{\underset{k}{\underset{}{lim}}𝒪(kY_{\eta _0+1})_{𝒪_{Y_{\eta _0}}}(/J_{\eta _0})}{/J_{\eta _0}}=\underset{k>0}{\underset{}{lim}}((/J_{\eta _0})_{𝒪_{Y_{\eta _0}}}(𝒪_{Y_{\eta _0}}(kY_{\eta _0+1})/𝒪_{Y_{\eta _0}}))\text{.}$$ Denote the sheaf $`_k=(/J_{\eta _0})_{𝒪_{Y_{\eta _0}}}(𝒪_{Y_{\eta _0}}(kY_{\eta _0+1})/𝒪_{Y_{\eta _0}})`$. By induction on $`k`$ let us show that the maps $$A_{_0(\sigma _1)}(_k)A_{_0(\sigma _1)\sigma _2}(_k)$$ (46) are imbeddings. If $`k=1`$, then the sheaf $`_1`$ is locally free on $`Y_{\eta _0+1}`$. Therefore in this case (46) follows from theorem 4 applied to $`Y_{\eta _0+1}`$. If $`k>1`$, then from the exact sequence $$0\frac{𝒪_{Y_{\eta _0}}((k1)Y_{\eta _0+1})}{𝒪_{Y_{\eta _0}}}\frac{𝒪_{Y_{\eta _0}}(kY_{\eta _0+1})}{𝒪_{Y_{\eta _0}}}\frac{𝒪_{Y_{\eta _0}}(kY_{\eta _0+1})}{𝒪_{Y_{\eta _0}}((k1)Y_{\eta _0+1})}0$$ it follows the commutative diagram: $$\begin{array}{ccccccc}0& A_{_0(\sigma _1)}(_{k1})& & A_{_0(\sigma _1)}(_k)& & A_{_0(\sigma _1)}(_k/_{k1})& 0\\ & \text{}& & \text{}& & \text{}& \\ 0& A_{_0(\sigma _1)\sigma _2}(_{k1})& & A_{_0(\sigma _1)\sigma _2}(_k)& & A_{_0(\sigma _1)\sigma _2}(_k/_{k1})& 0\text{.}\end{array}$$ The sheaf $`_k/_{k1}=(/J_{\eta _0})_{𝒪_{Y_{\eta _0}}}(𝒪_{Y_{\eta _0}}(kY_{\eta _0+1})/𝒪_{Y_{\eta _0}}((k1)Y_{\eta _0+1}))`$ is locally free on $`Y_{\eta _0+1}`$. Therefore the map $`\gamma `$ is injective by theorem 4 applied to $`Y_{\eta _0+1}`$. The map $`\alpha `$ is injective by induction hypothesis. Hence we have that the map $`\beta `$ is injective as well. Thus (46) is proved. After passage in (46) to the direct limit on $`k`$ we obtain (45) in the case $`m=1`$. Now let us show (45) in the case $`m>1`$. From the exact sequences: $$0\frac{J_{\eta _0}^{m1}}{J_{\eta _0}^m}/J_{\eta _0}^m/J_{\eta _0}^{m1}0$$ and $$0(j_{\eta _0})_{}(j_{\eta _0})^{}(\frac{J_{\eta _0}^{m1}}{J_{\eta _0}^m})(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m)(j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^{m1})0$$ it follows the exact sequence $$0\frac{(j_{\eta _0})_{}(j_{\eta _0})^{}(\frac{J_{\eta _0}^{m1}}{J_{\eta _0}^m})}{\frac{J_{\eta _0}^{m1}}{J_{\eta _0}^m}}𝒢_m𝒢_{m1}0\text{.}$$ Applying the exact functors $`A_{_0(\sigma _1)}`$ and $`A_{_0(\sigma _1)\sigma _2}`$, we obtain the diagram $$\begin{array}{ccccccc}0& A_{_0(\sigma _1)}(\frac{(j_{\eta _0})_{}(j_{\eta _0})^{}(\frac{J_{\eta _0}^{m1}}{J_{\eta _0}^m})}{\frac{J_{\eta _0}^{m1}}{J_{\eta _0}^m}})& & A_{_0(\sigma _1)}(𝒢_m)& & A_{_0(\sigma _1)}(𝒢_{m1})0& \\ & \text{}& & \text{}& & \text{}& \\ 0& A_{_0(\sigma _1)\sigma _2}(\frac{(j_{\eta _0})_{}(j_{\eta _0})^{}(\frac{J_{\eta _0}^{m1}}{J_{\eta _0}^m})}{\frac{J_{\eta _0}^{m1}}{J_{\eta _0}^m}})& & A_{_0(\sigma _1)\sigma _2}(𝒢_m)& & A_{_0(\sigma _1)\sigma _2}(𝒢_{m1})0\text{.}& \end{array}$$ The sheaf $`\frac{J_{\eta _0}^{m1}}{J_{\eta _0}^m}=\frac{J_{\eta _0}^{m1}}{J_{\eta _0}^m}_{Y_{\eta _0}}(/J_{\eta _0})`$ is locally free on $`Y_{\eta _0}`$. Therefore the map $`\alpha `$ is injective by the same reasons as in the case $`m=1`$. The map $`\gamma `$ is injective by the inductive hypothesis. Therefore the map $`\beta `$ is injective as well (this follows from non complicated diagram search). Thus we have shown (45). Now consider the following diagram. $$\begin{array}{ccccc}0A_{_0(\sigma _1)}(/J_{\eta _0}^m)& & A_{_0(\sigma _1)}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))& \stackrel{\psi }{}& A_{_0(\sigma _1)}(𝒢_m)0\\ \text{}& & \text{}& & \text{}\\ 0A_{_0(\sigma _1)\sigma _2}(/J_{\eta _0}^m)& & A_{_0(\sigma _1)\sigma _2}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))& & A_{_0(\sigma _1)\sigma _2}(𝒢_m)0\text{.}\end{array}$$ (Note also that here the map $`\theta `$ is injective. This follows from the fact that the map $`\beta `$ is injective by the above, and the statement $`\varphi `$ is injective by lemma 13.) Now let an element $$xA_{_0(\sigma _1)}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))\text{, \hspace{1em}but}$$ $$xA_{_0(\sigma _1)}(/J_{\eta _0}^m)\text{.}$$ Then since the map $`\beta `$ is injective, we have that $$\beta \psi (x)0\text{.}$$ (47) Now consider the diagram $$\begin{array}{ccccc}0A_{\sigma _2}(/J_{\eta _0}^m)& & A_{\sigma _2}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))& \stackrel{\psi _1}{}& A_{\sigma _2}(𝒢_m)0\\ \text{}& & \text{}& & \text{}\\ 0A_{_0(\sigma _1)\sigma _2}(/J_{\eta _0}^m)& & A_{_0(\sigma _1)\sigma _2}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))& & A_{_0(\sigma _1)\sigma _2}(𝒢_m)0\text{.}\end{array}$$ (Similarly to the previous reasonings we have that in this diagram all the vertical arrows are injective.) And if an element $$xA_{\sigma _2}(/J_{\eta _0}^m)\text{, \hspace{1em}then}$$ $$\beta _1\psi _1(x)=0\text{.}$$ (48) Now if $$xA_{_0(\sigma _1)}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))A_{\sigma _2}(/J_{\eta _0}^m)$$ (where the intersection is possible to be taken in $`A_{_0(\sigma _1)\sigma _2}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))`$, because all the appearing maps are injective), then from functoriality we have $$\beta \psi (x)=\beta _1\psi _1(x)\text{.}$$ Therefore, comparing this with (47) and (48), we obtain that in $`A_{_0(\sigma _1)\sigma _2}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))`$ is satisfied $$A_{_0(\sigma _1)}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))A_{\sigma _2}(/J_{\eta _0}^m)=A_{_0(\sigma _1)}(/J_{\eta _0}^m)A_{\sigma _2}(/J_{\eta _0}^m)\text{.}$$ Now from the definition of $`A_\sigma `$ we have that $$A_{\sigma _1}()=\underset{}{lim}_mA_{_0(\sigma _1)}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))\text{,}$$ from lemma 10 we have that $$A_{\sigma _2}()=\underset{}{lim}_mA_{\sigma _2}(/J_{\eta _0}^m)\text{.}$$ Therefore, $$A_{\sigma _1}()A_{\sigma _2}()=\left(\underset{}{lim}_mA_{_0(\sigma _1)}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))\right)\left(\underset{}{lim}_mA_{\sigma _2}(/J_{\eta _0}^m)\right)=$$ $$=\underset{}{lim}_m\left(A_{_0(\sigma _1)}((j_{\eta _0})_{}(j_{\eta _0})^{}(/J_{\eta _0}^m))A_{\sigma _2}(/J_{\eta _0}^m)\right)=$$ $$=\underset{}{lim}_m\left(A_{_0(\sigma _1)}(/J_{\eta _0}^m)A_{\sigma _2}(/J_{\eta _0}^m)\right)=$$ $$=\left(\underset{}{lim}_mA_{_0(\sigma _1)}(/J_{\eta _0}^m)\right)\left(\underset{}{lim}_mA_{\sigma _2}(/J_{\eta _0}^m)\right)=A_{_0(\sigma _1)}()A_{\sigma _2}()\text{.}$$ Acting further in this manner, i. e., eliminating the minmal number in the union of indices every time, we obtain that $$A_{\sigma _1}()A_{\sigma _2}()=A_{(i)}()A_\sigma ()\text{, \hspace{1em}where}$$ $`\sigma =(\zeta _0,\mathrm{})`$ and $`\zeta _0>i`$. But in this case by the reasonings, which is completely analogous to the above, we obtain at once that $$H^0(X,(j_i)_{}(j_i)^{}(/J_i^m))A_\sigma (/J_i^m)=$$ $$=H^0(X,/J_I^m)A_\sigma (/J_i^m)=H^0(X,/J_i^m)\text{.}$$ (Note that in contrast to the reasonings above with the functor $`A_\gamma `$, the functor $`H^0(X,)`$ is a left exact functor only. But the key diagram works in this case as well: $$\begin{array}{ccccccc}0& A_1& & A_2& & A_3& \\ & \text{}& & \text{}& & \text{}& \\ 0& B_1& & B_2& & B_3\text{.}& \end{array}$$ If the maps $`\alpha `$ and $`\gamma `$ are injective, then the map $`\beta `$ is injective as well.) Now $$A_i()A_\sigma ()=\left(\underset{}{lim}_mH^0(X,(j_i)_{}(j_i)^{}(/J_i^m))\right)\left(\underset{}{lim}_mA_\sigma (/J_i^m)\right)=$$ $$=\underset{}{lim}_m\left(H^0(X,(j_i)_{}(j_i)^{}(/J_i^m))A_\sigma (/J_i^m)\right)=\underset{}{lim}_mH^0(X,/J_i^m)=H^0(X,)\text{,}$$ where the last equality follows from lemma 12. Item 1 of theorem 5 is proved. Now let us show item 2 of the theorem. Consider a few cases. Case 1. $`\sigma _1\sigma _2ø`$, $`0\sigma _1`$, $`0\sigma _2`$. By lemma 10 we have that $$A_{\sigma _1}()=\underset{}{lim}_mA_{\sigma _1}(/J_1^m)\text{,}A_{\sigma _2}()=\underset{}{lim}_mA_{\sigma _2}(/J_1^m)\text{,}$$ $$A_{\sigma _1\sigma _2}()=\underset{}{lim}_mA_{\sigma _1\sigma _2}(/J_1^m)\text{.}$$ Let us show that for any $`m1`$ $$A_{\sigma _1\sigma _2}(/J_1^m)=A_{\sigma _1}(/J_1^m)A_{\sigma _2}(/J_1^m)\text{,}$$ (49) where the last intersection is regarded in $`A_{\sigma _1\sigma _2}(/J_1^m)`$. (By lemma 13 we can imbed these groups there.) Let us prove (49) by induction on $`m`$. Let $`m=1`$. In this case $`/J_1`$ is a locally free sheaf on $`Y_1`$. And equality (49) turns into the analogous equality (49) on $`Y_1`$. The scheme $`Y_1`$ has lesser dimension than dimension of $`X`$. Applying induction on dimension of scheme, we can suppose that theorem 5 is already true for schemes of lesser dimension. (For schemes of dimension 1 theorem 5 follows at once from theorem 4 and theorem 3 by trivial reasons.) Therefore (49) is true when $`m=1`$. Let $`m>1`$. Then the exact sequence $$0J_1^{m1}/J_1^m/J_1^m/J_1^{m1}0$$ induces the following commutative diagram: $$\begin{array}{ccccc}0A_{\sigma _1\sigma _2}(J_1^{m1}/J_1^m)& & A_{\sigma _1\sigma _2}(/J_1^m)& & A_{\sigma _1\sigma _2}(/J_1^{m1})0\\ \text{}& & \text{}& & \text{}\\ 0A_{\sigma _1}(\frac{J_1^{m1}}{J_1^m})A_{\sigma _2}(\frac{J_1^{m1}}{J_1^m})& & H_m& & \frac{H_m}{A_{\sigma _1}\left(\frac{J_1^{m1}}{J_1^m}\right)A_{\sigma _2}\left(\frac{J_1^{m1}}{J_1^m}\right)}0\text{,}\end{array}$$ where $`H_m=A_{\sigma _1}(/J_1^m)A_{\sigma _2}(/J_1^m)\text{.}`$ From $`\sigma _1\sigma _2ø`$ it follows that the functor $`A_{\sigma _1\sigma _2}`$ is exact. Therefore the upper row of the diagram is exact. There is the natural map $`\theta `$: $$\frac{H_m}{A_{\sigma _1}\left(\frac{J_1^{m1}}{J_1^m}\right)A_{\sigma _2}\left(\frac{J_1^{m1}}{J_1^m}\right)}H_{m1}\text{.}$$ And from the exact sequences $$0A_{\sigma _1}(J_1^{m1}/J_1^m)A_{\sigma _1}(/J_1^m)A_{\sigma _1}(/J_1^{m1})$$ $$0A_{\sigma _2}(J_1^{m1}/J_1^m)A_{\sigma _2}(/J_1^m)A_{\sigma _2}(/J_1^{m1})$$ it follows at once that the map $`\theta `$ is an imbedding. Besides, the map $`\theta \gamma `$ is the natural map from $`A_{\sigma _1\sigma _2}(/J_1^{m1})`$ to $`H_{m1}`$; and, consequently, by the induction hypothesis it is possible to suppose that $`\theta \gamma `$ is an isomorphism. From the last two facts we obtain at once that $`\gamma `$ is an isomorphism. Since the sheaf $`J_1^{m1}/J_1^m=J_1^{m1}/J_1^m_{𝒪_X}`$ is locally free on $`Y_1`$ and $`dimY_1<dimX`$, we have that the map $`\alpha `$ is an isomorphism as well. Therefore from this commutative diagram it follows that the map $`\beta `$ is an isomorphism as well. Thus, equality (49) is proved. Now passage in (49) to the projective limit on $`m`$ concludes the proof of case 1. Case 2. $`0\sigma _1`$, $`0\sigma _2`$ (or vice versa), $`\sigma _1\sigma _2ø`$. Now by the analogous reasonings, as in the proof of item 1 of this theorem, we obtain at once the following $$A_{\sigma _1}()A_{\sigma _2}()=A_{_0(\sigma _1)}()A_{\sigma _2}()\text{;}$$ that reduces this case to the case 1, analyzed above. Case 3. $`0\sigma _1`$, $`0\sigma _2`$. Then $$A_{\sigma _1}()=\underset{}{lim}_kA_{_0(\sigma _1)}((kY_1)),A_{\sigma _2}()=\underset{}{lim}_kA_{_0(\sigma _2)}((kY_1))$$ Now from case 1 (or if $`_0(\sigma _1)_0(\sigma _2)=ø`$, then from item 1 of this theorem) we have that $$A_{_0(\sigma _1)}((kY_1))A_{_0(\sigma _2)}((kY_1))=A_{_0(\sigma _1)_0(\sigma _2)}((kY_1))\text{.}$$ (Here $`A_ø()=H^0(X,)`$). Therefore, $$A_{\sigma _1}()A_{\sigma _2}()=\left(\underset{}{lim}_kA_{_0(\sigma _1)}((kY_1))\right)\left(\underset{}{lim}_kA_{_0(\sigma _2)}((kY_1))\right)=$$ $$=\underset{}{lim}_k\left(A_{_0(\sigma _1)}((kY_1))A_{_0(\sigma _2)}((kY_1))\right)=\underset{}{lim}_kA_{_0(\sigma _1)_0(\sigma _2)}((kY_1))=A_{\sigma _1\sigma _2}()\text{.}$$ Theorem 5 is proved. In the sequel we shall assume that all the conditions of theorem 4 are satisfied, and a field $`k`$ is the field of definition of the scheme $`X`$. Also, let us assume that $`Y_n=x`$, where $`x`$ is a $`k`$-rational point on $`X`$ which is smooth on any $`Y_i`$ ($`0in`$). Let us choose and fix local parameters $`z_1,\mathrm{},z_n\widehat{𝒪}_{x,X}`$ such that $`z_{ni+1}\|_{Y_{i1}}=0`$ is a local equation of the divisor $`Y_i`$ in the formal neighbourhood of the point $`x`$ on the scheme $`Y_{i1}`$ ($`1in`$). Let $``$ be a rank 1 locally free sheaf on $`X`$. Fix a trivialization $`e_x`$ of the sheaf $``$ in the formal neghbourhood of the point $`x`$ on $`X`$. Now the done choice of local parameters and trivialization makes possible to identify $`A_{(0,1,\mathrm{},n)}()`$ with the $`n`$-dimensional local field $`k((z_1))\mathrm{}(((z_n))`$. Moreover, let us fix a collection of integers $`0j_1\mathrm{}j_kn1`$. Define $`\sigma S_{nk}`$ as the set $`\{i:\mathrm{\hspace{0.25em}0}in,ij_1,\mathrm{},ij_k\}`$. By theorem 4 we have the natural imbedding $`A_\sigma ()A_{(0,1,\mathrm{},n)}()`$. And under identifying of $`A_{(0,1,\mathrm{},n)}()`$ with the field $`k((z_1))\mathrm{}((z_n))`$ the space $`A_\sigma ()`$ converts to the following $`k`$-subspace in $`k((z_1))\mathrm{}((z_n))`$: $$\{a_{i_1,\mathrm{},i_n}z_1^{i_1}\mathrm{}z_n^{i_n}:a_{i_1,\mathrm{},i_n}k,i_{nj_1}0,i_{nj_2}0,\mathrm{},i_{nj_k}0\}\text{.}$$ (50) Thus, from theorem 5 we obtain that for determination of the images of $`A_\sigma ()`$ in $`k((z_1))\mathrm{}((z_n))`$ (for any $`\sigma S`$) it suffices to know only one image of $`A_{(0,1,\mathrm{},n1)}`$ in $`k((z_1))\mathrm{}((z_n))`$. (All the others are obtained by intersection of the image of $`A_{(0,1,\mathrm{},n1)}`$ in $`k((z_1))\mathrm{}((z_n))`$ with the standard subspaces (50) in $`k((z_1))\mathrm{}((z_n))`$.) It is clear that these reasonings is generalized at once to locally free sheaves $``$ of rank $`r`$ and spaces $`k((z_1))\mathrm{}((z_n))^r`$. Denote the described map $$(X,Y_1,\mathrm{},Y_n,(z_1,\mathrm{},z_n),,e_x)A_{(0,1,\mathrm{},n1)}()A_{(0,1,\mathrm{},n1,n)}()\stackrel{e_x}{}$$ $$\stackrel{e_x}{}A_{(0,1,\mathrm{},n)}(\widehat{𝒪}_{x,X})\stackrel{z_1,\mathrm{},z_n}{}k((z_1))\mathrm{}((z_n))^r$$ by $`\mathrm{\Psi }_r`$. Definition. $$\begin{array}{ccc}_n\hfill & \stackrel{\mathrm{def}}{=}& \{X,(Y_1,\mathrm{},Y_n),(z_1,\mathrm{},z_n),,e_{Y_n}\}\hfill \\ X\hfill & & \text{a projective equidimensional Cohen -Macaulay scheme}\hfill \\ & & \text{of dimension }n\text{ over a field }k\hfill \\ X=Y_0Y_1\mathrm{}Y_n\hfill & & \text{a flag of closed subschemes such that }Y_i\text{ is an ample}\hfill \\ & & \text{Cartier divisor on the scheme }Y_{i1}\text{ }(1in)\hfill \\ Y_n\hfill & & \text{a smooth }k\text{-rational point on all }Y_i\text{ }(0in)\hfill \\ z_1,\mathrm{},z_n\hfill & & \text{formal local parameters in the point }Y_n\hfill \\ & & \text{such that }\left(z_{ni+1}\|_{Y_{i1}}=0\right)=Y_i\text{ in the formal}\hfill \\ & & \text{neighbourhood of the point }Y_n\text{ on the scheme }Y_{i1}\hfill \\ \hfill & & \text{a rank }r\text{ vector bundle on }X\hfill \\ e_{Y_n}\hfill & & \text{a trivialization of }\text{ in the formal neighbourhood}\hfill \\ & & \text{of the point }Y_n\text{ on }X\hfill \end{array}$$ In the field $`K=k((z_1))\mathrm{}((z_n))`$ we have the following filtration $$K(m)=z_n^mk((z_1))\mathrm{}((z_{n1}))[[z_n]]\text{.}$$ Let $`K`$-space $`V=K^r`$, and let the filtartion $`V(m)=K(m)^r`$. ###### Theorem 6 There exists a canonical map $$\mathrm{\Phi }_n:_n\{k\text{-vector subspaces}BK,WV\}$$ such that 1. from the subspace $`BK`$ is uniquely reconstructed the complex $`A(𝒪_X)`$, which calculates cohomology of the sheaf $`𝒪_X`$ on $`X`$; 2. from the subspace $`WV`$ is uniquely reconstructed the complex $`A()`$, which calculates cohomology of the sheaf $``$ on $`X`$; 3. if $`(B,W)\mathrm{Im}\mathrm{\Phi }_n`$, then $`BBB`$, $`BWW`$; 4. for all $`m`$ the map $$\{Y_1,(Y_2,\mathrm{},Y_n),(z_1,\mathrm{},z_{n1})\|_{Y_1},(mY_1)\|_{Y_1},e_{Y_n}(m)\|_{Y_1}\}$$ $$\{\frac{BK(m)}{BK(m+1)}\frac{K(m)}{K(m+1)}=k((z_1))\mathrm{}((z_{n1}))\text{,}$$ $$\frac{WV(m)}{WV(m+1)}\frac{V(m)}{V(m+1)}=k((z_1))\mathrm{}((z_{n1}))^r\}$$ coincides with the map $`\mathrm{\Phi }_{n1}`$; 5. If $`q,q^{}_n`$ and $`\mathrm{\Phi }_n(q)=\mathrm{\Phi }_n(q^{})`$, then $`q`$ is isomorphic to $`q^{}`$. Proof. If $$q=\{X,(Y_1,\mathrm{},Y_n),(z_1,\mathrm{},z_n),,e_{Y_n}\}_n\text{,}$$ then to define the map $`\mathrm{\Phi }_n`$ we put $$B=\mathrm{\Psi }_1(X,Y_1,\mathrm{},Y_n,(z_1,\mathrm{},z_n),𝒪_X,id)\text{,}$$ $$W=\mathrm{\Psi }_r(X,Y_1,\mathrm{},Y_n,(z_1,\mathrm{},z_n),,e_{Y_n})\text{,}$$ $$\mathrm{\Phi }_n(q)=\{B,W\}\text{.}$$ Now items 1-4 of this theorem follows from theorems 3, 4, 5 and reasonings above about the map $`\mathrm{\Psi }`$, and also for item 4 is needed the fact that $`(j_0)_{}(j_0)^{}=\underset{m}{\underset{}{lim}}(mY_1)`$. Let us show item 5. Intersecting $`B`$ with the standard subspaces (50), we can uniquely reconstruct the algebra $`A_{(0)}(𝒪_X)K`$. Similarly, from $`W`$ we can reconstruct the $`k`$-subspace $`A_{(0)}()V`$. Then $$XY_1=SpecA_{(0)}(𝒪_X)$$ $$X=Proj\left(\underset{m0}{}(A_{(0)}(𝒪_X)K(m))\right)$$ (51) $$=Proj\left(\underset{m0}{}(A_{(0)}()V(m))\right)\text{,}$$ where the last equalities follow from the following statement (see \[17, lemma 7\]): if $`X`$ is a projective scheme over a field, $``$ is a coherent sheaf on $`X`$, and $`C`$ is an ample divisor on $`X`$, then $`XProj(S)`$, $`Proj(F)`$, where $`S=_{m0}H^0(X,𝒪_X(mC))`$, $`F=_{m0}H^0(X,(mC))`$. Besides, the image under the imbedding of $$\underset{m0}{}(A_{(0)}(𝒪_X)K(m+1))\text{into}\underset{m0}{}(A_{(0)}(𝒪_X)K(m))$$ is the homogeneous ideal determining the subscheme $`Y_1`$ in $`X`$. Now, using item 4 of this theorem, we can reconstruct all the geometrical data on the subscheme $`Y_1`$ in the analogous way, and, further, by induction we can reconstruct all the data from $`q`$ up to an isomorphism. Theorem 6 is proved. ###### Remark 9 Note that $`(nY_1)\|_{Y_1}=_{𝒪_X}𝒪(nY_1)\|_{Y_1}`$ and the sheaf $$𝒪(nY_1)\|_{Y_1}=𝒪(nY_1)/𝒪((n1)Y_1)=𝒩_{Y_1/X}^n\text{,}$$ where the bundle $`𝒩_{Y_1/X}`$ coincides with the normal budle of $`Y_1`$ in $`X`$ in some cases (for example, if $`X`$ and $`Y_1`$ are smooth). ###### Remark 10 From (51) and absence of divisors of zero in the field $`K`$ it follows at once that the schemes $`X,Y_1,\mathrm{},Y_n`$, satisfied the conditions of the definiton $`_n`$, are always integral schemes. ###### Remark 11 $`\mathrm{\Phi }_1`$ is a variant of the Krichever correspondence for curves (see , ). Besides, any integral noetherian scheme of dimension $`1`$ is a Cohen-Macaulay scheme. $`\mathrm{\Phi }_2`$ coincides with the map, constructed in . Note that any normal noetherian scheme of dimension $`2`$ is a Cohen-Macaulay scheme (see \[9, ch. II, th. 8.22A\]). Also in is analyzed an example of the Krichever map for $`X=𝐏_2`$. ###### Remark 12 In is discussed the problem of change of locally free sheaves to torsion free sheaves. Let $`X`$ be a smooth projective surface with a flag of irreducible subvarieties $`Y_1Y_2`$ such that $`Y_1`$ is an ample divisor on $`X`$, and $`Y_2`$ is a point. Let $`𝒢`$ be any torsion free sheaf on $`X`$. Then we have the imbedding $$0𝒢𝒢^{}\text{.}$$ (52) The sheaf $`𝒢^{}`$ is reflexive; and since $`dimX=2`$, we have that $`𝒢^{}`$ is a locally free sheaf. Now applying the exact functors $`A_\sigma `$ and the left exact functor $`H^0(X,)`$ to sequence (52), we obtain at once from the obtained sequences and theorem 4 for the sheaf $`𝒢^{}`$ that theorem 4 is true for the sheaf $`𝒢`$ as well. But theorem 5 is no longer true for torsion free sheaves on $`X`$. Indeed, let $`𝒢=m_Q`$, where $`m_Q`$ is the ideal sheaf of a point $`QX`$. Then, applying the exact functors $`A_\sigma `$ to the exact sequence $$0m_Q𝒪_Xk_Q0$$ we obtain at once that * if $`QY_1`$, then $$A_{(0)}(𝒢)A_{(01)}(𝒢)A_{(02)}(𝒢)\text{,}$$ because $`A_{(01)}(𝒢)A_{(01)}(𝒪_X)`$, but $`A_{(01)}(𝒢)=A_{(01)}(𝒪_X)`$, $`A_{(02)}(𝒢)=A_{(02)}(𝒪_X)`$; * if $`QY_1`$, $`QY_2`$, then $$A_{(1)}(𝒢)A_{(01)}(𝒢)A_{(12)}(𝒢)\text{,}$$ because $`A_{(1)}(𝒢)A_{(1)}(𝒪_X)`$, but $`A_{(01)}(𝒢)=A_{(01)}(𝒪_X)`$, $`A_{(12)}(𝒢)=A_{(12)}(𝒪_X)`$. Steklov Mathematical Institute E-mail : d\_osipov@mi.ras.ru E-mail : d\_osipov@mail.ru
warning/0003/cond-mat0003078.html	ar5iv	text	# The induced charge in a Fröhlich polaron: Sum rule and spatial extent. ## I Introduction An electron added to an insulating polar crystal forms a quasiparticle called dielectric polaron after Fröhlich. This has been recognized as a fundamental field theoretical problem. More recently a variety of novel materials have emerged which present interesting properties when doped away from an insulating phase, like colossal magnetoresistent manganites and high temperature superconducting cuprates. The fact that these are polar crystals has produced a renewed interest in Fröhlich polaron problem. Roughly speaking a dielectric polaron is composed of an electron and the (opposite) charge that it induces in the lattice. Electron and induced charge attract each other so that for the electron to move it has to drag the induced charge resulting in an increase of the quasiparticle mass. In this work we study correlation functions which measure the magnitude and spatial extent of the induced charge and associated polarization field. We derive a rigorous sum rule which states that the total induced charge equals the charge induced by a classical (static) electron and it is independent of temperature. This result, well known within perturbative and variational approaches , is proven here to be exact. The large distance behavior of the electric field is determined by this sum rule. In addition we discuss the short distance and the high temperature asymptotic limit of these quantities. These results provide constrains to approximations on the polaron problem. Real space path integral Monte Carlo (PIMC) method is used to evaluate correlation functions. These are compared with Feynman’s variational approximation (FVA) and analytical results in weak and strong coupling. We find that the polaron radius is determined at low temperatures by the electron-phonon coupling $`\alpha `$ alone while at high temperatures it is proportional to the de Broglie thermal wave length ($`\lambda _T=\mathrm{}\sqrt{2\pi \beta /m}`$ with $`\beta `$ the inverse temperature) and becomes independent of coupling. We define a polaron crossover temperature $`T^{}(\alpha )`$. Although the electron induces a temperature independent charge in the lattice the induced charge hinders the electron motion only below $`T^{}(\alpha )`$. At high temperatures thermal effects wash out the hindering effect of the induced charge but, we remark again, not the induced charge itself. ## II Analytical Results Our starting point is the effective action for the Fröhlich polaron problem, after the phonons have been eliminated with the usual path integral techniques, $`S=S_0+S_I`$ with $`S_0[𝐱]`$ $`={\displaystyle _0^\mathrm{}\beta }𝑑\tau {\displaystyle \frac{1}{2}}m\dot{x}^2`$ (1) $`S_I[𝐱]`$ $`={\displaystyle \frac{\alpha }{2\sqrt{2}}}{\displaystyle \frac{(\mathrm{}\omega _L)^{3/2}}{m^{1/2}}}{\displaystyle _0^\mathrm{}\beta }𝑑\tau _1{\displaystyle _0^\mathrm{}\beta }𝑑\tau _2{\displaystyle \frac{D(\tau _1\tau _2)}{\|𝐱(\tau _1)𝐱(\tau _2)\|}}`$ (2) here $`m`$ is the electron mass, $`\alpha =e^2m^{1/2}/\mathrm{}\overline{ϵ}(2\mathrm{}\omega _L)^{1/2}`$, is the coupling constant with $`\omega _L`$ the phonon frequency, $`1/\overline{ϵ}=1/ϵ_{\mathrm{}}1/ϵ_0`$ and $`ϵ_{\mathrm{}}`$ ($`ϵ_0`$) is the high (zero) frequency dielectric constant. $$D(\tau )=\frac{\mathrm{exp}(\omega _L\|\tau \|)+\mathrm{exp}[(\mathrm{}\beta \|\tau \|)\omega _L]}{\mathrm{exp}(\mathrm{}\beta \omega _L)1}$$ (3) is the phonon propagator. We are interested in the correlation function between the electron charge density $`e\stackrel{~}{n}(𝐫)`$ and the charge induced in the lattice $`e\stackrel{~}{n_i}(𝐫)`$ normalized to the probability density to find an electron at a given point i.e. the inverse of the volume $`V`$. Dropping the charges this is defined as, $$\stackrel{~}{n}(0)\stackrel{~}{n_i}(𝐫)/V^1g(r)/\overline{ϵ}$$ (4) Averages are defined as path integrals weighted by $`S`$ $$\mathrm{}=\frac{𝒟𝐱e^{S/\mathrm{}}(\mathrm{})}{𝒟𝐱e^{S/\mathrm{}}}$$ (5) where the paths entering in Eq. (5) depart and arrive at the same point. Further integration over such a point is not performed and this assigns to Eq. (5) the meaning of a constrained average with $`𝐱(0)=𝐱(\beta )=0`$. Those averages are however equivalent to the unconstrained ones because of translational invariance. We used the spherical symmetry of the problem to define $`g(r)`$, and we divided by $`\overline{ϵ}`$ on the right hand side of Eq. (4) for later convenience. Another quantity of interest is the induced lattice polarization $$𝐏(𝐫)\stackrel{~}{n}(0)\stackrel{~}{𝐏}(𝐫)/V^1$$ (6) $`\stackrel{~}{𝐏}(𝐫)=e\stackrel{~}{n_i}(𝐫)`$ is the density of polarization operator. The correlation function in Eq. (6) is related to the induced electrostatic potential \[$`^2V(𝐫)=4\pi eg(r)/\overline{ϵ}`$\] considered in Ref. and which will not be discussed here. We stress that these quantities have the meaning of correlation functions measuring average induced charge, polarization and induced potential at a distance $`r`$ from the electron position. The charge density operator for the phonons is, $$e\stackrel{~}{n_i}(𝐫)=\sqrt{\frac{\mathrm{}\omega _L}{4\pi V\overline{ϵ}}}\underset{𝐤}{}k\stackrel{~}{Q}_𝐤e^{i𝐤.𝐫}.$$ (7) where $`\stackrel{~}{Q}_𝐤`$ is the dimensionless displacement for momentum $`𝐤`$ phonons. Inserting Eq. (7) in Eq. (4) we obtain an equation for $`g(r)`$ as a function of the density displacement correlation function $`\stackrel{~}{n}(0)\stackrel{~}{Q}_𝐤`$. The phonon variables can be traced out by standard methods. We have found that it is possible to give an exact expression for the correlation functions in terms of path integrals weighted by the effective electronic action of Eq. (1). We find for the density-induced-density correlation function $$g(r)=_0^\beta 𝑑\tau U(\tau )\delta [𝐫𝐱(\tau )]$$ (8) and for the polarization field $$𝐏(𝐫)=\frac{e}{4\pi \overline{ϵ}}_0^\beta 𝑑\tau U(\tau )\frac{\widehat{𝐫}}{\|𝐫𝐱(\tau )\|^2}$$ (9) where $`\widehat{𝐫}𝐫/r`$ and $$U(\tau )=\mathrm{}\omega _L\frac{\mathrm{sinh}(\omega _L\tau )+\mathrm{sinh}[\omega _L(\mathrm{}\beta \tau )]}{2\mathrm{tanh}(\beta \mathrm{}\omega _L/2)\mathrm{sinh}(\beta \mathrm{}\omega _L)}.$$ (10) Within FVA, the variational quadratic action can be exploited in Eqs. (8),(9) to analytically perform the averages and to recover the results of Refs. and . Eqs. (8),(9) have a simple physical interpretation. The induced charge can be seen as “distributed” along the electron path with weight $`U(\tau )`$. The polarization is the superposition of polarizations associated with those elementary induced charges. Eq. (8) can be integrated in the whole space using the properties of the Dirac’s $`\delta `$ function. Since $`𝑑\tau U(\tau )=1`$, we conclude that $`g(r)`$ is normalized to one. The total induced charge $`q`$ is computed by integrating the density-induced density correlation function in Eq. (4): $$q=e_0^{\mathrm{}}𝑑r4\pi r^2g(r)/\overline{ϵ}=\frac{e}{\overline{ϵ}}.$$ (11) which completes the proof of the sum rule. The total induced charge amounts to the charge the electron would induce if it where a static classical particle. In other words there are no quantum corrections to the total induced charge. An alternative derivation can be work out following Quémerais. From the time derivative $`i\mathrm{}\ddot{\stackrel{~}{𝐏}}=[\dot{\stackrel{~}{𝐏}},\stackrel{~}{H}]`$ one obtains $$\ddot{\stackrel{~}{𝐏}}=\omega _L^2\left(\frac{1}{4\pi \overline{ϵ}}\stackrel{~}{𝐃}\stackrel{~}{𝐏}\right)$$ (12) Here $`\stackrel{~}{H}`$ is the Hamiltonian, $`\stackrel{~}{𝐃}`$ is the electric displacement operator due to the electron ($`.\stackrel{~}{𝐃}=4\pi e\stackrel{~}{n}`$) and $`i\mathrm{}\dot{\stackrel{~}{𝐏}}=[\stackrel{~}{𝐏},\stackrel{~}{H}]`$. Taking the divergence we obtain a relation for the charge operators: $$.\ddot{\stackrel{~}{𝐏}}(𝐫)=\omega _L^2e\left(\frac{\stackrel{~}{n}(𝐫)}{\overline{ϵ}}+\stackrel{~}{n_i}(𝐫)\right)$$ (13) We can integrate this expression in all space and take the thermodynamic average. The left hand side is proportional to the average of the net force felt by the lattice at the boundary of the system which should vanish at equilibrium. The right hand side gives Eq. (11). Eq. (11) shows that the induced charge is independent of temperature. This contradicts the naive argument that all polaron effects should disappear at high temperatures. To understand this behavior one can do an analogy with the behavior of an harmonic oscillator in an external field. In that case, because of harmonicity, one gets a displacement which is temperature independent. Here roughly speaking the harmonic oscillator represents the phonon coordinates and the external field is the field produced by the electron on the phonons. The induced charge is a measure of how much the ions displace from their bare equilibrium positions in the presence of the electron. As for the single harmonic oscillator, this “displacement” is independent of temperature. Only anharmonicities can make the induced charge temperature dependent. Using the sum rule it is easy to see that at distances much larger than the polaron radius, as defined below, the polarization field goes as $`𝐏(𝐫)=e\widehat{𝐫}/(4\pi \overline{ϵ}r^2)`$. Clearly the distortion produced by the electron is long range, a fact that is not always recognized in the literature. The total electric field (always in the sense of a correlation function) is given by $`𝐄=𝐃4\pi 𝐏`$ where we should include in $`𝐃=e\widehat{𝐫}/(ϵ_{\mathrm{}}r^2)`$ the high frequency screening. At long distances we have $`𝐄(𝐫)=e\widehat{𝐫}/(ϵ_0r^2)`$ which means that the electric field generated by the electron gets screened by the static dielectric constant. This is generally expected but to the best of our knowledge has never been proven for all coupling and temperatures. Now we discuss the short distance behavior. At distances much smaller than the polaron radius we expect that the effect of the interaction becomes irrelevant in the functional integrals. This is because the latter are dominated by electron paths with short wave length or equivalently high kinetic energy. We can then replace the total action by the free electron action in Eqs. (8),(9). We obtain the asymptotic result: $$\underset{r0}{lim}𝐏(𝐫)=\frac{e}{8\pi \overline{ϵ}l^2\mathrm{tanh}(\beta \mathrm{}\omega _L/2)}\widehat{𝐫}.$$ (14) where $`l=\sqrt{\mathrm{}/2m\omega _L}`$ is the harmonic oscillator characteristic length. Using the same argument we obtain that $`g(r)r^1`$ for $`r0`$ and the proportionality coefficient can also be obtained with the same method. The latter behavior has been obtained in Ref. within the FVA. These results coincide with lowest order perturbation theory. At high temperatures we also expect that the effect of the interaction becomes irrelevant and so we can replace again the total action by the free electron action in the functional integrals. The high temperature asymptotic result for $`g(r)`$ is $$4\pi r^2g(r)=\frac{2r}{l^2\beta \mathrm{}\omega _L}\mathrm{exp}\left(\frac{r^2}{l^2\beta \mathrm{}\omega _L}\right)\beta \mathrm{}\omega _L<<1.$$ (15) This result has also been obtained in Ref. within FVA. We remark that although the density-induced density correlation function does not vanish for large temperatures the polaron effective mass tends to the bare electron mass. ## III Numerical Results Now we discuss the spatial extent of the induced charge at general couplings and temperatures. We have evaluated averages in Eq. (8) using PIMC. Eqs. (8),(9) being expressed in real space rather that in Fourier components are more suitable for this purpose. We have performed Metropolis PIMC calculations within the imaginary time discretization scheme. In order to regularize the attractive divergence of the retarded action at short distance, and to improve the convergence with the number of the imaginary time discretization points, we have developed a preaveraging procedure similar to the one used for local actions. Details of the method and more extensive results will be given in a separate publication. Here we just say that the results for the $`g(r)`$ are well converged as checked by doubling the number of imaginary time slices. Previous MC studies of the Fröhlich polaron were limited to small ($`\alpha 4`$) or intermediate ($`\alpha 7`$) couplings and were focused to the calculation of the ground state energy and effective mass. Following Ref. we have also computed $`g(r)`$ in the FVA, i.e. using Feynman’s quadratic action to evaluate the average appearing in Eq. (8). In Fig. 1 we show $`4\pi r^2g(r)`$ for weak, intermediate and strong coupling. We also show the weak coupling result obtained by perturbation theory at $`T=0`$ and the strong coupling result in the Landau-Pekar approximation. For all couplings the correlation function decays exponentially with distance as expected for a polaron. The area under the curves is one according to the sum rule Eq. (11). The short distance asymptotic behavior Eq. (14) is exactly satisfied in FVA and is also satisfied in the PIMC within the numerical error. We define the polaron radius $`r_m`$ as the distance at which $`4\pi r^2g(r)`$ is maximum. As the coupling increases the polaron shrinks indicating the progressively more localized nature. In Fig. 2 we show the temperature dependence of $`4\pi r^2g(r)`$ at intermediate coupling. In Figs. 1,2 we see that, apart from a small overestimate of the polaron radius, FVA reproduces fairly well the PIMC data. We can then safely use the FVA to discuss the spatial properties of the polaron. Notice that the good agreement found between PIMC and FVA is not obvious since in principle only the free energy is expected to be accurate for the latter. Contrary to the naive expectation the effect of temperature is to shrink the polaron (Fig. 2). Our physical explanation is the following: At low temperatures phonons relax shrinking the spatial extent of the electron till the increase in electron kinetic energy balances the gain in electron-phonon interaction energy. At high temperatures, however, a typical electron has energy $`\overline{E}=3/2\beta `$ and momentum $`\mathrm{}\overline{k}=\sqrt{3m/\beta }`$. One can construct a wave packet of width $`\mathrm{\Delta }k`$ in momentum space using plane waves with higher and smaller energy without affecting the electron internal energy. The biggest $`\mathrm{\Delta }k`$ which will not affect the electron internal energy is of order of $`\overline{k}`$ itself. One can then achieve a localization of the electron of order $`1/\overline{k}\mathrm{}/\sqrt{3m/\beta }\lambda _T`$. It follows that the phonons can relax at practically no cost till the polaron radius stabilizes at a value of this order. In fact at high temperatures the asymptotic value of the polaron radius can be obtained from Eq. (15): $`lim_{\beta 0}r_m=l\sqrt{\beta \mathrm{}\omega _L}/2=0.2\lambda _T`$. This scaling has been found by Sethia et al. for the mean square displacement of the electron in imaginary time. Notice that the polaron radius becomes independent of the coupling. The same behavior of $`r_m`$ has been obtained in Ref. within FVA. The authors of Ref. ascribe the high temperature behavior of the polaron radius to the increased fluctuations of the phonon field. We conversely think that the thermal fluctuations of the electron are responsible of the high temperature behavior of the polaron radius and the phonon field acts only as a probe of the intrinsic thermal length of the electron as explained above. To characterize the temperature and coupling dependence of the polaron size we plot in Fig. 3 the polaron radius as a function of coupling for different inverse temperatures. We see that when the temperature is such that $`0.2\lambda _T>l`$ (low temperatures, see the curves for $`\beta \mathrm{}\omega _L=20,\mathrm{}`$), the polaron radius exhibits little temperature dependence at all couplings and simply interpolates between the weak coupling polaron radius $`r_m=l`$ and the Landau-Pekar polaron radius $`r_m=3l\sqrt{\pi /2}/\alpha `$. When $`0.2\lambda _T<l`$ two different regimes occur. At small coupling the polaron radius tends to saturate at $`0.2\lambda _T`$ (horizontal arrow) whereas at high couplings one recovers the low temperature polaron radius $`r_m(\alpha ,T=0)`$. We can define a crossover line when these two lengths are of equal magnitude so that the crossover temperature as a function of $`\alpha `$ is given by the equation $`r_m(\alpha ,T=0)=0.2\lambda _T(T^{})`$. In the FVA approximation we find the crossover temperature in energy units to be $`T^{}(\alpha )=0.15\mathrm{}\omega _Lv(\alpha ,T=0)`$ with $`v`$ the usual Feynman variational parameter. For $`\alpha 0`$, $`T^{}`$ goes to a constant of order of $`0.5\mathrm{}\omega _L`$ whereas for large $`\alpha `$ it increases quadratically with $`\alpha `$. In the low temperature regime ($`T<T^{}`$) the polaron radius becomes almost temperature independent and is determined by the coupling. The high temperature regime $`T>T^{}`$ is characterized by a polaron radius which is independent of the coupling and is determined by the temperature alone $`(r_m=0.2\lambda _T)`$. ## IV Conclusions We have studied the charge induced in the lattice by an electron in Fröhlich polaron problem and the associated polarization field. We have derived relations which express the charge-induced density correlation function in terms of path integral involving only the electronic degree of freedom which are suitable to be evaluated by PIMC method. A rigorous sum rule was derived that determines the total induced charge and the long distance behavior of the polarization field. We give also the asymptotic limits of these quantities at short distances and at high temperatures. We have compared results obtained using FVA through the lines of Refs. and those obtained by PIMC method. To the best of our knowledge this is the first PIMC computation of real space correlation functions in Fröhlich model. From the spatial dependence of the induced charge we obtained a polaron radius. The polaron radius is determined by the coupling at low temperatures and by the thermal wave length at high temperatures with a crossover temperature that we evaluated in FVA. At high temperature a polaron with small radius and small effective mass is achieved. These results are not in contradiction because the small radius at high temperatures is a thermal effect of the electron and it is not related to the lattice response. The lattice acts only as a probe of the intrinsic electron localization radius namely $`\lambda _T`$. Obviously this small radius polaron has nothing to do with the Holstein zero-temperature small polaron which induces almost local lattice displacements and moves coherently with a large effective mass. The PIMC polaron radius is always smaller than the FVA calculation in the range of coupling and temperature studied. This effect is more pronounced at intermediate couplings. The overall temperature dependence agrees with the findings of Ref. however our physical interpretation is different. ###### Acknowledgements. We acknowledge the third anonymous referee to bring us the attention on Ref. (which suggested us the second proof of the induced-density sum rule) and Ref. . We acknowledge useful discussions and suggestions from S. Fratini and J. T. Titantah. J.L. thanks P. Calvani’s group for hospitality during this work. We acknowledge partial support from the MURST 1997 matching funds program no.9702265437.
warning/0003/nlin0003008.html	ar5iv	text	# Signatures of quantum integrability and nonintegrability in the spectral properties of finite Hamiltonian matrices ## I Introduction An autonomous classical Hamiltonian system with two degrees of freedom, specified by some analytic function $`H(p_1,q_1;p_2,q_2)`$ of canonical coordinates, is either integrable or nonintegrable – tertium non datur. If a second integral of the motion can be found, i.e. an analytic function $`I(p_1,q_1;p_2,q_2)`$ which is functionally independent of $`H`$ and satisfies $`dI/dt=\{H,I\}=0`$, the system is proven integrable. If chaotic trajectories can be detected in the phase flow, the system is demonstrably nonintegrable. Although it may happen that neither evidence can be ascertained in practice for a given $`H`$, one or the other status is guaranteed to apply. A question of long-standing interest has been whether an equally clear-cut classification of systems exists in quantum mechanics. Translating the criterion of classical integrability into quantum mechanics for systems with few degrees of freedom opens up loopholes of ambiguity that are not easily closed. Quantum mechanically, a second integral of the motion, i.e. an operator $`I`$ with $`[H,I]=0`$ can always be constructed, for example, via time average of an arbitrary operator $`A`$. Performing the time average in the energy representation eliminates all off-diagonal matrix elements of $`A`$. Which attributes of quantum invariants are most sensitive to the integrability status of the system? Quantum chaos research has identified a catalog of attributes that distinguish quantized nonintegrable from quantized integrable systems. The most widely studied distinctive properties pertain to level statistics. However, in the extreme quantum limit of a typical model system, where the density of energy levels is low, this distinction is blurry at best or altogether unrecognizable. Only in the energy range where the level density is high, which includes the semiclassical regime, do the contrasting level spacing distributions come into focus. Other indicators of quantum chaos are similarly ambiguous. One unequivocal discriminant between quantized integrable and nonintegrable systems was recently identified in a study of level crossing manifolds in the parameter space of a two-spin model. The system is specified by the quadratic Hamiltonian $$H=\underset{\alpha =x,y,z}{}\left\{J_\alpha S_1^\alpha S_2^\alpha +\frac{1}{2}A_\alpha \left[(S_1^\alpha )^2+(S_2^\alpha )^2\right]\right\}$$ (1) for two quantum spins $`𝐒_1,𝐒_2`$ of equal length $`\sqrt{\sigma (\sigma +1)}`$ ($`\sigma =\frac{1}{2},1,\frac{3}{2},\mathrm{}`$). In the classical limit $`\mathrm{}0`$, $`\sigma \mathrm{}`$, $`\mathrm{}\sqrt{\sigma (\sigma +1)}=s`$, the operators $`𝐒_i`$ turn into 3-component vectors, $`𝐒_i=s(\mathrm{sin}\vartheta _i\mathrm{cos}\phi _i`$, $`\mathrm{sin}\vartheta _i\mathrm{sin}\phi _i`$, $`\mathrm{cos}\vartheta _i)`$, and Eq. (1) then describes the energy function of an autonomous Hamiltonian system with two degrees of freedom and canonical coordinates $`p_i=s\mathrm{cos}\vartheta _i,q_i=\phi _i,i=1,2`$. The classical integrability condition was shown to have the form $`(A_xA_y)(A_yA_z)(A_zA_x)`$ (2) $`+{\displaystyle \underset{\alpha \beta \gamma =\mathrm{cycl}(xyz)}{}}J_\alpha ^2(A_\beta A_\gamma )=0.`$ (3) Quantum mechanically, the Hamiltonian (1) is expressible as a real symmetric block-diagonal matrix, where each of the infinitely many finite-dimensional blocks is associated with one spin-$`\sigma `$ realization of an irreducible representation of the underlying (discrete) symmetry group (see Appendix A). The main conclusions of the level crossing study for this system may be summarized as follows: (i) In the six-dimensional (6D) parameter space of (1), level degeneracies occur on smooth 4D structures. (ii) For an invariant block of $`H`$ with $`K`$ levels, this 4D structure consists of $`K1`$ sheets, each representing one twofold $`[k,k+1]`$ level degeneracy in the sequence $`E_1E_2\mathrm{}E_K`$. (iii) All 4D level crossing sheets are completely embedded in the 5D integrability hypersurface. (iv) Under mild assumptions, the integrability condition (2) can be determined analytically from the conditions of level degeneracy in low-dimensional invariant Hilbert subspaces of $`H`$. These results strongly suggest that the notion of integrability remains meaningful for quantum systems described by finite Hamiltonian matrices, notwithstanding the fact that there exist universal algorithms for the diagonalization of finite symmetric matrices. For a deeper understanding of this subtle notion of quantum integrability, we note that classical integrability guarantees the existence of a canonical transformation $`(p_1,q_1;p_2,q_2)(J_1,\theta _1;J_2,\theta _2)`$ to action-angle coordinates. It transforms the Hamiltonian $`H(p_1,q_1;p_2,q_2)`$ and the second integral of the motion $`I(p_1,q_1;p_2,q_2)`$ into analytic functions $`H_C(J_1,J_2)`$, $`I_C(J_1,J_2)`$. Each point $`(J_1,J_2)`$ on the action plane specifies a torus in phase space. In the nonintegrable case, the actions $`J_1,J_2`$ are only defined for the surviving tori. Since the tori are no longer dense anywhere in phase space, no smooth functions $`H_C,I_C`$ on $`J_1,J_2`$ exist anymore. In a companion paper we have postulated that the underlying cause for the embedment of $`(d_I1)`$-dimensional level crossing manifolds in a $`d_I`$-dimensional integrability manifold of the parameter space (with dimensionality $`dd_I`$) is linked to the existence of action operators as constituent elements of the Hamiltonian. In that study we have demonstrated for two distinct model systems the explicit functional dependence $`H_Q(J_1,J_2)`$, $`I_Q(J_1,J_2)`$ of the Hamiltonian and the second integral of the motion on two action operators, and compared it to the similar yet different functional dependence $`H_C(J_1,J_2)`$, $`I_C(J_1,J_2)`$ of the corresponding classical invariants on the classical action coordinates. A direct comparison was facilitated by the fact that for the model systems considered there we knew not only the second integral of the motion but also found a set of separable canonical coordinates for the description of the classical time evolution. ## II Method A more indirect but no less compelling method for demonstrating the existence of action operators as constituent elements of the quantum invariants $`H,I`$ in some regions of parameter space, namely on the integrability hypersurface, and their nonexistence elsewhere is pursued here for the two-spin model (1). We investigate the functional dependence of the eigenvalues of quantum invariants on the Hamiltonian parameters, in particular across lines demarcating changes in symmetry and/or integrability status. On the integrability hypersurface (2), the natural quantum numbers of the eigenstates within any invariant Hilbert subspace of $`H`$ are the integer pairs $`(m_1,m_2)`$ specifying the eigenvalues (in units of $`\mathrm{}`$) of the action operators $`J_1,J_2`$. Henceforth we call them action quantum numbers. Elsewhere in parameter space, where level crossings between eigenstates of the same parameter space are prohibited, the natural quantum number is a single integer, the energy sorting quantum number $`n`$. What consequences do these conflicting assignments of quantum numbers in the two regions of parameter space have for the functional dependence of quantum invariants on the Hamiltonian parameters? Consider the case of a $`K`$-dimensional invariant subspace of (1) spanned by the basis given in Appendix A. The $`K`$ eigenstates $`\|k,k=1,\mathrm{},K`$ then form a star of orthonormal vectors pointing in oblique directions with respect to the coordinate axes. A tiny change of the parameters $`J_\alpha ,A_\alpha `$ causes the star of eigenvectors to rotate slightly. By monitoring the inner product between eigenvectors before and after every infinitesimal parameter change, we can keep track of all eigenvectors along the entire loop in parameter space. At the same time, we monitor the effect of the gradually transforming eigenvectors on the eigenvalues of two quantum invariants. For this purpose we choose the energy expectation value $`E_k=k\|H\|k`$ and the expectation value $`I_k=k\|A\|k`$, where $`A`$ is some function of the $`S_i^\alpha `$. When the Hamiltonian parameters $`J_\alpha ,A_\alpha `$ are varied along a path in 6D parameter space, the vector $`\|k`$ traces a path on the surface of a $`K`$-dimensional unit sphere, and the point $`(E_k,I_k)`$ leaves a trace in the plane of invariants. What if two eigenvectors are accidentally degenerate ($`E_k=E_k^{}`$), which happens when their energy eigenvalues cross each other at some point on the path in parameter space? Generically, the eigenvalues of the second invariant are different at the point of level degeneracy ($`I_kI_k^{}`$). We can always choose the second invariant such that this is the case. At the crossing point the orientation of the two eigenvectors is not fixed. However that ambiguity is removed if we impose the condition that the path of every point $`(E_k,I_k)`$ in the plane of invariants must be continuous. We shall see that varying $`J_\alpha ,A_\alpha `$ along a closed path in parameter space does not guarantee that the trace of every eigenstate in the $`(E_k,I_k)`$-plane is also closed. It may happen, for example, that two eigenvectors transform into each other in the course of one parameter-space loop, thus leaving an open trace in the plane of invariants, which will be closed only after a second traversal of the loop. The two kinds of quantum numbers assigned to eigenstates in different regions of parameter space as discussed previously, suggest the following scenario. (i) If the closed path in parameter space lies entirely on the integrability hypersurface, then the traces of all eigenstates in the plane of invariants will be closed. Along the loop, level crossings occur frequently, but the labeling of all eigenstates by the action quantum numbers $`m_1`$, $`m_2`$ remains valid on every stretch of it. (ii) If the path in parameter space lies entirely off the integrability hypersurface, the traces of all eigenstates will again be closed but for a different reason. Level crossings are prohibited in this region. All states are labeled by the energy sorting quantum number $`n`$. That label is valid along the entire loop. (iii) If the closed path in parameter space consists of a leg $`A`$ on and a leg $`B`$ off the integrability hypersurface, then the conflicting assignment of quantum numbers has the consequence that some of the traces in the plane of invariants remain open. An eigenstate $`\|k`$ may undergo one or several level crossings on leg $`A`$ of the path and thus end up at a different position in the energy-level sequence at the beginning of leg $`B`$ when the energy-sorting quantum number kicks in. As the parameters are varied along leg $`B`$ back to their starting values, the point $`(E_k,I_k)`$ is prevented from finding its way back to the original position in the plane of invariants because level crossings are now prohibited. Not surprisingly, physical reality turns out to be more complicated. However, the observations made by this method of analysis prove to be highly illuminating in regard to the relations between symmetry, integrability, and the assignment of quantum numbers. ## III Results To facilitate comparison with results obtained previously, we use the same reduced 3D parameter space as in Ref. . It is spanned by $`J_y,J_z,A_xA_y2A`$ at $`J_x=1,A_x+A_y=0,A_z=0`$. The integrability condition (2), which becomes $$A(1+J_y^22J_z^22A^2)=0,$$ (4) is satisfied on a 2D surface consisting of the plane $`A=0`$ and a hyperboloid with axis at $`A=0,J_z=0`$. Embedded in this integrability surface are 1D level crossing manifolds in patterns whose complexity increases with the number of levels in the invariant (Hilbert) subspaces under consideration. Individual eigenstates $`\|k`$ will now be tracked along closed paths in this reduced parameter space. Each path selected displays distinct characteristic features in the traces on the plane of invariants $`(E_k,I_k)`$. Here we use $`I_k=k\|(S_1^z+S_2^z)^2\|k`$. We consider invariant (Hilbert) subspaces of symmetry class A1A with $`K=6,10`$ levels corresponding to spin quantum numbers $`\sigma =4,5`$, respectively (see Appendix A). Figure 1 depicts the reduced parameter space projected onto the integrability plane $`A=0`$. The dot-dashed lines represent the level crossing manifold of $`H_{A1A}^5`$ with $`K=10`$ levels in the plane $`A=0`$. None of the intersection points of two dot-dashed lines involves triple or quadruple degeneracies. Each level crossing line can thus be labeled $`[k,k+1]`$ by the positions in the level sequence $`E_1E_2\mathrm{}E_K`$ of the two levels involved in the crossing. The integrability hyperboloid intersects the integrability plane along the two dashed lines. There exist 30 $`H_{A1A}^5`$ level crossing lines on the hyperboloid. These lines intersect the plane $`A=0`$ at seven points on each dashed line, namely on the intersection points with dot-dashed lines and on the symmetry points at $`\|J_y\|=\|J_z\|=1`$. The solid circles represent projections of paths along which we track the quantum invariants $`E_k,I_k`$. A different projection of the reduced parameter space is shown in Fig. 2. The larger circle represents a path along the intersection of the integrability hyperboloid with the plane $`J_y=0.4`$. The squares on that circle mark the locations where the ten level crossing lines on the hyperboloid for $`H_{A1A}^4`$ intersect the plane $`J_y=0.4`$. The smaller (concentric) circle represents a path that is located in the nonintegrable region of parameter space except for the two points where it intersects the integrability plane $`A=0`$ (dashed line). ### A Hallmark of integrability The first path considered is the circle $`J_y^2+J_z^2=\frac{1}{4}`$ in the plane $`A=0`$ as shown in Fig. 1. This path does not come close to any of the symmetry points (pentagons). In Fig. 3 we have plotted the ten levels of $`H_{A1A}^5`$ versus angular distance $`\alpha `$ on the circular path. We observe 20 pairwise crossings between six levels at the angles where the path intersects the dot-dashed lines in Fig. 1. No instances of level repulsion can be discerned in this plot, which is not to say that the $`\alpha `$-dependence of adjacent levels is uncorrelated. Take the six levels near the center of the spectrum. They can be divided into two groups of three levels undergoing similar oscillations along the path. The synchronicity of these oscillations is, in fact, a consequence of the (postulated) smooth dependence of the functions $`H_Q(J_1,J_2)`$ and $`I_Q(J_1,J_2)`$ on $`\alpha `$ for this path embedded in the integrability plane. In Fig. 4 we show the traces in the $`(E_k,I_k)`$-plane of the two eigenstates whose levels undergo four crossings along the path (thick lines in Fig. 3). The traces are continuous, closed, and smooth. The square and the arrow indicate the starting point and the direction of the trace. Every level crossing is represented by two vertically displaced asterisks, one on each trace. It is important to note that the traces remain perfectly smooth at the points of level crossing. The level crossings have no impact on the eigenvectors, or on the expectation values $`I_k`$. Every eigenvector loops around and returns to its original orientation in Hilbert space. Its path is largely unaffected by the presence of other eigenvectors which become instantaneously degenerate with it. It is as if vectors undergoing level crossings belonged to different invariant subspaces. The behavior of energy levels as observed in Fig. 3 and the properties of traces as seen in Fig. 4 reflect what we expect for a typical situation in an integrable system with two degrees of freedom. The two invariants $`E_k,I_k`$ are functions of two quantized actions $`J_1,J_2`$ with a smooth dependence on the Hamiltonian parameters. The discrete values of the actions define the natural quantum numbers of all levels, and each eigenstate maintains its identity along any path in parameter space notwithstanding the presence of level crossings. All traces produced along closed paths are therefore closed as well. There are two sources of complication forcing on us a refinement of this description without undermining the postulated link between quantum integrability and action operators. These two complications will be discussed next before we investigate the effects of nonintegrability. ### B Level repulsion near symmetry points The second path considered is the circle $`J_y^2+J_z^2=\frac{9}{4}`$ in the integrability plane $`A=0`$ (see Fig. 1). What makes it different from the previous path is that it passes close to the four points $`\|J_y\|=\|J_z\|=1`$, where additional degeneracies occur, caused by a higher symmetry. The ten levels of $`H_{A1A}^5`$ versus $`\alpha `$ are plotted in Fig. 5. As in Fig. 3 for the previous path, we observe 20 level crossings, each one associated with a point where the circular path intersects one of the dot-dashed lines in Fig. 1. In addition to these crossings we observe instances of level collisions at $`\alpha =n\pi /2,n=1,3,5,7`$, i.e. in the vicinity of the symmetry points. It is instructive to compare the effects of level crossings and level collisions on the traces in the plane of invariants. In Fig. 6 we show again the trace of the point $`(E_k,I_k)`$ for two states that are involved in four levels crossings (thick lines in Fig. 5), now along the second path. These traces exhibit features not seen in Fig. 4. We again observe that none of the level crossings leaves any mark on the traces, implying that the wave functions of the two eigenstates are completely unperturbed by the instantaneous level degeneracies (see asterisks). On any stretch between successive mutual crossings, both levels collide with one neighboring level, and each collision does have a dramatic effect on the traces of the states involved in the collision. Level collisions produce precipitous changes in the second invariant $`I_k`$ near the closest encounter of the colliding levels. The rapid variation of expectation values signals a strong perturbation of the wave functions in a level collision. The presence of this characteristic signature of level collisions is as conspicuous in the traces shown in Fig. 6 as is their absence in the traces shown in Fig. 4. In what might be called a hard level collision, the two states exchange wave functions in a manner like two billiard balls exchange momenta in a head-on collision. This makes it hard to distinguish a hard collision from a crossing in a plot such as Fig. 5 because of graphical resolution. A plot of one invariant versus the other (Fig. 6) is much more sensitive to that distinction. Here a hard level collision produces a variation in $`I_k`$ that looks almost like a discontinuity. The phenomena observed in Figs. 5 and 6 are not in contradiction with the assertion that the invariants $`E_k,I_k`$ are functions of two quantum actions. It tells us, however, that the dependence of these functions on the Hamiltonian parameters is singular at the symmetry points of $`H`$. The phenomenon of level repulsion in the immediate vicinity of symmetry points is then caused by invariants pertaining to the higher symmetry and by the associated additional level degeneracies. The traces of all levels depicted in Fig. 5 are closed as were all traces of the levels shown in Fig. 3. The implication is that the number of crossings between any pair of levels must be an even number. The fact is that neither the level crossings nor the level collisions can cause any confusion in the labeling of the levels by action quantum numbers along a path in the integrability plane $`A=0`$ as long as it avoids the points $`\|J_y\|=\|J_z\|=1`$ of higher symmetry with symmetry induced level degeneracies. Each eigenstate maintains its identity along such paths, or so it seems. ### C Open traces caused by a change in symmetry The third path considered is the circle $`J_y^2+J_z^2=2`$ at $`A=0`$ (see Fig. 1). It is embedded in the integrability plane and passes through the points $`\|J_y\|=\|J_z\|=1`$. The impact of these symmetry points on the energy levels is depicted in Fig. 7. What were level collisions in Fig. 5 have now turned into additional level crossings. At the symmetry points, the ten levels combine into a singlet, a doublet, a triplet, and a quadruplet. No instances of level repulsion are observable anymore. The absence of level collisions along this path is confirmed by a study of the traces in the $`(E_k,I_k)`$-plane. In Fig. 8 we show the traces of the two states that again start in the seventh and eighth positions of the level sequence. Gone are the rapid near-vertical displacements which we have identified in Fig. 6 and which were caused by level collisions. The traces in Fig. 8 are as unaffected by the new symmetry-induced level crossings as they are oblivious of crossings elsewhere in the integrability plane. However, a striking new feature makes its appearance in Fig. 8. The traces do not close in themselves after one loop around the circular path in parameter space. The eighth level becomes the seventh level after one loop, and then turns into the second level after two loops. Only after the third loop does it end up in the original eighth position of the level sequence. In Fig. 7 the three levels involved in that loop are drawn as thick lines. Inspection shows that there are two further groups of three states which transform into each other as the parameter values loop around the circle. That leaves one state (near the center of the spectrum) whose trace closes in itself after one loop. Is this phenomenon of levels transforming into each other compatible with the notion that the invariants are functions of the quantized actions with a smooth dependence on the Hamiltonian parameters? Yes if we allow the dependence on the parameters to be singular at points of higher symmetry within the integrability manifold. The presence of such singularities was already suggested by the level collisions observed in Figs. 5 and 6. The results of Figs. 7 and 8 confirm the singular parameter dependence from a different vantage point. When we start with the second path in parameter space (Sec. III B) and increase the radius of the circle gradually toward that of the third path, we observe a gradual hardening of the level collisions near the symmetry points. The hardening is characterized by increasingly sharp curvatures in the graphs of $`E_k`$ versus $`\alpha `$ (Fig. 5) and by increasingly rapid vertical variations in the graphs $`I_k`$ versus $`E_k`$ (Fig. 6). In the limiting case of this path, the sharply curved but smooth bends in the graph $`E_k`$ versus $`\alpha `$ turn into cusps, and the fast but smooth vertical variations in the graphs $`I_k`$ versus $`E_k`$ turn into discontinuities. An infinitely hard level collision is indistinguishable from a level crossing. In Figs. 7 and 8 smooth segments of graphs between singularities that belong to different colliding levels are rejoined to form entirely smooth graphs of crossing levels. Hence, if we insist that all levels maintain their identity along any closed path in the integrability plane $`A=0`$, we must interpret all level crossings that take place at the points of higher symmetry, $`\|J_y\|=\|J_z\|=1`$, as infinitely hard level collisions. All the evidence accumulated thus far still supports the existence of the functions $`H_Q(J_1,J_2)`$ and $`I_Q(J_1,J_2)`$ with a smooth parameter dependence on the integrability manifold, provided we allow for singularities at points of higher symmetry. Before we discuss the strongly contrasting properties of quantum invariants along paths that are not fully embedded in the integrability manifold of (1), we should report on yet another feature that complicates the interpretation of the integrable cases. ### D Open traces caused by topology The circle $`A^2+J_z^2=0.58`$ with center at $`J_y=0.4`$ is the fourth path along which we study the behavior of quantum invariants. This path represents a circular section of the integrability hyperboloid (4) (see Fig. 2). Like the first path considered, it does not pass near any point in parameter space where symmetry induced level degeneracies occur. The angular dependence of the six $`H_{A1A}^4`$ levels, depicted in Fig. 9, does indeed not show any level collisions just as was the case in Fig. 3 for the first path. All levels undergo several crossings along this path, and none of the crossings has any noticeable effect on the quantum invariants $`E_k,I_k`$ plotted in Fig. 10. Nevertheless, there is a major difference between the evolution of eigenstates along these two paths. Each one of the six levels shown in Fig. 9 transforms into a different level in the course of one loop of the path around the integrability hyperboloid. It takes three loops for every eigenstate to return to its original position in the level sequence. On the plane of invariants this phenomenon is reflected in open traces that connect to form two rings of three segments each as shown in Fig. 10. The two sets of levels are distinguished by line thickness. Unlike in the previous situation (Sec. III C), here the open-trace phenomenon cannot be attributed to a change of symmetry along the path. What distinguishes the first path, where open traces do not occur from the fourth path, where they do occur, is that only the former can be shrunk to a point without leaving the integrability manifold. Hence the multiple connectedness of the integrability hyperboloid forces us to allow for functions $`H_Q(J_1,J_2)`$ and $`I_Q(J_1,J_2)`$ whose dependence on the Hamiltonian parameters is still smooth but multiple-valued. With these concessions, the signature properties of quantum integrability postulated above remain fully intact. The quantum invariants $`E_k,I_k`$ exhibit strongly contrasting features when observed along paths that are not embedded in the integrability manifold. Visualizing these differences does not depend on a statistical analysis. They are unmistakenly identifiable in systems for systems with very few levels. ### E Level repulsion due to nonintegrability For a direct comparison with the previous situation, we now choose a circle with the same center as the fourth path and a somewhat smaller radius, $`J_z^2+A^2=0.3712`$. This fifth path lies off the integrability manifold except for two points where it intersects the integrability plane $`A=0`$ (see Fig. 2). However, no level degeneracies occur at these intersection points. The six $`H_{A1A}^4`$ levels versus $`\alpha `$ along the fifth path are plotted in Fig. 11. Even though the resulting pattern is vaguely similar to that observed in Fig. 9, the differences are clear-cut. All level crossings have turned into level collisions. Most of the collisions are fairly soft. The two hardest collisions are barely resolved as such on the scale of Fig. 12. None of the levels transform into each other any more. The levels are now naturally labeled by the energy sorting quantum number. Each open segment of the traces shown in Fig. 10 has turned into a closed trace. All level collisions, especially the hard ones, leave the characteristic marks on the traces in the form of a rapidly varying second invariant $`I_k`$. If we were to move the fifth path closer to the integrability hyperboloid by increasing its radius (see Fig. 2), we could observe a gradual hardening of all level collisions. The level configurations as shown in Fig. 11 would increasingly resemble those in Fig. 9. The traces as shown in Fig. 12, however, would remain very different from those pertaining to the integrable case (Fig. 10). Only in the limiting case where the fifth path merges with the fourth path would the closed traces of the nonintegrable model break into segments connected by vertical lines. The ends of each segment would then rejoin ends of other segments to form the smooth rings of open traces shown in Fig. 10. Similar observations are made upon lifting the first path off the integrability plane $`A=0`$ to a plane at $`A0`$. All the level crossings that exist in Fig. 3, for example, turn into level collisions. The closed traces such as shown in Fig. 4 break into pieces whose ends rejoin via near vertical lines into a new set of closed traces. Along the second path we had observed (in Fig. 5) level crossings (due to integrability) and level collisions (due to nearby points of higher symmetry). Lifting this path off the integrability plane again removes all level crossings and results in a set of closed traces. The characteristic marks of level collisions on the traces in the $`(E_k,I_k)`$-plane are the same no matter whether they are caused by a reduced symmetry or by nonintegrability. Lifting the third path off the integrability plane has the same effects on the level crossings attributed to integrability and the level crossings attributed to the higher symmetry at selected points in parameter space (Fig. 7). All are removed indiscriminately. ### F Open traces caused by nonintegrability The conflicting assignments of quantum numbers to eigenstates for parameter values on and off the integrability manifold is most compellingly documented when we pick a path in parameter space that is only partially embedded in the integrability manifold. The sixth path considered in this study of quantum invariants is a modification of the first path (Sec. III A) with the same projection in Fig. 1. Whereas the first path was embedded in the integrability plane $`A=0`$, the sixth path has a variable height relative to that plane: $`A(\alpha )=0.3\mathrm{cos}^2(\alpha /2)`$. It touches down to the integrability plane at a single point $`(\alpha =180^{})`$, where a level crossing takes place. Along this path there exist no other level crossings. All the other crossings that existed in Fig. 3 for the first path are now replaced by level collisions (see Fig. 13). The inevitable consequence of having a single level crossing along a closed path in parameter space is the existence of a pair of open traces in the plane of invariants, namely the traces of the states that undergo the crossing at $`\alpha =180^{}`$. These traces are shown in Fig. 14. The ends of the solid and dashed lines form a single loop, which is traced in the direction indicated. What causes here an open trace in the plane of invariants is obviously akin to what had caused an open trace in the situation described in Sec. III C. In both cases two levels cross once due to particular circumstances at one point of the path, and are thus prevented from crossing back to their original position in the level sequence on the remaining stretch of the path. In Sec. III C the particular circumstance was a higher symmetry, here it is integrability. ## IV Interpretation The study of quantum invariants along closed paths through parameter space indicates that a change in symmetry and a change in integrability status produce related phenomena. In some dynamical systems, the conservation laws that guarantee integrability are direct consequences (via Noether’s theorem) of continuous symmetries. Switching from integrability to nonintegrability is then accompanied by a reduction in symmetry. In the two-spin model (1), the presence of a (continuous rotational) O(2) or higher symmetry in spin space does indeed imply the existence of a second integral of the motion, namely the component of the total spin along the symmetry axis, and integrability is guaranteed. However, a second integral of the motion was shown to exist for certain parameter values even in the absence of a continuous rotational symmetry. Does integrability in that case indicate the presence of a hidden symmetry? Classical integrability guarantees that the Hamiltonian (1) can be expressed as a function of the two action variables: $`H=H_C(J_1,J_2)`$. The cyclical nature of the angle coordinates thus implies that $`H_C`$ is invariant with respect to continuous rotation-like transformations in phase-space. Since this is not related to a continuous symmetry in configuration space, it is appropriate to call it a hidden symmetry. For a description of the impact of symmetries on the level spectrum of the quantum two-spin model, it is useful to distinguish three kinds of symmetry: discrete symmetries, continuous symmetries, and hidden symmetries. Discrete symmetries have no bearing on the classical integrability property, but they do affect the shapes of phase-space trajectories. Quantum mechanically, they divide the Hilbert space of $`H`$ into invariant subspaces. In general, this does not result in symmetry-induced level-degeneracies, but it does lead to accidental degeneracies between levels belonging to different invariant subspaces. Such level crossings exist independently of whether or not $`H`$ is integrable. Hidden symmetries, which guarantee classical integrability, cause additional accidental level degeneracies, namely between states within one of the invariant subspaces pertaining to any existing discrete symmetry. Continuous symmetries, in essence, combine the effects of the discrete and hidden symmetries, and allow accidental inter-subspace degeneracies. In addition to these effects, continuous symmetries (sometimes in tandem with discrete symmetries) produce level degeneracies of a permanent nature, the so-called symmetry-induced level degeneracies. There exists a hierarchy of symmetries in the two-spin model (1): (S0) In the absence of any symmetry, there are no level degeneracies. All levels will collide when Hamiltonian parameters are varied. This situation can be realized by an external magnetic field. (S1) The existence of discrete symmetries alone produces finite-D invariant Hilbert subspaces. Level crossings exist between states belonging to different subspaces. Levels within any subspace collide. (S2) The existence of hidden symmetries in addition to discrete symmetries produces level crossings between states in the same invariant subspace. (S3) The continuous symmetries produce permanent degeneracies in certain regions of parameter space. There exists a hierarchy of level collisions which corresponds to the hierarchy of symmetries. (S1$``$S0) Inter-subspace level crossings in the presence of discrete symmetries turn into level collisions when discrete symmetries are removed. (S2$``$1) Intra-subspace level crossings turn into level collisions when the hidden symmetries are removed, i.e. when the integrability is destroyed. (S3$``$S2) Symmetry-induced level degeneracies associated with a continuous symmetry are removed outside the range of that symmetry irrespective of the presence or absence of the hidden symmetry. Some level crossings along paths through symmetry points in parameter space turn into level collisions along nearby paths that miss the symmetry point. Other level crossings are insensitive to whether the path hits or misses the symmetry point. They are the product of the hidden symmetry. All phenomena observed in the quantum invariants $`E_k,I_k`$ along closed paths on, off, and across the integrability manifold, indicate that the effects of a change in integrability status are akin to the effects of a change in symmetry. All observations point to the existence of a hidden symmetry that accompanies quantum integrability. In the classical limit, this hidden symmetry manifests itself in phase space when viewed from a particular coordinate system – the action-angle coordinates. The same hidden symmetry must also exist in the quantum system, but only on the integrability manifold. Even though nonintegrability is not to be taken literally in the quantum case, the presence or absence of that hidden symmetry has consequences that are equally clear-cut as in the classical limit. ###### Acknowledgements. This work was supported by the Research Office of the University of Rhode Island. We are very grateful to Joachim Stolze for his comments and suggestions relating to this work. ## A Discrete symmetries The (discrete) symmetry group relevant for the general 2-spin Hamiltonian (1) is $`D_2S_2`$, where $`D_2`$ contains the three twofold rotations $`C_2^\alpha `$, $`\alpha =x,y,z`$ about the coordinate axes, and $`S_2`$ the permutations of the two spins. The eight irreducible representations of $`D_2S_2`$ are named A1S, A1A, B1S, B1A, B2S, B2A, B3S, B3A, where S (A) stand for (anti-)symmetric under permutation and A1, B1, B2, B3 for $`(C_2^x,C_2^y,C_2^z)=(1,1,1)`$, $`(1,1,1)`$, $`(1,1,1)`$, $`(1,1,1)`$, respectively. The basis vectors with transformation properties corresponding to the eight different irreducible representations $`R`$ are listed in Table I for integer $`\sigma `$ and in Table II for half-integer $`\sigma `$. The Hamiltonian matrix can then be expressed in the form $$H=\underset{R,\sigma }{}H_R^\sigma $$ (A1) with blocks of dimensionalities $`K=1,3,6,10,\mathrm{}`$ in 16 different realizations, two for each symmetry class (one with integer $`\sigma `$ and one with half-integer $`\sigma `$).
warning/0003/math0003226.html	ar5iv	text	# Untitled Document On Holomorphic Jet Bundles by Pit-Mann Wong and Wilhelm Stoll Department of Mathematics University of Notre Dame Notre Dame, Indianna 46556 USA Introduction In this article we provide a more detailed discussion (see \[W4\]) of the jet bundles introduced by Green and Griffiths \[G-G\]. In section 1 some basic facts about these jet bundles (which are different from the usual jet bundles used in analysis) are established with the most important one being a Theorem of Green and Griffiths concerning the natural filtration of the sheaf, denoted $`𝒥_k^mX`$, of $`k`$-jet differentials of weight $`m`$. With this reuslt many computations (of Chern classes) and properties of $`𝒥_k^mX`$ can be obtained or inferred from the more familiar objects $`^{i_1}T^{}X\mathrm{}^{i_k}T^{}X`$ (satisfying the condition $`i_1+2i_2=\mathrm{}+ki_k=m`$). The calculation of Chern classes of $`𝒥_k^mX`$ are carried out in section 1 for curves and in section 2 for surfaces. These are needed later in applying the Riemann-Roch Theorem for $`𝒥_k^mX`$ and its corresponding line sheaf, $`𝒪_{𝐏(J^kX)}(m)`$, over the projectivized bundle $`𝐏(J^kX)`$. There are some complications in working with these sheaves due to the fact that the fibers of $`𝐏(J^kX)`$ are weighted projective spaces and hence not smooth and moreover, the natural sheaves $`𝒪_{𝐏(J^kX)}(m)`$ are not necessarily locally free if $`m`$ is not divisible by $`k!`$. These minor difficulties are clarified in section 3 and is readily seen to be rather harmless. In section 4 we consider the case of surfaces of general type and here there is another complication due to the fact that, as oppose to the bundles $`^{i_1}T^{}X\mathrm{}^{i_k}T^{}X`$, the sheaves $`𝒥_k^mX`$ are not semi-stable (with respect to the canonical bundle of $`X`$). This difficulty, however, can be overcome rather easily and as a result we obtain applications in the theory of holomorphic curves in surfaces of general type (hypersurfaces in $`𝐏^3`$ in particular). These are presented in section 4. We also include two appendices. In appendix A the lemma of logarithmics derivatives and a version of Schwarz Lemma are presented (see \[L\], \[L-Y\], \[DSW1\], \[DSW2\], \[S-Y\], \[W3\] and \[J\]). Some combinatorics related to the symmetric groups which we used in the computation of Chern classes (this comes in, for example, in counting the number of positive integer solutions of the equation $`i_1+2i_2+\mathrm{}+ki_k=m`$) are presented in appendix B. For higher dimensional manifolds the approach of Nevanlinna Theory appears to work better (see \[W5\]). Nevanlinna Theory for symmetric and exterior products of the cotangent bundle can be found in \[St\]. § 1 Holomorphic Jet Bundles We examine two concepts of ”jet bundles” of a complex manifold. The first is the jet bundles used by analysts (PDE) and also by Faltings in his work on rational points of an ample subvariety of an abelian variety and integral points of complement of an ample divsior of an abelian variety \[F\]. The second is the jet bundles introduced by Green and Griffiths \[G-G\]. The first notion of jet bundle shall henceforth be referred to as the full jet bundle while the second notion of jet bundle shall be referred to as the restricted jet bundle. The reason for these terminologies is that the fiber dimension of the full jet bundle is much larger than that of the restricted jet bundle. For a complex manifold $`X`$ the (locally free) sheaf of germs of holomorphic tangent vector fields (differential operators of order 1) of $`X`$ shall be denoted by $`T^1X`$ or simply $`TX`$. An element of $`T^1X`$ acts on the sheaf of germs of holomorphc functions by differentiation: $$(D,f)T^1X\times 𝒪_XDf𝒪_X$$ and the action is linear over the complex number field $`𝐂`$, i.e., $$Dom_𝐂(𝒪_X,𝒪_X).$$ This concept can be extended as follows: Definition 1.1 Let $`X`$ be a complex manifold of dimension $`n`$ the sheaf of germs of holomorphic $`k`$-jets (differential operators of order $`k`$), denoted $`𝒯^kX`$, is the subsheaf of the sheaf of homomorphisms $`om_𝐂(𝒪_X,𝒪_X)`$ consisting of elements (differential operators) of the form $$\underset{j=1}{\overset{k}{}}\underset{i_j𝐍}{}D_{i_1}\mathrm{}D_{i_j}$$ where $`D_{i_j}T^1X`$. In terms of holomorphic coordinates $`z_1,\mathrm{},z_n`$ an element of $`𝒯^kX`$ is expressed as: $$\underset{j=1}{\overset{k}{}}\underset{1i_1\mathrm{}i_jn}{}a_{i_1,\mathrm{},i_j}\frac{^j}{z_{i_1}\mathrm{}z_{i_j}}$$ where the coefficients $`a_{i_1,\mathrm{},i_j}`$ are holomorphic functions. We can also drop the reqirement that the indices be non-decreasing by requiring symmetry in the coefficients, in other words, the elements of $`𝒯^kX`$ can also be expressed as: $$\underset{j=1}{\overset{k}{}}\underset{1i_1,\mathrm{},i_jn}{}a_{i_1,\mathrm{},i_j}\frac{^j}{z_{i_1}\mathrm{}z_{i_j}}$$ where the coefficients $`a_{i_1,\mathrm{},i_j}`$ are symmetric in the indices $`i_1,\mathrm{},i_j`$, i.e., if $`\sigma `$ is an element of the symmetric group of $`j`$ elements then $$a_{i_{\sigma (1)},\mathrm{},i_{\sigma (j)}}=a_{i_1,\mathrm{},i_j}.$$ The effect of holomorphic change of coordinates from $`z=(z_1,\mathrm{},z_n)`$ to $`w=(w_1,\mathrm{},w_n)`$ is given by the transistion function (for $`k=2`$): $`\left(\begin{array}{c}(\frac{}{z_i})_{1in}\\ (\frac{^2}{z_iz_k})_{1ikn}\end{array}\right)=\left(\begin{array}{cc}A0& \\ BC& \end{array}\right)\left(\begin{array}{c}(\frac{}{w_j})_{1jn}\\ (\frac{^2}{w_jw_l})_{1jln}\end{array}\right)`$ (7) where $`A`$ is the $`n`$ by $`n`$ matrix: $$A=(\frac{w_j}{z_i})_{1i,jn},$$ while $`B`$ is the $`C_2^{n+1}`$ by $`n`$ matrix: $$B=(\frac{^2w_j}{z_iz_k})_{1ikn}$$ $`C`$ is the $`C_2^{n+1}`$ by $`C_2^{n+1}`$ matrix: $$C=(\frac{w_j}{z_i}\frac{w_l}{z_k})_{1ikn,1jln}$$ and $`0=0_{n\times C_2^{n+1}}`$ is the $`n`$ by $`C_2^{n+1}`$ zero-matrix (here $`C_2^{n+1}=(n+1)!/(n1)!2!`$ is the usual binomial coefficient). Note that the matrix $`A`$ is the transistion function for the tangent bundle $`TX`$ while the matrix $`C`$ is the transistion function of $`^2TX`$, the 2-fold symmetric product of the tangent bundle. For general $`k`$ the transistion function of $`T^kX`$ is of the form: $`\left(\begin{array}{cccccc}A_100000& & & & & \\ A_20000& & & & & \\ .000& & & & & \\ .00& & & & & \\ .0& & & & & \\ A_k& & & & & \end{array}\right)`$ where the $`C_j^{n+j1}`$ by $`C_j^{n+j1}`$ matrix $`A_j`$ is the transistion function of the bundle $`^jTX`$, the $`j`$-fold symmetric product of the tangent bundle. Here $`C_j^{n+j1}=(n+j1)!/j!(n1)!`$ is the usual binomial coefficient. The linear (and invertible) nature of the transistion functions implies that $`T^kX`$ is locally free. This can also be seen by observing that $`T^{k1}X`$ injects into $`T^kX`$ and there is an exact sequence of sheaves: $`0T^{k1}XT^kXT^kX/T^{k1}X0`$ (9) where $`T^kX/T^{k1}X^kT^1X`$ (10) is the sheaf of germs of $`k`$-fold symmetric product of $`T^1X`$, i.e., sheaf of germs of operators of the form: $$\underset{1i_1\mathrm{}i_jn}{}a_{i_1,\mathrm{},i_j}\frac{^k}{x_{i_1}\mathrm{}x_{i_k}}.$$ These exact sequences imply, by induction, that $`𝒯^kX`$ is locally free as the sheaves $`^kT^1X`$, being the symmetric product of the tangent sheaf, is locally free. We include here the proof of the isomorphism (3). Proposition 1.2 With the notations above we have: $$T^kX/T^{k1}X^kTX$$ where $`TX`$ is the $`k`$-fold symmetric product of the tangent bundle. Proof. We shall define a surjection from the $`k`$-fold tensor product of $`TX`$ onto the quotient $`T^kX/T^{k1}X`$: $$\mu :^kTXT^kX/T^{k1}X$$ and then show that the surjection factors through the symmetric product resulting in a bijection. The map $`\mu `$ is defined by: $$\mu (D_1\mathrm{}D_k)=[D_1\mathrm{}D_k]$$ where $`D_i`$ is (the germ of) a vector field and $`[]:T^kXT^kX/T^{k1}X`$ is the quotient map. By definition the map $`\mu `$ is surjective. To see that the map factors through to the symmetric product it is sufficient to show that the map is invariant by any transposition, i.e., $$\mu (D_1\mathrm{}D_iD_{i+1}\mathrm{}D_k)=\mu (D_1\mathrm{}D_{i+1}D_i\mathrm{}D_k).$$ This follows from the fact that the Lie bracket $`D_iD_{i+1}D_{i+1}D_i`$ of the vector fields $`D_i`$ and $`D_{i+1}`$ is again a vector field and not a 2-jet. Thus we have: $$D_1\mathrm{}(D_iD_{i+1}D_{i+1}D_i)\mathrm{}D_kT^{k1}X$$ which implies that the map $`\mu `$ descends to the symmetric product $`^kTX`$. More precisely, if we denote the symmetrization operator by $`\sigma _k`$ then $$\overline{\mu }(D_1\mathrm{}D_k)=\overline{\mu }(\sigma _k(D_1\mathrm{}D_k))\stackrel{\mathrm{def}}{=}\mu (D_1\mathrm{}D_k)$$ is well-defined. It is clear that $`\overline{\mu }`$ is surjective and it remains to show that $`\overline{\mu }:^kTXT^kX/T^{k1}X`$ is injective. Let $`(z_1,\mathrm{},z_n)`$ be a local coordinate near a point $`xX`$ then $$\frac{^k}{x_{i_1}\mathrm{}x_{i_k}}=\frac{}{x_{i_1}}\mathrm{}\frac{}{x_{i_k}},1i_1\mathrm{}i_kn$$ is a basis of $`^kTX`$ at the point $`xX.`$ If $`\overline{\mu }`$ is not injective then there exists a differential operator $`\mathrm{\Psi }`$ of the form $$\mathrm{\Psi }=\underset{1i_1\mathrm{}i_kn}{}a_{i_1\mathrm{}i_k}\frac{^k}{x_{i_1}\mathrm{}x_{i_k}}\mathrm{\Phi }$$ where $`\mathrm{\Phi }`$ is a differential operator of order at most $`k1`$ and such that $`\overline{\mu }(\mathrm{\Psi })=0`$. Apply the operator $`\mathrm{\Psi }`$ to the function $`f=x_{i_1}\mathrm{}x_{i_k}`$ shows that this is possible only if all the coefficients $`a_{i_1\mathrm{}i_k}`$ are zero. This shows that $`\overline{\mu }`$ is injective and completes the proof of the Proposition. QED The restricted $`k`$-jet bundles are introduced by Green-Griffiths in \[G-G\]. It is defined as follows. Denote by $`_x,xX`$ the sheaf of germs of holomorphic curves $`\{f:\mathrm{\Delta }_rX,f(0)=x\}`$. Define, for $`k𝐍`$, an equivalence relation as follows. Let $`z_1,\mathrm{},z_n`$ be holomorphic coordinates near $`x`$ and for $`f_x`$ let $`f_i=z_if`$. Two elements $`f,g_x`$ are said to be $`k`$-equivalent, denoted $`f_kg`$, if $`f_j^{(p)}(0)=g_j^{(p)}(0)`$ for all $`1pk`$. The sheaf of restricted $`k`$-jets is defined to be $`J^kX=_{xX}_x/_k`$. Elements of $`J^kX`$ will be denoted by $`j^kf(0)=(f(0),f^{^{}}(0),\mathrm{},f^{(k)}(0))`$. It is also clear that $`J^1X=T^1X=TX`$ is the tangent bundle. Note that the definition above does not depend on the choices of the coordinates near $`x`$; for if $`(z_jf)^{(p)}(0)=(z_jg)^{(p)}(0)`$ for $`0pk`$ then $`(w_jf)^{(p)}(0)=(w_jg)^{(p)}(0)`$ for any other coordinates $`w_1,\mathrm{},w_n`$. The effect of change of coordinates is as follows: $`(w_jf)^{^{}}={\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{w_j}{z_i}}(f)(z_if)^{^{}},`$ $`(w_jf)^{^{\prime \prime }}={\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{w_j}{z_i}}(f)(z_if)^{^{\prime \prime }}+{\displaystyle \underset{i,k=1}{\overset{n}{}}}{\displaystyle \frac{^2w_j}{z_iz_k}}(f)(z_if)^{^{}}(z_kf)^{^{}}`$ and for general $`k`$ we have $`(w_jf)^{(k)}={\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{w_j}{z_i}}(f)(z_if)^{(k)}+P({\displaystyle \frac{^lw_j}{z_{i_1}\mathrm{}z_{i_l}}}(f),(w_jf)^{(l)})`$ (11) where $`P`$ is a polynomial with integer coefficients in $`^lw_j/z_{i_1}\mathrm{}z_{i_l},(w_jf)^{(l)}`$ for $`j=1,\mathrm{},n`$ and $`l=1,\mathrm{},k`$. Note that the quadratic nature in $`(z_if)^{^{}}`$ in the formula for $`(w_if)^{^{\prime \prime }}`$ means that the sheaf $`J^kX`$ is not locally free. It is instructive to compare this with the matrix $`C`$ in the transition formula (1) for $`T^kX`$ under the change of coordinates. In that case the formula is quadratic in the partial derivatives $`w_j/z_i`$ but linear in $`(z_if)^{^{}}`$ hence the transformation can still be represented as a linear transformation while this is not the case for $`J^kX`$. There is however, a natural $`𝐂^{}`$-action on $`J^kX`$ defined via parametrization. Namely, for $`\lambda 𝐂^{}`$ and $`f_x`$ a map $`f_\lambda _x`$ is defined by $`f_\lambda (t)=f(\lambda t)`$; then $`j^kf_\lambda (0)=(f_\lambda (0),f_\lambda ^{^{}}(0),\mathrm{},f_\lambda ^{(k)}(0))=(f(0),\lambda f^{^{}}(0),\mathrm{},\lambda ^kf^{(k)}(0)).`$ In other words the $`C^{}`$-action is given by $`\lambda .j^kf(0)=(f(0),\lambda f^{^{}}(0),\mathrm{},\lambda ^kf^{(k)}(0)).`$ (12) Definition 1.3 The restricted $`k`$-jet bundle is defined to be $`J^kX`$ together with the $`𝐂^{}`$-action defined above and shall simply be denoted by $`J^kX`$. Another difference between the full and restricted $`k`$-jet bundles is that there is, in general, no natural way of injecting $`J^{k1}X`$ into $`J^kX`$. For instance, the coordinates transformations shows that $$(f(0),(f_1^{^{}}(0),\mathrm{},f_n^{^{}}(0)))(f(0),(f_1^{^{}}(0),\mathrm{},f_n^{^{}}(0)),(0,\mathrm{},0))$$ is not a well-defined map of $`J^1X`$ into $`J^2X`$ as the condition $`f^{^{\prime \prime }}(0)=0`$ is not preserved by a general change of coordinates. On the other hand, the transformation formulas show that $`\{j^2f=(f(0),f^{^{}}(0),f^{^{\prime \prime }}(0))J^2X\|f^{^{}}(0)=0\}`$, more generally, $`Z_0=\{j^kf=(f(0),f^{^{}}(0),\mathrm{},f^{(k)}(0))J^kX\|f^{^{}}(0)=0\}`$ (13) is a well-defined subvariety of $`J^kX`$ as the conditon $`f^{^{}}(0)`$ is invariant under change of coordinates. Moreover, the transformation law actually says that even though the condition $`f^{^{\prime \prime }}(0)=0`$ is coordinate dependent the conditions that $`f^{^{}}(0)=f^{^{\prime \prime }}(0)=0`$ are independent of choices of coordinates, in other words, the zero-section of $`J^2X`$, more generally, the zero-section of $`J^kX`$: $`\{j^kf(0)J^kX\|f^{^{}}(0)=f^{^{\prime \prime }}(0)=\mathrm{}=f^{(k)}(0)=0\}`$ (14) is well-defined. Theorem 1.4 Let $`X`$ be a complex manifold of dimension $`n`$ then $`T^kX`$ is a holomorphic vector bundle of rank $`r=n+C_2^{n+1}+C_3^{n+2}+\mathrm{}+C_k^{n+k1}=_{i=1}^kC_i^{n+i1}`$ while $`J^kX`$ is a holomorphic $`𝐂^{}`$-bundle of rank $`r=kn`$ and the zero-section of $`J^kX`$ is well-defined. As noted above there is no natural inclusion map from $`J^{k1}X`$ into $`J^kX`$ there is however a natural projection map $$p_{kj}:J^kXJ^jX$$ for any $`jk`$ defined simply by $`p_{kj}(j^kf(0))=j^jf(0).`$ (15) The projection map clearly respect the $`𝐂^{}`$-action defined by (5) and so is a $`𝐂^{}`$-bundle morphism. If $`\mathrm{\Phi }:XY`$ is a holomorphic map between the complex manifolds $`X`$ and $`Y`$ then the usual differentail $`\mathrm{\Phi }_{}:T^1XT^1Y`$ is defined. The same is true for the $`k`$-jets as the $`k`$-th order differential $`\mathrm{\Phi }_k:T^kXT^kY`$ can be defined by $$\mathrm{\Phi }_k=(D_1\mathrm{}D_k)(g)\stackrel{\mathrm{def}}{=}D_1\mathrm{}D_k(g\mathrm{\Phi })$$ where $`g𝒪_Y`$. The $`k`$-th order differential, denoted $`J^k\mathrm{\Phi }:J^kXJ^kY`$ can also be defined: $$J^k\mathrm{\Phi }(j^kf(0))\stackrel{\mathrm{def}}{=}(\mathrm{\Phi }f)^{(k)}(0)$$ for any $`j^kf(0)J^kX`$. For the restricted jet bundle $`J^kX`$ there is another notion closely related to (but not the same) the differential: the natural lifting of a holomorphic curve. Namely, given any holomorphic map $`f:\mathrm{\Delta }_rX(0<r\mathrm{}`$), the lifting $`j^kf:\mathrm{\Delta }_{r/2}J^kX`$ is defined by: $$j^kf(\zeta )=j^kg(0),\zeta \mathrm{\Delta }_{r/2}$$ where $`g(\xi )=f(\zeta +\xi )`$ is holomorphic for $`\xi \mathrm{\Delta }_{r/2}`$. Consider the special case $`dimX=1`$ then $`T^kX`$ and $`J^kX`$ have the same rank and the underlying space of $`T^kX`$ and $`J^kX`$ are the same but the structures are different. Consider the map (for simplicity we write this out only for $`k=2`$): $`(f(0),f^{^{}}(0),f^{(^{\prime \prime })}(0))f^{^{\prime \prime }}(0){\displaystyle \frac{}{z}}+(f^{^{}}(0))^2{\displaystyle \frac{^2}{z^2}}.`$ (16) which is clearly holomorphic but is not biholomorphic. For if $$g(t)=f(0)f^{^{}}(0)t+f^{^{\prime \prime }}(0)t^2/2$$ then $`j^2g(0)=(f(0),f^{^{}}(0),f^{^{\prime \prime }}(0))`$ and under the identification above $`j^2f(0)`$ and $`j^2g(0)`$ are mapped onto the same element. Moreover the map is a $`𝐂^{}`$-bundle map because $`\lambda .(f(0),f^{^{}}(0),f^{(^{\prime \prime })}(0))=(f(0),\lambda f^{^{}}(0),\lambda ^2f^{(^{\prime \prime })}(0))`$ is mapped onto the element $$\lambda ^2\{f^{^{\prime \prime }}(0)\frac{}{z}+(f^{^{}}(0))^2\frac{^2}{z^2}\}.$$ More generally, for $`X`$ of arbitrary dimension, the map $`(f(0),f^{^{}}(0),f^{(^{\prime \prime })}(0)){\displaystyle \underset{i=1}{\overset{n}{}}}f_i^{^{\prime \prime }}(0){\displaystyle \frac{}{z_i}}+{\displaystyle \underset{1i,jn}{}}f_i^{^{}}(0)f_j^{^{}}(0){\displaystyle \frac{^2}{z_iz_j}}`$ (17) is a holomorphic $`𝐂^{}`$-bundle map from $`J^2X`$ onto a $`𝐂^{}`$ sub-bundle of $`T^2X`$. We have already seen the case of $`n=1`$; for $`n=2`$ the second sum above has 3 terms: $$(f_1^{^{}})^2,(f_2^{^{}})^2,f_1^{^{}}f_2^{^{}}.$$ Thus if $`j^2f(0)`$ and $`j^2g(0)`$ have the same image then $$(f_1^{^{}})^2=(g_1^{^{}})^2,(f_2^{^{}})^2=(g_2^{^{}})^2,f_1^{^{}}f_2^{^{}}=g_1^{^{}}g_2^{^{}}$$ so that $`f_1^{^{}}=\pm g_1^{^{}},f_2^{^{}}=\pm g_2^{^{}}`$. This means that the map is generically 2 to 1 onto its image and ramified along the subvariety $`Z_0`$ defined by (6). Returning to the case of a Riemann surface $`X`$ we define a map $`p_3:J^3XT^3X`$ by the formula: $$p_3(j^3f(0))=f^{(3)}(0)\frac{}{z}+f^{^{\prime \prime }}(0)f^{^{}}(0)\frac{^2}{z^2}+(f^{^{}}(0))^3\frac{^3}{z^3}$$ and in general $`p_k:J^kXT^kX`$ by the formula: $$p_k(j^kf(0))=\underset{j=1}{\overset{k}{}}f^{(j)}(0)(f^{^{}}(0))^{kj}\frac{^{kj+1}}{z^{kj+1}}.$$ For the higher dimensional manifold $`X`$ the maps are defined by $`p_3(j^3f(0))`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}f_i^{(3)}(0){\displaystyle \frac{}{z_i}}+{\displaystyle \underset{i,j=1}{\overset{n}{}}}f_i^{^{\prime \prime }}(0)f_j^{^{}}(0){\displaystyle \frac{^2}{z_iz_j}}+`$ $`+{\displaystyle \underset{i,j,k=1}{\overset{n}{}}}f_i^{^{}}(0)f_j^{^{}}(0)f_k^{^{}}(0){\displaystyle \frac{^3}{z_iz_jz_k}}`$ and in general $`p_k:J^kXT^kX`$ by the formula: $`p_k(j^kf(0))={\displaystyle \underset{j=1}{\overset{k}{}}}{\displaystyle \underset{i=1}{\overset{n}{}}}\{f_i^{(j)}(0){\displaystyle \underset{i_1,\mathrm{},i_{kj}i}{}}f_{i_l}^{^{}}(0)\}{\displaystyle \frac{^{kj+1}}{z_iz_{i_1}\mathrm{}z_{i_{kj}}}}.`$ (18) It is clear from the definition of the map $`p_k`$ that: Theorem 1.5 Let $`J^kX`$ and $`T^kX`$ be, respectively, the restriced and the full $`k`$-jet bundles over a complex manifold $`X`$. Then the map defined by $`(10)`$ is a holomorphic $`𝐂^{}`$-bundle map which is generically finite to $`1`$ onto its image. Moreover, the map is ramified precisly along the subvariety $`Z_0=\{j^kf(0)J^kX\|f^{^{}}(0)=0\}.`$ We consider now the ”dual” of the jet bundles. Definition 1.6 The dual of the full jet bundles $`T^kX`$ shall be referred to as the sheaf of germs of $`k`$-jet forms and shall be denoted by $`T_k^{}X`$. The global sections shall be referred to as $`k`$-jet forms. For $`m𝐍`$ the $`m`$-fold symmetric product shall be denoted by either $`^mT_k^{}X`$ and its global sections shall be referred to as $`k`$-jet forms of weight $`m`$. By definition, a $`k`$-jet form of weight $`m`$ assigns to each point $`xX`$ a homogeneous (with respect to the standard $`𝐂^{}`$-action of $`T^kX`$ as a vector bundle) polynomial of degree $`m`$ on the fiber $`T_x^kX`$ (where $`T^kX`$ is the $`k`$-jet bundle). Let $`(U,z_1,\mathrm{},z_n)`$ be a local holomorphic coordinates over $`U`$ then $`(e_i={\displaystyle \frac{}{z_i}})_{1in},`$ $`(e_{i_1i_2}={\displaystyle \frac{^2}{z_{i_1}z_{i_2}}})_{1i_1i_2n},`$ $`.`$ $`.`$ $`.`$ $`(e_{i_1\mathrm{}i_k}={\displaystyle \frac{^k}{z_{i_1}\mathrm{}z_{i_k}}})_{1i_1i_2\mathrm{}i_kn}`$ is a basis of $`T^kX\|_U`$. The dual basis shall be denoted, formally, by $`(e_i^{}=dz_i)_{1in},`$ $`(e_{i_1i_2}^{}=d^2z_{i_1}z_{i_2})_{1i_1i_2n},`$ $`.`$ $`.`$ $`.`$ $`(e_{i_1\mathrm{}i_k}^{}=d^kz_{i_1}\mathrm{}z_{i_k})_{1i_1i_2\mathrm{}i_kn}.`$ An element of $`T_k^{}X`$ is then of the form: $`\omega ={\displaystyle \underset{j=1}{\overset{k}{}}}{\displaystyle \underset{1i_1\mathrm{}i_jn}{}}a_{i_1,\mathrm{},i_j}e_{i_1,\mathrm{},i_j}^{}`$ where the coefficients $`a_{i_1,\mathrm{},i_j}`$ are holomorphic functions. Sometimes it is convenient to express the sum without the restriction as in the second sum above but insisting on the symmetry of the coefficients (see also definition 1.1): $`\omega ={\displaystyle \underset{j=1}{\overset{k}{}}}{\displaystyle \underset{1i_1,\mathrm{},i_jn}{}}a_{i_1,\mathrm{},i_j}e_{i_1,\mathrm{},i_j}^{}`$ where the coefficients $`a_{i_1,\mathrm{},i_j}`$ are symmetric in the indices. We can also write down a basis for the symmetric product $`^mT^kX`$ and its dual basis for $`^mT_k^{}X`$. It is convenient to use the following notations and conventions for the index set. Let $$_k=\{I=(i_1,\mathrm{},i_k)\|i_k𝐍\{0\},0i_1\mathrm{}i_kn\mathrm{and}\mathrm{not}\mathrm{all}i_j=0\}$$ be endowed with lexicographical order and let $$𝒥_m(_k)=\{J=(I_1,\mathrm{},I_m)\|I_j_k,I_1\mathrm{}I_m\}.$$ With these notations, for example, the basis for $`T_k^{}X`$ over $`U`$ is simply expressed as $`B_k^{}=\{e_I^{}=e_{i_1}^{}\mathrm{}e_{i_k}^{}\|I_k\}`$ with the conventions that $`e_0=1`$. Analogously, a basis for $`^mT_k^{}X`$ over $`U`$ is expressed as $`B_k^m=\{e_J^{}=e_{I_1}^{}\mathrm{}e_{I_k}^{}\|J𝒥_m(_k)\}`$. Moreover, a section $`\omega H^0(U,^mT_k^{}X)`$ is expressed as $`\omega ={\displaystyle \underset{J𝒥_m(_k)}{}}a_Je_J^{}`$ where the coefficients are holomorphic functins on $`U`$. Taking the dual of the sequence (2) we get an exact sequence: $`0^kT_1^{}XT_k^{}XT_{k1}^{}X0.`$ (19) For example, for $`k=3`$ we have two exact sequences: $$0^3T_1^{}XT_3^{}XT_2^{}X0,$$ $$0^2T_1^{}XT_2^{}XT_1^{}X0.$$ In particular, by Whitney’s Formula: $$c_1(T_3^{}X)=c_1(T_2^{}X)+c_1(^3T_1^{}X)=c_1(T_1^{}X)+c_1(^2T_1^{}X)+c_1(^3T_1^{}X).$$ In general, we have, by induction: Theorem 1.7 The first Chern number of the bundle of $`k`$-jet forms is given by the formula: $$c_1(T_k^{}X)=\underset{j=1}{\overset{k}{}}c_1(^jT_1^{}X).$$ In particular, if $`X`$ is a Riemann surface then $$c_1(T_k^{}X)=\underset{j=1}{\overset{k}{}}jc_1(T_1^{}X)=\frac{k(k+1)}{2}c_1(𝒦_X)$$ where $`𝒦_X=T_1^{}X`$ is the canonical bundle of $`X`$. Note that if $`X`$ is a Riemann surface then the rank of $`T_k^{}X`$ is $`k`$. Corollary 1.8 Let $`X`$ be a projective manifold and suppose that the cotangent bundle $`T_1^{}X`$ is ample then $`T_k^{}X`$ is ample for all $`k`$. Definition 1.9 The dual of $`J^kX`$, i.e., germs of $`\omega :j^kX\|_U𝐂`$ such that $`\omega (\lambda .j^kf)=\lambda ^m\omega (j^kf)`$ for some positive integer $`m`$, shall be referred to as the sheaf of germs of $`k`$-jet differentials and shall be denoted by $`𝒥_k^{}X`$. A jet differential $`\omega `$ satisfying the homogenity above with integer $`m`$ is said to a $`k`$-jet differential of weight $`m`$. The sheaf of $`k`$-jet differential of weight $`m`$ shall be denoted by $`𝒥_k^mX`$. It follows from the definition of the $`𝐂^{}`$-action on $`J^kX`$ that a $`k`$-jet differential $`\omega `$ of weight $`m`$ is of the form: $`\omega (j^kf)={\displaystyle \underset{\|I_1\|+2\|I_2\|+\mathrm{}+k\|I_k\|=m}{}}a_{I_1,\mathrm{},I_k}(z)(f^{^{}})^{I_1}\mathrm{}(f^{(k)})^{I_k}`$ (20) where $`a_{I_1,\mathrm{}I_k}`$ are holomorphic functions, $`I_j=(i_{1j},\mathrm{},i_{nj}),n=dimX`$ are the multi-indices with ech $`i_{lj}`$ being a non-negative integer and $`I_j\|=i_{1j}+\mathrm{}+i_{nj}`$. In terms of a local coordinate $`(z_1,\mathrm{},z_n)`$ near a point $`z`$, $$(f^{^{}})^{I_1}\mathrm{}(f^{(k)})^{I_k}=(f_1^{^{}})^{i_{11}}\mathrm{}(f_n^{^{}})^{i_{n1}}\mathrm{}(f_1^{(k)})^{i_{1k}}\mathrm{}(f_n^{(k)})^{i_{nk}}.$$ Moreover the coefficients $`a_{I_1,\mathrm{}I_k}(z)`$ are symmetric with respect to the indices in each $`I_j`$. More precisely, $$a_{(i_{\sigma _1(1)1},\mathrm{},i_{\sigma _1(n)1}),\mathrm{},(i_{\sigma _k(1)k},\mathrm{},i_{\sigma _k(n)k})}=a_{(i_{11},\mathrm{},i_{n1}),\mathrm{},(i_{1k},\mathrm{},i_{nk})}$$ where each $`\sigma _j:\{1,\mathrm{},n\}\{1,\mathrm{},n\},j=1,\mathrm{},n`$ is a permutation of $`n`$-elements. Let $`_k^mX`$ be the subsheaf of $`𝒥_k^mX`$ consisting of elements of the form: $`\omega (j^kf)={\displaystyle \underset{j=1}{\overset{k}{}}}{\displaystyle \underset{\|I_1\|+j\|I_j\|=m}{}}a_{I_1}(f)(f^{^{}})^{I_1}(f^{(j)})^{I_j}.`$ (21) Note that the coefficient of $`(f^{^{}})^{I_1}(f^{(j)})^{I_j}`$ depends only on $`I_1`$ but is independent of $`I_j`$. This sheaf shall be referred to as jet differentials of linear type. Lemma 1.10 The sheaf $`_k^mX`$ of jet differentials of linear type is well-defined. For $`m=k=2`$ we have $`_2^2X=𝒥_2^2X`$ and if $`X`$ is a Riemann surface then $`_3^3X=𝒥_3^3X`$. Proof. The change of variable formulas (4) shows that a jet differential of the form (14) is invariant by change of coordinates. QED There is a differentiatial operation $`d:𝒥_k^mX𝒥_{k+1}^{m+1}X`$ naturally defined by: $`d\omega (j^{k+1}f)\stackrel{\mathrm{def}}{=}(\omega (j^kf))^{^{}}.`$ (22) It should be noted that in contrast to exterior differentiation of forms $`dd0`$ on jet differentials. In particular, given a holomorphic function $`\varphi `$ defined on some open neighborhood $`U`$, $`d^{(k)}\varphi (j^kf)=(\varphi f)^{(k)}`$ (23) which is a non-trivial $`k`$-jet differential for general $`\varphi `$. Another difference between jet differentials and exterior differential forms is that a lower order jet differential can be naturally identified with a jet differential of higher order. More precisely, the natural projection $`p_{kj}:J^kXJ^jX`$ defined for $`kj`$ induces an injection $`p_{kj}^{}:𝒥_j^mX𝒥_k^mX`$ defined naturally by: $`p_{kj}^{}\omega (j^kf)\stackrel{\mathrm{def}}{=}\omega (p_{kj}(j^kf))=\omega (j^jf).`$ (24) We shall simply write $`\omega (j^kf)=\omega (j^jf)`$ if no confusion arises. Moreover, the symmetric product of a $`k`$-jet differential of weight $`m`$ and a $`k^{}`$-jet differential of weight $`m^{}`$ is a $`(k+k^{})`$-jet differential of weight $`m+m^{}`$. Consider first the case of $`k=2`$ and denote by $$p:J^2Xp(J^2X)T^2X$$ the generically 2 to 1 map onto its image as defined in the previous section. Let $`\omega H^0(U,T_2^{}X)`$ considered as a linear (along the fibers) functional $$\omega :T^2X\|_U𝐂.$$ Consider the composite map $$\omega p:J^2X\|_U𝐂.$$ By the definition of $`p`$ we observe that $$\omega p(\lambda .j^2f)=\omega (\lambda ^2.p(j^2f))=\lambda ^2\omega p(j^2f)$$ is homogeneous of degree 2. In other words, the composite $`\omega p`$ is a section of $`𝒥_2^2X`$ over $`U`$. Thus we have a well-defined $`𝐂^{}`$-bundle map $$q=p^{}:T_2^{}X𝒥_2^2X.$$ Consider again the special case of a Riemann surface $`X`$ and $`k=2`$. Let $`\omega H^0(U,T_2^{}X)`$ be a 2-jet form then locally $`\omega `$ is simply of the form $$\omega =adz+bd^2z^{(2)}$$ where $`dz`$ is the dual of $`/z`$ and $`d^2z^{(2)}`$ (the notation is formal and should not be confused with differentiating $`z^2`$ twice) with $`a`$ and $`b`$ being local holomorphic functions. By the definition of $`p`$ we have: $`p^{}\omega (j^2f)=a(f)f^{^{\prime \prime }}+b(f)(f^{^{}})^2.`$ (25) If $`p^{}\omega (j^2f)=a(f)f^{^{\prime \prime }}+b(f)(f^{^{}})^2=0`$ for all $`j^2f`$ then, by taking $`j^2f=(f(0),f^{^{}}(0),f^{^{\prime \prime }}(0))`$ such that $`f^{^{}}(0)=0`$ and $`f^{^{\prime \prime }}(0)0`$, we see that $`a(f(0))=0`$ (as $`x=f(0)`$ is an arbitrary point of $`U`$, we have $`a0`$) and so $`p^{}\omega (j^2f)=b(f)(f^{^{}})^2`$ for all $`j^2f`$. Now choosing $`j^2f(0)`$ so that $`f^{^{}}(0)0`$ this time shows that $`b(f(0))=0`$. In other words, the map $`p^{}:T_2^{}X𝒥_2^2X`$ is injective. On the other hand, any homogeneous polynomial of degree $`2`$ of $`J_x^2X`$ is the germ of a section of the form as in (18) where $`a`$ and $`b`$ are holomorphic functions defined on some open neighborhood of $`x`$ and that $`f(0)=x`$. This shows that $`p^{}`$ is actually an isomorphism between $`T_2^{}X`$ and $`𝒥_2^2X`$. Analogously, a section of $`^2T^{}X`$ is of the form $$\omega =b(dzdz)$$ where $`b`$ is a holomorphic function on $`U`$. Then the pull-back $$p^{}(\omega )(j^2f)=b(f)(f^{^{}})^2.$$ In other words $$p^{}(^2T^{}X)^2T^{}X$$ and hence we conclude that the pull-back of the sequence: $$0^2T^{}XT_2^{}XT^{}X0$$ yields an exact sequence: $$0p^{}(^2T^{}X)^2T^{}X𝒥_2^2X𝒥_2^2X/p^{}(^2T^{}X)T^{}X0.$$ The same argument works also for $`k=3`$; for general $`k`$ an analogous argument shows that $`p_k^{}(T_k^{}X)`$ is isomorphic to the sheaf of jet differentials of linear type $`_k^k`$: Theorem 1.11 For a complex manifold $`X`$ the pull-back $`p_k^{}(T_k^{}X)`$ is $`𝐂^{}`$-isomorphic to $`_k^kX`$ where $`p_k:J^kXT^kX`$ is the map defined by $`(11)`$. Then $`p_k^{}(T_k^{}X)`$ is $`𝐂^{}`$-isomorphic to $`_k^kX`$. Proof. We have already seen the case $`k=2`$ and suppose now that $`\omega H^0(U,T_k^{}X)`$ is a k-jet form then locally $`\omega `$ is of the form $$\omega =a_1dz+a_2d^2z^{(2)}+\mathrm{}+a_kd^kz^{(k)}$$ where $`d^jz^{(j)}`$ is the dual of the differential operator $`^j/z^j`$ and each $`a_j`$ is a holomorphic function on $`U`$. Pulling back we get: $$p_k^{}\omega (j^kf)=\underset{j=1}{\overset{k}{}}a_jf^{(j)}(f^{^{}})^{kj}$$ and suppose that $`p_k^{}\omega (j^kf)0`$. Consider first the case $`k=3`$ then for any $`xU`$ choosing $`j^3f`$ so that $`f(0)=x,f^{^{}}(0)=0`$ shows that $`a_3f^{(3)}(0)=0`$ so $`a_3(x)=0`$. Since $`x`$ is arbitray the function $`a_30`$. Thus $$0p_3^{}\omega (j^3f)=f^{^{}}\{a_2f^{^{\prime \prime }}+a_1(f^{^{}})^2\}$$ and so we have $$a_2f^{^{\prime \prime }}+a_1(f^{^{}})^20$$ on $`J^2X\|_U\{f^{^{}}0\}`$. Let $`\varphi :\mathrm{\Delta }_ϵ\mathrm{\Delta }_ϵ`$ be a holomorphic function such that $`\varphi (0)=0,\varphi ^{^{}}(0)=1`$ then $`(f\varphi )^{^{}}=f^{^{}}(\varphi )\varphi ^{^{}},(f\varphi )^{^{\prime \prime }}=f^{^{\prime \prime }}(\varphi )(\varphi ^{^{}})^2+f^{^{}}(\varphi )\varphi ^{^{\prime \prime }}`$ and the condition that $`f^{^{}}(0)0`$ implies that we may choose $`\varphi `$ so that the condition that the first jet is non-zero, (i.e., $`(f\varphi )^{^{}}(0)0`$) is preserved but $$(f\varphi )^{^{\prime \prime }}=f^{^{\prime \prime }}(\varphi )(\varphi ^{^{}})^2+f^{^{}}(\varphi )\varphi ^{^{\prime \prime }}=0$$ i.e., choose $`\varphi `$ so that $`\varphi ^{^{\prime \prime }}(0)=f^{^{\prime \prime }}(0)/f^{^{}}(0)`$. This yields: $$a_2(\varphi )(f\varphi )^{^{\prime \prime }}(0)=a_2(\varphi )(f\varphi )^{^{\prime \prime }}(0)+a_1(\varphi )((f\varphi )^{^{}})^2(0)=0$$ so $`a_20`$ (because $`x=f(0)`$ is an arbitrary point) and the original equation is reduced to the equation $`a_1(f)(f^{^{}})^30`$. Thus by choosing $`f^{^{}}(0)0`$ we conclude that $`a_1(f(0))=0`$; this implies that $`a_10`$ as well. This establishes injectivity of the map $`p_3^{}`$; surjectivity follows from the fact that an element of $`𝒥_3^3X`$ is of the form $$a_3f^{(3)}+a_2f^{^{}}f^{^{\prime \prime }}+a_1(f^{^{}})^3.$$ In general we have, by setting $`f^{^{}}=0`$, that $`a_k0`$ and then: $$0p_k^{}\omega =f^{^{}}\underset{j=1}{\overset{k1}{}}a_jf^{(j)}(f^{^{}})^{k1j}$$ and so $$0\underset{j=1}{\overset{k1}{}}a_jf^{(j)}(f^{^{}})^{k1j}$$ on $`J^kX\{f^{^{}}0\}`$. This shows injectivity; surjectivity now follows from the definition of $`_k^kX`$. The proof is then completed by induction and by reparametrization. QED The following Theorem can be found (without proof) in Green-Griffiths \[G-G\], we include a proof here for the sake of completeness: Theorem 1.12 There exists a filtration of $`𝒥_k^mX`$: $$𝒥_{k1}^mX=_k^0_k^1\mathrm{}_k^{[m/k]}=𝒥_k^mX$$ $`(`$where $`[m/k]`$ is the greatest integer smaller than or equal to $`m/k)`$ such that $$_k^i/_k^{i1}𝒥_{k1}^{mki}X(^iT^{}X).$$ Proof. The filtrations are defined as follows. Since a $`(k1)`$-jet differential of weight $`m`$ is also a $`k`$-jet differential of weight $`m`$ thus $$F_k^0=𝒥_{k1}^mX𝒥_k^mX$$ which in terms of the expression (13) for jet differentials consists of elements of which does not contain any terms involving $`f^{(k)}`$; put it another way the exponent $`I_k`$ for $`f^{(k)}`$ satisfies the condition $`\|I_k\|=0`$: $`\omega (j^kf)`$ $`=`$ $`{\displaystyle \underset{\|I_1\|+2\|I_2\|+\mathrm{}+(k1)\|I_{k1}\|=m}{}}a_{I_1,\mathrm{},I_{k1}}(z)(f^{^{}})^{I_1}\mathrm{}(f^{(k1)})^{I_{k1}}`$ $`=`$ $`{\displaystyle \underset{\|I_1\|+2\|I_2\|+\mathrm{}+k\|I_k\|=m,\|I_k\|=0}{}}a_{I_1,\mathrm{},I_k}(z)(f^{^{}})^{I_1}\mathrm{}(f^{(k)})^{I_k}.`$ For any $`0j[m/k]`$ we define $`F_j𝒥_k^mX`$ to be the sheaf of germs consisting elements so that $`\|I_k\|j`$: $`F_k^j=\{\omega \|\omega (j^kf)={\displaystyle \underset{\|I_1\|+2\|I_2\|+\mathrm{}+k\|I_k\|=m,\|I_k\|j}{}}a_{I_1,\mathrm{},I_k}(z)(f^{^{}})^{I_1}\mathrm{}(f^{(k)})^{I_k}\}.`$ (26) By definition, we have, for $`1j[m/k]`$: $`F_k^j/F_k^{j1}`$ $`=`$ $`\{\omega \|\omega (j^kf)`$ $`=`$ $`{\displaystyle \underset{\|I_1\|+2\|I_2\|+\mathrm{}+k\|I_k\|=m,\|I_k\|=j}{}}a_{I_1,\mathrm{},I_k}(z)(f^{^{}})^{I_1}\mathrm{}(f^{(k)})^{I_k}\}`$ and the claim is that $$F_k^j/F_k^{j1}𝒥_{k1}^{mkj}X(^jT^{}X).$$ We first establish the special case of a Riemann surface. In this case a $`k`$-jet differential of weight $`m`$ is of the form $$\omega (j^kf)=\underset{i_1+2i_2+\mathrm{}+ki_k=m}{}a_{i_1,\mathrm{},i_k}(z)((zf)^{^{}})^{i_1}\mathrm{}((zf)^{(k)})^{i_k}$$ where we identify $`f`$ with $`z`$ being a local coordinate on an open coordinate neighborhood $`UX`$ and $`i_j`$ are non-negative integers; the subsheaves $`F_k^j`$ is of the form: $$F_k^j=\{\omega \|\omega (j^kf)=\underset{i_1+2i_2+\mathrm{}+ki_k=m,i_kj}{}a_{i_1,\mathrm{},i_k}(z)(f^{^{}})^{i_1}\mathrm{}(f^{(k)})^{i_k}\}$$ for $`0j[m/k]`$ and $$F_k^j/F_k^{j1}=\{\omega \|\omega (j^kf)=\underset{i_1+2i_2+\mathrm{}+ki_k=m,i_k=j}{}a_{i_1,\mathrm{},i_k}(z)(f^{^{}})^{i_1}\mathrm{}(f^{(k)})^{i_k}\}$$ for $`1j[m/k]`$. We first define a map $$L_U:F_k^j/F_k^{j1}\|_U𝒥_{k1}^{mkj}X(^jT^{}X)\|_U$$ where $`L_U({\displaystyle \underset{i_1+2i_2+\mathrm{}+ki_k=m,i_k=j}{}}a_{i_1,\mathrm{},i_k}(z)(f^{^{}})^{i_1}\mathrm{}(f^{(k)})^{i_k})`$ $`=(f^{(k)})^j{\displaystyle \underset{i_1+2i_2+\mathrm{}+(k1)i_{k1}=mkj}{}}a_{i_1,\mathrm{},j}(z)(f^{^{}})^{i_1}\mathrm{}(f^{(k1)})^{i_{k1}})`$ The fact that $`L_U`$ is an isomorphism is clear and the fact that $`L=L_U`$ (where $`𝒰=\{U\}`$ is an open cover of $`X`$ by coordinate neighborhoods) follows from the following observation that (see (4)) if $`(V,w)`$ is another coordinate neighborhood then $`((wf)^{(k)})^j=((w/z)(zf)^{(k)}+P)^j=((w/z)(zf)^{(k)})^j+Q`$ where $`P`$ and $`Q`$ are polynomials in the variables $`^sw_i/z_l^s,1i,ln,1sk`$ and in $`(zf)^{(r)},1rk1`$. In particular, $`Q`$ is a $`(k1)`$-jet differential of total weight $`mkj`$. In orther words, $$((wf)^{(k)})^j=(w/z)^j((zf)^{(k)})^j\mathrm{mod}F_k^{j1}$$ and the transisistion function $`(w/z)^j`$ is the same as the transistioon function for $`^jT^{}X`$. The higher dimensional case is notationally more complicated but the proof is essentially the same. QED As an immediate consequence (see Green-Griffiths \[G-G\]) we have: Corollary 1.13 Let $`X`$ be a smooth projective variety then $`𝒥_k^mX`$ admits a composition series whose factors contain all bundles of the form: $$(^{i_1}T^{}X)\mathrm{}(^{i_k}T^{}X)$$ where $`i_j`$ ranges over all non-negative integers satisfying $$i_1+2i_2+\mathrm{}+ki_k=m.$$ The first Chern number of $`c_1(𝒥_k^mX)`$ is given by: $$c_1(𝒥_k^mX)=\underset{i_1+2i_2+\mathrm{}+ki_k=m,i_j𝐙_0}{}c_1((^{i_1}T^{}X)\mathrm{}(^{i_k}T^{}X)).$$ In particular, for a curve $`X=C`$, $`c_1(𝒥_k^mC)`$ $`=`$ $`{\displaystyle \underset{i_1+2i_2+\mathrm{}+ki_k=m,i_j𝐙_0}{}}(i_1+i_2+\mathrm{}+i_k)c_1(T^{}C).`$ The preceding Theorem can be used in calculating the Chern classes of $`𝒥_k^mX`$. Example 1.14 For example, for $`m=k=2`$, the filtration is given by: $$^2T^{}X=𝒥_1^2X=𝒮_2^0𝒮_2^1=𝒥_2^2X,𝒮_2^1/𝒮_2^0T^{}X$$ we have the following exact sequence: $$0^2T^{}X𝒥_2^2XT^{}X0.$$ Thus the first Chern numbers are related by the formula: $$c_1(𝒥_2^2X)=c_1(^2T^{}X)+c_1(T^{}X).$$ The filtration of $`𝒥_3^3X`$ is as follows: $$𝒥_3^3X=𝒮_3^1S_3^0=𝒥_2^3X,𝒥_3^3X/𝒥_2^3X=𝒮_3^1/𝒮_3^0T^{}X.$$ Hence we have an exact sequence: $`0𝒥_2^3X𝒥_3^3XT^{}X0.`$ Now the filtration of $`𝒥_2^3X`$ is $$𝒥_2^3X=𝒮_2^1𝒮_2^0=𝒥_1^3X,𝒥_2^3X/𝒥_1^3XT^{}XT^{}X$$ and, since $`𝒥_1^3X=^3T^{}X`$, we have an exact sequence: $`0^3T^{}X𝒥_2^3XT^{}XT^{}X0.`$ From these 2 exact sequences we get $`c_1(𝒥_3^3X)=c_1(T^{}X)+c_1(T^{}XT^{}X)+c_1(^3T^{}X).`$ From basic representation Theory (or just simple liner algebra in this special case) we know that $`T^{}XT^{}X=^2T^{}X^2T^{}X`$ hence, $`c_1(𝒥_3^3X)=c_1(T^{}X)+c_1(^2T^{}X)+c_1(^3T^{}X)+c_1(^2T^{}X).`$ In representation theory $`^2T^{}X`$ is the Weyl module $`W_{1,1}^{}X`$ associate to the partition $`\{1,1\}`$ (see \[F-H\]). Thus we have: $`c_1(𝒥_3^3X)={\displaystyle \underset{j=1}{\overset{3}{}}}c_1(^jT^{}X)+c_1(W_{1,1}^{}X).`$ (27) In the special case of a Riemann surface $`^2T^{}X`$ is the zero-sheaf. Thus for a curve we have $$c_1(𝒥_3^3X)=(1+2+3)c_1(T^{}X)=6c_1(T^{}X).$$ For $`m=k=4`$, we have the following filtrations: $$𝒥_4^4X=S_4^1S_4^0=𝒥_3^4X,𝒥_4^4X/𝒥_3^4X=S_4^1/S_4^0T^{}X,$$ $$𝒥_3^4X=S_3^1S_3^0=𝒥_2^4X,𝒥_3^4X/𝒥_2^4X=S_3^1/S_3^0T^{}XT^{}X,$$ and $$𝒥_2^4X=S_2^2S_2^1S_2^0=𝒥_1^4X,$$ with $$𝒥_2^4X/𝒮_2^1=^2T^{}X,S_2^1/S_2^0T^{}X(^2T^{}X).$$ The exact sequences associate to the filtration for $`𝒥_4^4X`$ are: $`0𝒥_3^4X𝒥_4^4XT^{}X0;`$ $`0𝒥_2^4X𝒥_3^4XT^{}XT^{}X0;`$ $`0S_1𝒥_2^4X^2T^{}X0;`$ $`0^4T^{}XS_1T^{}X(^2T^{}X)0.`$ Thus the Chern number formula: $`c_1(𝒥_4^4X)`$ $`=`$ $`c_1(T^{}X)+c_1(T^{}XT^{}X)+c_1(^2T^{}X)`$ $`+c_1(T^{}X(^2T^{}X))+c_1(^4T^{}X).`$ Note that (by elementary representation theory) $$T^{}X(^kT^{}X)=W_{k,1}^{}X(^{k+1}T^{}X)$$ where $`W_{k,1}^{}`$ is the Weyl module associate to the partition $`\{k,1\}`$ thus: $`c_1(𝒥_4^4X)=c_1(^2T^{}X)+{\displaystyle \underset{i=1}{\overset{4}{}}}c_1(^iT^{}X)+{\displaystyle \underset{i=1}{\overset{2}{}}}c_1(W_{j,1}^{}X).`$ (28) In particulr, if $`X`$ is a curve then $$c_1(𝒥_4^4X)=(1+2+2+3+4)c_1(T^{}X)=12c_1(T^{}X).$$ Recall that $`c_1(T_4^{}X)=10c_1(T^{}X)`$. For $`m=k=5`$, we have the following filtrations: $$𝒥_5^5X=S_5^1S_5^0=𝒥_4^5X,𝒥_5^5X/𝒥_4^5X=S_5^1/S_5^0T^{}X,$$ $$𝒥_4^5X=S_4^1S_4^0=𝒥_3^5X,𝒥_4^5X/𝒥_3^5X=S_4^1/S_4^0T^{}XT^{}X,$$ $$𝒥_3^5X=S_3^1S_3^0=𝒥_2^5X,𝒥_3^5X/𝒥_2^5XT^{}X(𝒥_2^2X),$$ $$𝒥_2^5X=S_2^2S_2^1S_2^0=𝒥_1^5X,𝒥_2^5X/𝒮_2^1=(^2T^{}X)T^{}X,$$ $$S_2^1/S_2^0T^{}X(^3T^{}X).$$ The exact sequences associate to the filtration for $`𝒥_5^5X`$ are: $`0𝒥_4^5X𝒥_5^5XT^{}X0;`$ $`0𝒥_3^5X𝒥_4^5XT^{}XT^{}X0;`$ $`0𝒥_2^5X𝒥_3^5XT^{}X𝒥_2^2X0;`$ $`0𝒮_2^1𝒥_2^5X(^2T^{}X)T^{}X0,`$ $`0^5T^{}X𝒮_2^1T^{}X(^3T^{}X)0.`$ This yields the formula: $`c_1(𝒥_5^5X)`$ $`=c_1(T^{}X)+c_1(T^{}XT^{}X)+c_1(T^{}X𝒥_2^2X)`$ $`+c_1((^2T^{}X)T^{}X)+c_1(T^{}X(^3T^{}X))+c_1(^5T^{}X)`$ $`=c_1(T^{}X)+c_1(T^{}XT^{}X)+c_1(T^{}X(^2T^{}X))`$ $`+c_1(T^{}XT^{}X)+c_1((^2T^{}X)T^{}X)`$ $`+c_1(T^{}X(^3T^{}X))+c_1(^5T^{}X)`$ where we have used the fact that $$c_1(T^{}X𝒥_2^2X)=c_1(T^{}X(^2T^{}X))+c_1(T^{}XT^{}X).$$ Recall that $$T^{}XT^{}X=^2T^{}X^2T^{}X,$$ $$T^{}X(^2T^{}X)=^3T^{}XW_{2,1}^{}X,$$ $$T^{}X(^3T^{}X)=^4T^{}XW_{3,1}^{}X$$ (in general we have $$T^{}X(^dT^{}X)=^{d+1}T^{}XW_{d,1}^{}X.)$$ Thus we have: $`c_1(𝒥_5^5X)`$ $`={\displaystyle \underset{j=2}{\overset{3}{}}}c_1(^jT^{}X)+{\displaystyle \underset{j=1}{\overset{5}{}}}c_1(^jT^{}X)+{\displaystyle \underset{j=1}{\overset{2}{}}}c_1(W_{j,1}^{}X)+{\displaystyle \underset{j=1}{\overset{3}{}}}c_1(W_{j,1}^{}X).`$ In particulr, if $`X`$ is a curve then $$c_1(𝒥_5^5X)=(1+2+3+2+3+4+5)c_1(T^{}X)=20c_1(T^{}X).$$ For $`m=k=6`$, we have the following filtrations: $$𝒥_6^6X=S_6^1S_6^0=𝒥_5^6X,𝒥_6^6X/𝒥_5^6X=S_6^1/S_6^0T^{}X,$$ $$𝒥_5^6X=S_5^1S_5^0=𝒥_4^6X,𝒥_5^6X/𝒥_4^6X=S_5^1/S_5^0T^{}XT^{}X,$$ $$𝒥_4^6X=S_4^1S_4^0=𝒥_3^6X,𝒥_4^6X/𝒥_3^6XT^{}X𝒥_2^3X,$$ $$𝒥_3^6X=S_3^2S_3^1S_3^0=𝒥_2^6X,$$ with factors $$𝒥_3^6X/𝒮_3^1=^2T^{}X,S_3^1/S_3^0T^{}X𝒥_2^3X,$$ $$𝒥_2^6X=S_2^3S_2^2S_2^1S_2^0=^6T^{}X,𝒥_2^6X/𝒮_2^2=^3T^{}X,$$ $$S_2^2/S_2^1(^2T^{}X)(^2T^{}X),S_2^1/S_2^0T^{}X(^4T^{}X).$$ The exact sequences associate to the filtration for $`𝒥_6^6X`$ are: $`0𝒥_5^6X𝒥_6^6XT^{}X0;`$ $`0𝒥_4^6X𝒥_5^6XT^{}XT^{}X0;`$ $`0𝒥_3^6X𝒥_4^6XT^{}X𝒥_2^3X0;`$ $`0𝒮_3^1𝒥_3^6X^2T^{}X0,`$ $`0𝒥_2^6X𝒮_3^1T^{}X𝒥_2^3X0,`$ $`0𝒮_2^2𝒥_2^6X^3T^{}X0,`$ $`0𝒮_2^1𝒮_2^2(^2T^{}X)(^2T^{}X)0,`$ $`0^6T^{}X𝒮_2^1T^{}X(^4T^{}X)0.`$ This yields the formula: $`c_1(𝒥_6^6X)`$ $`=`$ $`c_1(T^{}X)+c_1(T^{}XT^{}X)+2c_1(T^{}X𝒥_2^3X)`$ $`+c_1(^2T^{}X)+c_1(^3T^{}X)+c_1((^2T^{}X)(^2T^{}X))`$ $`+c_1(T^{}X(^4T^{}X))+c_1(^6T^{}X)`$ $`=`$ $`c_1(T^{}X)+3c_1(T^{}XT^{}X)+2c_1(T^{}X(^3T^{}X))`$ $`+c_1(^2T^{}X)+c_1(^3T^{}X)+c_1((^2T^{}X)(^2T^{}X))`$ $`+c_1(T^{}X(^4T^{}X))+c_1(^6T^{}X).`$ where we have used the fact that $$c_1(T^{}X𝒥_2^3X))=c_1(T^{}X(^3T^{}X))+c_1(T^{}XT^{}X).$$ In particulr, if $`X`$ is a curve then $$c_1(𝒥_6^6X)=(1+6+8+2+3+4+5+6)c_1(T^{}X)=35c_1(T^{}X).$$ The calculation before can be carried out in a much simpler fashion as follows. A partition of a natural number $`m`$ is a set of positve integers $`k_1,\mathrm{},k_q`$ such that $`m=k_1+\mathrm{}+k_q`$. A partition can be expressed as $$m=\underset{j=1}{\overset{k}{}}ji_j$$ where the integers $`i_j=\mathrm{\#}`$ of $`j`$’s in $`\{k_1,\mathrm{},k_q\}`$ are non-negative. Obviously we have $`1qk`$ and $`1k_ik`$ for all $`i`$. The following result is well-known in representation theory and in combinatorics (see \[H-W\]): Theorem 1.15 The number, denoted $`p(m)`$, of classes of $`S_m(`$the symmetric group on $`m`$ elements$`)`$ is equal to the number of partitions of $`m`$ and also to the number of $`(`$inequivalent$`)`$ irreducible representations of $`S_m`$. The number $`p(m)`$ is asymptotically approximated by the formula of Hardy-Ramanujan $$p(m)\frac{e^{\pi \sqrt{2m/3}}}{4m\sqrt{3}}.$$ Remark 1.16 The first few numbers are as follows: $$p(1)=1,p(2)=2,p(3)=3,p(4)=5,p(5)=7,p(6)=11,p(7)=15,$$ $$p(8)=22,p(9)=30,p(10)=42,p(11)=56,p(12)=77,p(13)=101.$$ We are interested in the case $$k=\lambda _1+\lambda _2+\mathrm{}+\lambda _{\rho _\lambda }$$ where $`\lambda _1\lambda _2\mathrm{}\lambda _{\rho _\lambda }1`$. Define $`l_i=\lambda _i+\rho _\lambda i,i=1,\mathrm{},\rho _\lambda `$. Then the dimension $`d_\lambda `$ of the representation $`V_\lambda ,\lambda =(\lambda _1,\mathrm{},\lambda _{\rho _\lambda })`$ associated to the partition $`\lambda `$ is given by the formula $`d_\lambda =1`$ if $`\rho _\lambda =1`$ and for $`\rho _\lambda 1`$ (see \[F-H\], p. 50): $`d_\lambda ={\displaystyle \frac{k!}{l_1!\mathrm{}l_{\rho _\lambda }!}}{\displaystyle \underset{1i<j\rho _\lambda }{}}(l_il_j)`$ (29) The number $`\rho _\lambda `$ shall be referred to as the length of the partition $`\lambda `$. We consider also the case of partitioning a number by a partition of fixed length $`k`$. Denote by $`p_k(m)`$ the number of solutions of the equation $`x_1+\mathrm{}+x_k=m`$ with the condition that $`1x_kx_{k1}\mathrm{}x_1`$. This number is obviously equal to the number of solutions of the equation $`y_1+\mathrm{}+y_k=mk`$ with the condition that $`0y_ky_{k1}\mathrm{}y_1`$. If there are exactly $`i`$ of the integers $`y_i`$ which are positive then these are the solutions of $`x_1+\mathrm{}+x_i=mk`$ ($`x_iy_i+1`$) and so there are $`p_i(mk)`$ of such solutions; consequently we have: Theorem 1.17 With the notations above we have $$p_k(m)=\underset{i=0}{\overset{k}{}}p_i(mk)$$ if $`1km`$ and with the convention that $`p_0(0)=1,p_0(m)=0`$ if $`m>0`$ and $`p_k(m)=0`$ if $`k>m`$. The following identity is easily established by induction: Theorem 1.18 The number $`p_k(m)`$ satisfies the following recursive relation: $`p_k(m)=p_{k1}(m1)+p_k(mk)`$. Obviously we have $`p_1(m)=p_m(m)=1`$ and $`p_2(m)=m/2`$ or $`(m1)/2`$ according to $`m`$ being even or odd. Thus Theorem 1.19 yields $`p_3(m)=p_2(m1)+p_3(m3)`$, $`p_4(m)=p_3(m1)+p_4(m4),`$ $`p_5(m)=p_4(m1)+p_5(m5)`$ and we get for example $$p_1(6)=1,p_2(6)=3,p_6(6)=1$$ $$p_3(6)=p_2(5)+p_3(3)=3,$$ $$p_4(6)=p_3(5)=p_2(4)=2,$$ $$p_5(6)=p_4(5)=p_3(4)=p_2(3)=1$$ hence as $`p(m)=_kp_k(m)`$ we have $$p(6)=\underset{k=1}{\overset{6}{}}p_k(6)=1+3+3+2+1+1=11.$$ For $`m=7`$ we have $$p_1(7)=1,p_2(7)=3,p_7(7)=1$$ $$p_3(7)=p_2(6)+p_3(4)=p_2(6)+p_2(3)=4,$$ $$p_4(7)=p_3(6)=3,$$ $$p_5(7)=p_4(6)=2,$$ $$p_6(7)=p_5(6)=1$$ $$p(7)=\underset{k=1}{\overset{7}{}}p_k(7)=1+3+4+3+2+1+1=15.$$ The total length of all partitions $`L(m)`$ of a positive integer $`m`$ is defined to be $`L(m)={\displaystyle \underset{j=1}{\overset{m}{}}}jp_j(m).`$ For example if $`m=6`$ then $`L(6)=1+6+9+8+5+6=35`$ and for $`m=7,L(7)=1+6+12+12+10+6+7=54.`$ More generally for $`km`$ $`L_k(m)={\displaystyle \underset{j=1}{\overset{k}{}}}jp_j(m)`$ (30) shall be referred to as the total length of partitions of $`m`$ of length at most $`k`$. The following Lemma is easily established from the definitions: Lemma 1.19 With the notations above we have $$\underset{\lambda }{}\underset{j=1}{\overset{k}{}}i_j=\underset{\lambda }{}\rho _\lambda =\underset{j=1}{\overset{k}{}}jp_j(m)$$ where the sum on the right is taken over all partition $`\lambda =(\lambda _1,\mathrm{},\lambda _{\rho _\lambda })`$ of $`m,1\lambda _{\rho _\lambda }\mathrm{}\lambda _2\lambda _1,\rho _\lambda k`$ and $`i_j=\mathrm{\#}`$ of $`j`$’s in $`\{\lambda _1,\mathrm{},\lambda _{\rho _\lambda }\}`$. One has the following well-known asymptotic formula: Theorem 1.20 For $`m\mathrm{}`$ the number $`p_k(m)`$ is asymptotically given by: $$p_k(m)\frac{m^{k1}}{(k1)!k!}.$$ The preceding discussions yield the following Theorem: Theorem 1.21 Let $`X`$ be a non-singular pojective curve then the Chern number of $`𝒥_k^mX`$ is given by $$c_1(𝒥_k^mX)=L_k(m)c_1(𝒦_X)=\underset{j=1}{\overset{k}{}}jp_j(m)c_1(𝒦_X)=\underset{j=1}{\overset{k}{}}jp_j(m)c_1(𝒦_X)$$ where $`𝒦_X`$ is the canonical bundle of $`X`$. If we fix $`k`$ and let $`m\mathrm{}`$ then asymptotically: $$c_1(𝒥_k^mX)kp_k(m)\frac{m^{k1}}{(k1)!(k1)!}.$$ We give as examples the explicit calculation of the above. For $`m=k=3`$, we have $`p(3)=3`$ and the possible indices are tabulated below: \| \| $`\lambda `$ \| $`\rho _\lambda `$ \| $`d_\lambda `$ \| $`i_1`$ \| $`i_2`$ \| $`i_3`$ \| $`_{j=1}^ki_j`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| 1 \| (1, 1, 1) \| 3 \| 1 \| 3 \| 0 \| 0 \| 3 \| \| 2 \| (2, 1) \| 2 \| 2 \| 1 \| 1 \| 0 \| 2 \| \| 3 \| (3) \| 1 \| 1 \| 0 \| 0 \| 1 \| 1 \| The Chern number of a curve $`X`$ is obtained by summing the last column: $$c_1(𝒥_3^3X)=(1+2+3)c_1(T^{}X)=6c_1(T^{}X).$$ For $`m=k=4`$, we have $`p(4)=5`$ and the possible indices are listed below \| \| $`\lambda `$ \| $`\rho _\lambda `$ \| $`d_\lambda `$ \| $`i_1`$ \| $`i_2`$ \| $`i_3`$ \| $`i_4`$ \| $`_{j=1}^ki_j`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| 1 \| (1, 1, 1, 1) \| 4 \| 1 \| 4 \| 0 \| 0 \| 0 \| 4 \| \| 2 \| (2, 1, 1) \| 3 \| 3 \| 2 \| 1 \| 0 \| 0 \| 3 \| \| 3 \| (3, 1) \| 2 \| 3 \| 1 \| 0 \| 1 \| 0 \| 2 \| \| 4 \| (2, 2) \| 2 \| 2 \| 0 \| 2 \| 0 \| 0 \| 2 \| \| 5 \| (4) \| 1 \| 1 \| 0 \| 0 \| 0 \| 1 \| 1 \| The Chern numberof a curve $`X`$ is obtained by summing the last column: $$c_1(𝒥_4^4X)=12c_1(T^{}X).$$ For $`m=k=5`$, we have $`p(5)=7`$ and the possible indices are listed below \| \| $`\lambda `$ \| $`\rho _\lambda `$ \| $`d_\lambda `$ \| $`i_1`$ \| $`i_2`$ \| $`i_3`$ \| $`i_4`$ \| $`i_5`$ \| $`_{j=1}^ki_j`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| 1 \| (1, 1, 1, 1, 1) \| 5 \| 1 \| 5 \| 0 \| 0 \| 0 \| 0 \| 5 \| \| 2 \| (2, 1, 1, 1) \| 4 \| 4 \| 3 \| 1 \| 0 \| 0 \| 0 \| 4 \| \| 3 \| (3, 1, 1) \| 3 \| 6 \| 2 \| 0 \| 1 \| 0 \| 0 \| 3 \| \| 4 \| (2, 2, 1) \| 3 \| 5 \| 1 \| 2 \| 0 \| 0 \| 0 \| 3 \| \| 5 \| (4, 1) \| 2 \| 4 \| 1 \| 0 \| 0 \| 1 \| 0 \| 2 \| \| 6 \| (3, 2) \| 2 \| 15 \| 0 \| 1 \| 1 \| 0 \| 0 \| 2 \| \| 7 \| (5) \| 1 \| 1 \| 0 \| 0 \| 0 \| 0 \| 1 \| 1 \| The Chern numberof a curve $`X`$ is obtained by summing the last column: $$c_1(𝒥_5^5X)=20c_1(T^{}X).$$ For $`m=k=6`$, we have $`p(6)=11`$ and the possible indices are listed below \| \| $`\lambda `$ \| $`\rho _\lambda `$ \| $`d_\lambda `$ \| $`i_1`$ \| $`i_2`$ \| $`i_3`$ \| $`i_4`$ \| $`i_5`$ \| $`i_6`$ \| $`_{j=1}^ki_j`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| 1 \| (1, 1, 1, 1, 1, 1) \| 6 \| 1 \| 6 \| 0 \| 0 \| 0 \| 0 \| 0 \| 6 \| \| 2 \| (2, 1, 1, 1, 1) \| 5 \| 5 \| 4 \| 1 \| 0 \| 0 \| 0 \| 0 \| 5 \| \| 3 \| (3, 1, 1, 1) \| 4 \| 10 \| 3 \| 0 \| 1 \| 0 \| 0 \| 0 \| 4 \| \| 4 \| (2, 2, 1, 1) \| 4 \| 9 \| 2 \| 2 \| 0 \| 0 \| 0 \| 0 \| 4 \| \| 5 \| (4, 1, 1) \| 3 \| 10 \| 2 \| 0 \| 0 \| 1 \| 0 \| 0 \| 3 \| \| 6 \| (3, 2, 1) \| 3 \| 36 \| 1 \| 1 \| 1 \| 0 \| 0 \| 0 \| 3 \| \| 7 \| (2, 2, 2) \| 3 \| 5 \| 0 \| 3 \| 0 \| 0 \| 0 \| 0 \| 3 \| \| 8 \| (5, 1) \| 2 \| 30 \| 1 \| 0 \| 0 \| 0 \| 1 \| 0 \| 2 \| \| 9 \| (4, 2) \| 2 \| 9 \| 0 \| 1 \| 0 \| 1 \| 0 \| 0 \| 2 \| \| 10 \| (3, 3) \| 2 \| 5 \| 0 \| 0 \| 2 \| 0 \| 0 \| 0 \| 2 \| \| 11 \| (6) \| 1 \| 1 \| 0 \| 0 \| 0 \| 0 \| 0 \| 1 \| 1 \| The Chern numberof a curve $`X`$ is obtained by summing the last column: $$c_1(𝒥_6^6X)=35c_1(T^{}X).$$ For $`m=k=7`$, we have $`p(7)=15`$ and the possible indices are listed below \| \| $`\lambda `$ \| $`\rho _\lambda `$ \| $`d_\lambda `$ \| $`i_1`$ \| $`i_2`$ \| $`i_3`$ \| $`i_4`$ \| $`i_5`$ \| $`i_6`$ \| $`i_7`$ \| $`_{j=1}^ki_j`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| 1 \| (1, 1, 1, 1, 1, 1, 1) \| 7 \| 1 \| 7 \| 0 \| 0 \| 0 \| 0 \| 0 \| 0 \| 7 \| \| 2 \| (2, 1, 1, 1, 1, 1) \| 6 \| 6 \| 5 \| 1 \| 0 \| 0 \| 0 \| 0 \| 0 \| 6 \| \| 3 \| (3, 1, 1, 1, 1) \| 5 \| 15 \| 4 \| 0 \| 1 \| 0 \| 0 \| 0 \| 0 \| 5 \| \| 4 \| (2, 2, 1, 1, 1) \| 5 \| 14 \| 3 \| 2 \| 0 \| 0 \| 0 \| 0 \| 0 \| 5 \| \| 5 \| (4, 1, 1, 1) \| 4 \| 20 \| 3 \| 0 \| 0 \| 1 \| 0 \| 0 \| 0 \| 4 \| \| 6 \| (3, 2, 1, 1) \| 4 \| 35 \| 2 \| 1 \| 1 \| 0 \| 0 \| 0 \| 0 \| 4 \| \| 7 \| (2, 2, 2, 1) \| 4 \| 14 \| 1 \| 3 \| 0 \| 0 \| 0 \| 0 \| 0 \| 4 \| \| 8 \| (5, 1, 1) \| 3 \| 15 \| 2 \| 0 \| 0 \| 0 \| 1 \| 0 \| 0 \| 3 \| \| 9 \| (4, 2, 1) \| 3 \| 35 \| 1 \| 1 \| 0 \| 1 \| 0 \| 0 \| 0 \| 3 \| \| 10 \| (3, 3, 1) \| 3 \| 21 \| 1 \| 0 \| 2 \| 0 \| 0 \| 0 \| 0 \| 3 \| \| 11 \| (3, 2, 2) \| 3 \| 21 \| 0 \| 2 \| 1 \| 0 \| 0 \| 0 \| 0 \| 3 \| \| 12 \| (6, 1) \| 2 \| 6 \| 1 \| 0 \| 0 \| 0 \| 0 \| 1 \| 0 \| 2 \| \| 13 \| (5, 2) \| 2 \| 14 \| 0 \| 1 \| 0 \| 0 \| 1 \| 0 \| 0 \| 2 \| \| 14 \| (4, 3) \| 2 \| 14 \| 0 \| 0 \| 1 \| 1 \| 0 \| 0 \| 0 \| 2 \| \| 15 \| (7) \| 1 \| 1 \| 0 \| 0 \| 0 \| 0 \| 0 \| 0 \| 1 \| 1 \| The Chern numberof a curve $`X`$ is obtained by summing the last column: $$c_1(𝒥_7^7X)=54c_1(T^{}X).$$ We list below the next few values of $`L(k)`$: $$L(8)=86,L(9)=128,L(10)=192,L(11)=275,L(12)=399,L(13)=556$$ $$L(14)=780,L(15)=1068,L(16)=1463.$$ § 2 Computation of Chern Classes in Complex Surfaces We now treat the case of a complex surface (i.e., complex dimension 2). First we establish some basic facts: Lemma 2.1 Let $`X`$ be a nonsingular complex surface then $`c_1(^mT^{}X)={\displaystyle \frac{m(m+1)}{2}}c_1(T^{}X),`$ $`c_2(^mT^{}X)=a(m)c_1^2(T^{}X)+b(m)c_2(T^{}X)`$ where $`a(m)=m(m^21)(3m+2)/24,b(m)=m(m+1)(m+2)/6`$. Proof. the case of the first Chern class is straight forward and the calculation is omitted (see section 4 for a slightly more general calculation. To compute the Chern numbers of $`^2E`$ we proceed formally by writing the total Chern class $`c(E)=(1+(\lambda _1+\lambda _2)x+\lambda _1\lambda _2x^2)`$ then the total Chern class of $`^2E`$ is (keep in mind that rank $`^2E=3`$): $$(1+2\lambda _1x)(1+2\lambda _2x)(1+(\lambda _1+\lambda _2)x)$$ and a calculation (mod $`x^3`$) yields: $$1+3(\lambda _1+\lambda _2)x+[4\lambda _1\lambda _2+2(\lambda _1+\lambda _2)^2]x^2.$$ This shows that $`c_1(^2E)=3c_1(E),c_2(^2E)=2c_1^2(E)+4c_2(E)`$ and the Lemma is verified in this case. Next we compute the Chern numbers of $`^3E`$. With a similar formalism (and keep in mind that the rank of $`^3E`$ is 4), we have: $`c(^3E)`$ $`=`$ $`(1+3\lambda _1x)(1+3\lambda _2x)(1+(2\lambda _1+\lambda _2)x)(1+(\lambda _1+2\lambda _2)x)`$ $`=`$ $`1+6(\lambda _1+\lambda _2)x+\{11(\lambda _1+\lambda _2)^2+10\lambda _1\lambda _2\}x^2(\mathrm{mod}x^3).`$ This shows that $$c_1(^3E)=6c_1(E),c_2(^3E)=11c_1^2(E)+10c_2(E).$$ For general $`m`$ we observe that $`c_1(^mE)=\{\begin{array}{cc}(mp_2(m)+\frac{m}{2})c_1(E),\hfill & \text{if }m\text{ is even,}\hfill \\ (mp_2(m)+m)c_1(E),\hfill & \text{if }m\text{ is odd}\hfill \end{array}`$ where $`p_2(m)`$ is the number of solutions of $`m`$ with partitions of fixed length 2 defined above. The Lemma follows by recalling that $`p_2(m)=m/2`$ (resp. $`(m1)/2`$) if $`m`$ is even (resp. odd). For the tensor product we observe that $$c_1(^mE)=\underset{i=0}{}\frac{i+(mi)}{2}C_i^mc_1(E)$$ which follows from the fact that a partition $`m`$ of length 2 can be written simply as $`l=(l_1=i,l_2=mi)`$. Previously, this was defined by requiring that $`l_1l_21`$ but in the preceeding formula we include all parititions $`l=(l_1,l_2),l_i0,l_1+l_2=m`$. This accounts for the extra term $`m/2`$ (resp. $`m`$) in the formula for the symmetric product and also the factor $`1/2`$ in the formula for tensor product. The second Chern class is somewhat more complicated. Given an integer $`m`$ the non-negative partitions of $`m`$ of length 2 are $`\{(mi,i),i=0,\mathrm{},m\}`$ and $$c(^mE)=\underset{i=0}{\overset{m}{}}(1+((mi)\lambda _1+i\lambda _2)x)(\mathrm{mod}x^3).$$ The coefficients of $`x^2`$ is the second Chern class and is given by the following sums if $`m`$ is even: $$s_0=\underset{i=0}{\overset{\frac{m}{2}1}{}}((mi)\lambda _1+i\lambda _2)(i\lambda _1+(mi)\lambda _2)$$ $$s_1=\underset{i=1}{\overset{m1}{}}\{m\lambda _1((mi)\lambda _1+i\lambda _2)+m\lambda _2(i\lambda _1+(mi)\lambda _2)\}$$ $`s_2`$ $`=`$ $`{\displaystyle \underset{i=2}{\overset{m2}{}}}\{((m1)\lambda _1+\lambda _2)((mi)\lambda _1+i\lambda _2)+`$ $`+(\lambda _1+(m1)\lambda _2)(i\lambda _1+(mi)\lambda _2)\}`$ $$\mathrm{}$$ $$\mathrm{}$$ $$\mathrm{}$$ $`s_j`$ $`=`$ $`{\displaystyle \underset{i=j}{\overset{mj}{}}}\{((mj)\lambda _1+j\lambda _2)((mi)\lambda _1+i\lambda _2)+`$ $`+(j\lambda _1+(mj)\lambda _2)(i\lambda _1+(mi)\lambda _2)\}`$ $$\mathrm{}$$ $$\mathrm{}$$ $$\mathrm{}$$ $$s_{(m/2)1}=\{(\frac{m}{2}+1)\lambda _1+(\frac{m}{2}1)\lambda _2)\}\{\frac{m}{2}\lambda _1+\frac{m}{2}\lambda _2)\}$$ $$c_2(^mE)=s_0+s_1+\mathrm{}+s_{(m/2)1}.$$ By simple algebra, we have $`s_0`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{\frac{m}{2}1}{}}}i(mi)(\lambda _1^2+\lambda _2^2)+{\displaystyle \underset{i=0}{\overset{\frac{m}{2}1}{}}}(i^2+(mi)^2)\lambda _1\lambda _2`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{\frac{m}{2}1}{}}}i(mi)(\lambda _1+\lambda _2)^2+{\displaystyle \underset{i=0}{\overset{\frac{m}{2}1}{}}}(m2i)^2\lambda _1\lambda _2`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{\frac{m}{2}1}{}}}i(mi)c_1(E)^2+{\displaystyle \underset{i=0}{\overset{\frac{m}{2}1}{}}}(m2i)^2c_2(E).`$ For $`s_j,j1`$ the main observation is that each of these can be expressed as $`(\lambda _1+\lambda _2)^2`$ and so invoves only $`c_1^2`$, indeed we have, for $`1j(m/2)1`$, $`s_j`$ $`=`$ $`{\displaystyle \underset{i=j}{\overset{mj}{}}}(m^2m(i+j)+2ij)(\lambda _1+\lambda _2)^2`$ $`=`$ $`{\displaystyle \underset{i=j}{\overset{mj}{}}}(m^2m(i+j)+2ij)c_1(E)^2.`$ If $`m`$ is odd: $$s_0=\underset{i=0}{\overset{\frac{m1}{2}}{}}((mi)\lambda _1+i\lambda _2)(i\lambda _1+(mi)\lambda _2)$$ and $`s_j,j=1,\mathrm{},\frac{m1}{2}`$ are defined as before with $`c_2(^mE)=s_0+s_1+\mathrm{}+s_{\frac{m1}{2}}.`$ By simple algebra, we have $`s_0={\displaystyle \underset{i=0}{\overset{\frac{m1}{2}}{}}}i(mi)c_1(E)^2+{\displaystyle \underset{i=0}{\overset{\frac{m1}{2}}{}}}(m2i)^2c_2(E).`$ Thus $`c_2(^mE)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{\frac{m}{2}1}{}}}(m2i)^2c_2(E)+{\displaystyle \underset{i=0}{\overset{\frac{m}{2}1}{}}}i(mi)c_1(E)^2+`$ $`+{\displaystyle \underset{j=1}{\overset{\frac{m}{2}1}{}}}{\displaystyle \underset{i=j}{\overset{mj}{}}}(m^2m(i+j)+2ij)c_1(E)^2`$ if $`m`$ is even and $`c_2(^mE)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{\frac{m1}{2}}{}}}(m2i)^2c_2(E)+{\displaystyle \underset{i=0}{\overset{\frac{m1}{2}}{}}}i(mi)c_1(E)^2+`$ $`+{\displaystyle \underset{j=1}{\overset{\frac{m1}{2}}{}}}{\displaystyle \underset{i=j}{\overset{mj}{}}}(m^2m(i+j)+2ij)c_1(E)^2`$ if $`m`$ is odd. The Lemma follows by simplifying the preceding formulas. QED Lemma 2.2 Let $`E_i,i=1,\mathrm{},k`$ be holomorphic vector bundles, of rank $`r_i`$ respectively, over a non-singular complex surface $`X`$ then $`(i)c_1(_{i=1}^kE_i)={\displaystyle \underset{i=1}{\overset{k}{}}}(r_1\mathrm{}r_{i1}r_{i+1}\mathrm{}r_k)c_1(E_i),`$ $`(ii)c_2(_{i=1}^kE_i)={\displaystyle \underset{i=1}{\overset{k}{}}}r_i{\displaystyle \underset{i=1}{\overset{k}{}}}({\displaystyle \frac{c_2(E_i)}{r_i}}{\displaystyle \frac{c_1^2(E_i)}{2r_i}})+{\displaystyle \frac{\underset{i=1}{\overset{k}{}}r_i^2}{2}}{\displaystyle \underset{i=1}{\overset{k}{}}}({\displaystyle \frac{c_1(E_i)}{r_i}})^2.`$ Proof. Consider first the case $`k=2`$ then by expressing formally $`E_1=L_1\mathrm{}L_{r_1},E_2=F_1\mathrm{}F_{r_2}`$ as direct sums of line bundles we get $$E_1E_2=\underset{i=1}{\overset{r_1}{}}L_i(F_1\mathrm{}F_{r_2})$$ hence the first Chern class is given by $$c_1(E_1E_2)=\underset{i=1}{\overset{r_1}{}}(r_2c_1(L_i)+c_1(E_2))=\underset{i=1}{\overset{r1}{}}r_2c_1(E_1)+r_1c_2(E_2).$$ The case of general $`k`$ is similar. For $`c_2(E_1E_2)`$ we have $`c_2(E_1E_2)`$ $`=`$ $`c_2({\displaystyle \underset{i=1}{\overset{r_1}{}}}L_iE_2)`$ $`=`$ $`{\displaystyle \underset{i<j}{}}c_1(L_iE_2)c_1(L_jE_2)+{\displaystyle \underset{i}{}}c_2(L_iE_2).`$ The formula of the Lemma follows from the above and the following formulas, $$c_l(L_iE_2)=\underset{p=0}{\overset{l}{}}C_{lp}^{r_2p}c_1^{lp}(L_p)c_i(E_2).$$ The calculation of the general case is achieved via induction. QED The next formulas are consequences of the preceding lemmas: Corollary 2.3 Let $`X`$ be a non-singular complex surface $`X`$ then $$c_1(^{i_1}T^{}X\mathrm{}^{i_k}T^{}X)=\frac{i_1+\mathrm{}+i_k}{2}(i_1+1)\mathrm{}(i_k+1),$$ $`c_2(^{i_1}T^{}X\mathrm{}^{i_k}T^{}X)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{k}{}}}{\displaystyle \frac{\underset{l=1}{\overset{k}{}}(i_l+1)}{i_j+1}}c_2(^{i_j}T^{}X)`$ $`+{\displaystyle \underset{j=1}{\overset{k}{}}}{\displaystyle \frac{\underset{lj}{}(i_l+1)\{(\underset{lj}{}(i_l+1))1\}}{2}}c_1^2(^{i_j}T^{}X)`$ $`+{\displaystyle \underset{1j_1<j_2k}{}}{\displaystyle \frac{\underset{l}{}(i_l+1)^2}{(i_{j_1}+1)(i_{j_2}+1)}}c_1(^{i_{j_1}}T^{}X)c_1(^{i_{j_2}}T^{}X).`$ Let $`m`$ be a positive integer and for each fixed positive integer $`k`$ denote by $`q_k(m)`$ be the number of solutions of the equation $$i_1+2i_2+\mathrm{}+ki_k=m.$$ A solution of the preceding equation shall be referred to as a weighted partition of $`m`$ of length $`k`$. It is easy to see that Lemma 2.4 With the notations above we have $`q_k(m)=p_1(m)+\mathrm{}+p_k(m)`$. With this we have the following asymptotic estimate: Theorem 2.5 For $`m\mathrm{}`$ the number $`q_k(m)`$ is asymptotically given by: $$q_k(m)\frac{m^{k1}i^2)m^{j1}}{(k1)!(k1)!}.$$ Proof. By Theorem 1.20, we have $$p_j(m)\frac{m^{j1}}{(j1)!j!}.$$ By Lemma 2.4, $$q_k(m)=\underset{j=1}{\overset{k}{}}p_j(m)\underset{j=1}{\overset{k}{}}\frac{m^{j1}}{(j1)!j!}\frac{m^{k1}i^2)m^{j1}}{(k1)!(k1)!}.$$ QED With the preceding results the computation of the Chern numbers for $`𝒥_k^mX`$ can now be carried out by using the Theorem of Green on Griffiths. First we compute the Chern classes for the sheaves of each of the weighted partitions. Then the Chern numbers of $`𝒥_k^mX`$ is computed from these by the following Lemma. To state the Lemma we denote by $$_{km}=\{I=(i_1,\mathrm{},i_k)\|i_j𝐍,i_1+2i_2+\mathrm{}+ki_k=m\}.$$ Moreover fixing an ordering of the set $`_{km}`$ then Lemma 2.6 Let $`X`$ be a non-singular surface then $$c_1(𝒥_k^mX)=\underset{I_{km}}{}c_1(𝒮_I),$$ $$c_2(𝒥_k^mX)=\underset{I}{}c_2(𝒮_I)+\underset{I<J,I,J_{km}}{}c_1(𝒮_I)c_1(𝒮_J)$$ where $`𝒮_I=^{i_1}T^{}X\mathrm{}^{i_k}T^{}X`$. Proof. This is a consequence of Theorem 1.13: $$𝒥_{k1}^mX=_k^0_k^1\mathrm{}_k^{[m/k]}=𝒥_k^mX$$ $`(`$where $`[m/k]`$ is the greatest integer smaller than or equal to $`m/k)`$ such that $$_k^i/_k^{i1}𝒥_{k1}^{mki}X(^iT^{}X).$$ From the exact sequence $$0_k^{[m/k]1}𝒥_k^mX𝒥_{k1}^{mk[m/k]}X(^{[m/k]}T^{}X)0$$ we see that $$c_1(𝒥_k^mX)=c_1(_k^{[m/k]1})+c_1(𝒥_{k1}^{mk[m/k]}X(^{[m/k]}T^{}X)),$$ $`c_2(𝒥_k^mX)`$ $`=`$ $`c_1(_k^{[m/k]1})c_1(𝒥_{k1}^{mk[m/k]}X(^{[m/k]}T^{}X))`$ $`+c_2(_k^{[m/k]1})+c_2(𝒥_{k1}^{mk[m/k]}X(^{[m/k]}T^{}X)).`$ We then use filtrations of $`_k^{[m/k]1}`$ and of $`𝒥_{k1}^{mk[m/k]}X`$ to compute the Chern classes. Eventually the Chern classes are expressed by the Chern classos of the bundles $`𝒮_I=^{i_1}T^{}X\mathrm{}^{i_k}T^{}X`$ for each $`I_{km}`$. QED We shall compute the explicit numbers for the following cases (I) $`k=2,1m6`$, (II) $`k=3,m=6`$ which will be needed later. We shall aso compute (III) $`k=m5`$ for comparison with the result of section 1. We shall write, for simplicity: $$c_1=c_1(T^{}X),c_2=c_2(T^{}X).$$ $`(I_{22})k=2,m=2`$ There are two weighted partitions $`P_1=(i_1=2,i_2=0)`$ and $`P_2=(i_1=0,i_2=1)`$ corresponding to the two solutions of $`i_1+2i_2=2`$. The corresponding sheaves are $`𝒮_1=^2T^{}X,𝒮_2=T^{}X`$. Denote by $`\mathrm{\Delta }(𝒮_i)=c_1(𝒮_i)c_2(𝒮_i)`$ and $`\mu (𝒮_i)=c_1(\mu (𝒮_i))/\mathrm{rank}\mu (𝒮_i)`$. Thus $`c_1(𝒥_2^2X)=4c_1(T^{}X),c_2(𝒥_2^2X)=5c_1^2(T^{}X)+5c_2(T^{}X)`$, hence $$\mathrm{\Delta }(𝒥_2^2X)=c_1^2(𝒥_2^2X)c_2(𝒥_2^2X)=11c_1^2(T^{}X)5c_2(T^{}X),\mu (𝒥_2^2X)=4/5.$$ We remark that the formula given in \[G-G\] is $`c_1^2(𝒥_2^2X)c_2(𝒥_2^2X)=7c_1^2(T^{}X)5c_2(T^{}X).`$ $`(I_{23})k=2,m=3`$ There are two weighted partitions $`P_1=(i_1=3,i_2=0)`$ and $`P_2=(i_1=1,i_2=1)`$ corresponding to the two solutions of $`i_1+2i_2=3`$. Thus $`c_1(𝒥_2^3X)=10c_1(T^{}X),c_2(𝒥_2^3X)=41c_1^2(T^{}X)+14c_2(T^{}X)`$, hence $$\mathrm{\Delta }(𝒥_2^3X)=59c_1^2(T^{}X)14c_2(T^{}X),\mu (𝒥_2^3X)=5/4.$$ $`(I_{24})k=2,m=4`$ There are 3 weighted partitions $`P_1=(i_1=4,i_2=0),P_1=(i_1=2,i_2=1)`$ and $`P_3=(i_1=0,i_2=2)`$ corresponding to the 3 solutions of $`i_1+2i_2=4`$. Thus $`c_1(𝒥_2^4X)=22c_1(T^{}X),c_2(𝒥_2^4X)=203c_1^2(T^{}X)+35c_2(T^{}X)`$, hence $$\mathrm{\Delta }(𝒥_2^4X)=281c_1^2(T^{}X)35c_2(T^{}X),\mu (𝒥_2^4X)=11/7.$$ $`(I_{25})k=2,m=5`$ There are 3 weighted partitions $`P_1=(i_1=5,i_2=0),P_1=(i_1=3,i_2=1)`$ and $`P_3=(i_1=1,i_2=2)`$ corresponding to the 3 solutions of $`i_1+2i_2=5`$. Thus $`c_1(𝒥_2^5X)=40c_1(T^{}X),c_2(𝒥_2^5X)=750c_1^2(T^{}X)+70c_2(T^{}X)`$, hence $$\mathrm{\Delta }(𝒥_2^5X)=c_1^2(𝒥_2^5X)c_2(𝒥_2^5X)=850c_1^2(T^{}X)70c_2(T^{}X),\mu (𝒥_2^5X)=2.$$ $`(I_{26})k=2,m=6`$ There are 4 weighted partitions $`P_1=(i_1=6,i_2=0),P_2=(i_1=4,i_2=1),P_3=(i_1=2,i_2=1)`$ and $`P_4=(i_1=0,i_2=3)`$ corresponding to the 3 solutions of $`i_1+2i_2=6`$. Thus $`c_1(𝒥_2^6X)=70c_1(T^{}X),c_2(𝒥_2^6X)=662c_1^2(T^{}X)+135c_2(T^{}X)`$, hence $$\mathrm{\Delta }(𝒥_2^6X)=4238c_1^2(T^{}X)135c_2(T^{}X),\mu (𝒥_2^6X)=7/3.$$ $`(II_{36})k=3,m=6`$ There are 7 weighted partitions $`P_1=(i_1=6,i_2=0,i_3=0),P_2=(i_1=4,i_2=1,i_3=0),P_3=(i_1=3,i_2=0,i_3=1),P_4=(i_1=2,i_2=2,i_3=0),P_5=(i_1=1,i_2=1,i_3=1),P_6=(i_1=0,i_2=3,i_3=0)`$ and $`P_7=(i_1=0,i_2=0,i_3=2)`$ corresponding to the 7 solutions of $`i_1+2i_2+3i_3=6`$. us $`c_1(𝒥_3^6X)=101c_1(T^{}X),c_2(𝒥_3^6X)=5026c_1^2(T^{}X)+175c_2(T^{}X)`$ and $$\mathrm{\Delta }(𝒥_3^6X)=5175c_1^2(T^{}X)175c_2(T^{}X),\mu (𝒥_3^6X)=101/49.$$ $`(III_{33})k=m=3`$ In this case there are 3 weighted partitions: $`P_1=(3,0,0),P_2=(1,1,0)`$ and $`P_3=(0,0,3)`$. The tabulation is given by Thus we have $$c_1(𝒥_3^3X)=11c_1^2(T^{}X),c_2(𝒥_3^3X)=51c_1^2(T^{}X)+15c_2(T^{}X)$$ and so $`\mu (𝒥_3^3X)=11/10`$ and $$c_1^2(𝒥_3^3X)c_2(𝒥_3^3X)=70c_1^2(T^{}X)15c_2(T^{}X).$$ The formula given in \[G-G\] is $`c_1^2(𝒥_3^3X)c_2(𝒥_3^3X)=85c_1^2(T^{}X)49c_2(T^{}X).`$ $`(III_{44})k=m=4`$ In this case there are 5 weighted partitions: $`P_1=(4,0,0,0),P_2=(2,1,0,0),P_3=(1,0,1,0),P_4=(0,2,0,0)`$ and $`P_5=(0,0,0,1)`$. The tabulation is given by Thus we have $$c_1(𝒥_4^4X)=27c_1^2(T^{}X),c_2(𝒥_4^4X)=338c_1^2(T^{}X)+40c_2(T^{}X)$$ and so $`\mu (𝒥_4^4X)=27/20`$ and $$c_1^2(𝒥_4^4X)c_2(𝒥_4^4X)=391c_1^2(T^{}X)40c_2(T^{}X).$$ $`(III_{55})k=m=5`$ In this case there are 3 weighted partitions: $`P_1=(5,0,0,0,9),P_2=(3,1,0,0,0),P_3=(2,0,1,0,0),P_4=(1,2,0,0,0),P_5=(1,0,0,1,0)P_6=(1,2,0,0,0)`$ and $`P_7=(1,2,0,0,0)`$. The tabulation is given by Thus we have $$c_1(𝒥_5^5X)=58c_1^2(T^{}X),c_2(𝒥_5^5X)=1622c_1^2(T^{}X)+90c_2(T^{}X)$$ and so $`\mu (𝒥_5^5X)=29/18`$ and $$c_1^2(𝒥_5^5X)c_2(𝒥_4^4X)=1742c_1^2(T^{}X)90c_2(T^{}X).$$ We remark that the inequality (1.21) in \[G-G\] is incorrect (for example set $`k=2`$ or $`k=3`$ and compare these to the formulas obtained above; indeed for $`k=2`$ (see (1.21) in \[G-G\]) reduces to $`c_1^2(X)c_2(X)>0`$). § 3 Weighted Projective Spaces and Projectivized Jet Bundles For a vector bundle, e.g., the $`k`$-jet bundle $`T^kX`$, a standard approach of studying the bundle is to projectivized it and then study the line bundles over the projectivization. We are going to do the same for the $`𝐂^{}`$-bundle $`J^kX`$ using the well-known results in the former case as a guide. The fiber of the projectivized bundle are certain types of weighted projective space. Thus we shall first recall some basic facts about weighted projective spaces. For more detailed discussions and further references the readers are referred to the articles \[B-R\], \[Do\] and the monograph \[Di\]. Let $`Q=(q_0,q_1,\mathrm{},q_r)`$ ($`r1`$) be an $`(r+1)`$-tuple of positive integers. The tuple $`Q`$ is said to be reduced if the greatest common divisor (gcd) of $`(q_0,q_1,\mathrm{}.q_r)`$ is 1. In general if the gcd is $`d`$ the tuple $$Q_{\mathrm{red}}=Q/d=(q_0/d,\mathrm{},q_r/d)$$ is called the reduction of $`Q`$. Let $`d_0=gcd(q_1,\mathrm{},q_r),d_r=gcd(q_0,\mathrm{},q_{r1})`$ and $$d_i=\mathrm{gcd}(q_0,q_1,\mathrm{},q_{i1},q_{i+1},\mathrm{},q_r),1ir1.$$ Let $`a_0=lcm(d_1,\mathrm{},d_r),a_r=lcm(d_0,\mathrm{},d_{r1})`$ and $$a_i=\mathrm{lcm}(d_0,d_1,\mathrm{},d_{i1},d_{i+1},\mathrm{},d_r),1ir1$$ where $`lcm`$ means least common multiple. Define the normalization of $`Q`$ by $$Q_{\mathrm{norm}}=(q_0/a_0,\mathrm{},q_r/a_r).$$ A tuple $`Q`$ is said to be normalized if $`Q=Q_{\mathrm{norm}}`$. Let $`(𝐂^{r+1},Q)`$ be the $`(r+1)`$-dimensional complex vector space such that the variable $`z_i`$ is assigned the weight (or degree) $`q_i`$. A $`𝐂^{}`$-action is defined on $`(𝐂^{r+1},Q)`$ by: $`\lambda .(z_0,\mathrm{},z_r)=(\lambda ^{q_0}z_0,\mathrm{},\lambda ^{q_r}z_r),\lambda 𝐂^{}.`$ (41) The quotient space, $`𝐏(Q)=(𝐂^{r+1},Q)/𝐂^{}`$, is called the weighted projective space of type $`Q`$. The equivalence class of an element $`(z_0,\mathrm{},z_r)`$ is denoted by $`[z_0,\mathrm{},z_r]_Q`$. For $`Q=(1,\mathrm{},1)=\mathrm{𝟏},𝐏(Q)=𝐏^r`$ is the usual complex projective space of dimension $`r`$ and an element of $`𝐏^r`$ is denoted simply by $`[z_0,\mathrm{},z_r]`$. Indeed for the special case $`r=1`$ it can be shown that, for any tuple $`(q_0,q_1)`$, $`𝐏(q_0,q_1)𝐏^1`$. This is not so if $`r2`$, however, we do have: Theorem 3.1 Let $`Q=(q_0,\mathrm{},q_r)`$ be an $`(r+1)`$-tuple of positive integers then $$𝐏(Q)𝐏(Q_{\mathrm{red}})𝐏(Q_{\mathrm{norm}}).$$ Example 3.2 It is clear that a normalized tuple is reduced. The converse is not true in general. Let $`Q=(4,6,12)`$ then $`Q_{\mathrm{red}}=(2,3,6)`$ is reduced but is not normalized. In fact $`Q_{\mathrm{norm}}=(Q_{\mathrm{red}})_{\mathrm{norm}}=(1,1,6)`$. The tuple $`(6,10,15)`$ is reduced but is not normalized, in fact its normalization is $`(1,1,1)`$ hence $`𝐏(6,10,15)𝐏^2`$. Define a map $`\rho _Q:(𝐂^{r+1},\mathrm{𝟏})(𝐂^{r+1},Q)`$ by $`\rho _Q(z_0,\mathrm{},z_r)=(z_0^{q_0},\mathrm{},z_r^{q_r}).`$ (42) It is easily seen that $`\mu _Q`$ is compatible with the respective $`𝐂^{}`$-actions and hence descends to a well-defined morphism: $`\overline{\rho }_Q:𝐏^r𝐏(Q),\overline{\rho }_Q([z_0,\mathrm{},z_r])=[z_0^{q_0},\mathrm{},z_r^{q_r}]_Q.`$ (43) The weighted projective space can aso be described as follows. Denote by $`\mathrm{\Theta }_{q_i}`$ the group of $`q_i`$-th roots of unity. Then the group $`\mathrm{\Theta }_Q=_{i=0}^r\mathrm{\Theta }_{q_i}`$ acts on $`𝐏^r`$ by coordinate wise multiplication: $$(\theta _0,\mathrm{},\theta _r).[z_0,\mathrm{},z_r]=[\theta _0z_0,\mathrm{},\theta _rz_r],\theta _i\mathrm{\Theta }_{q_i}$$ and it is easily verified that $`𝐏(Q)=𝐏^r/\mathrm{\Theta }_Q.`$ Theorem 3.3 The weighted projective space $`𝐏(Q)`$ is isomorphic to the quotient $`𝐏^r/\mathrm{\Theta }_Q`$. In particular, $`𝐏(Q)`$ is irreducible and normal (the singularities are cyclic quotients and hence rational). Denote by $`S_Q(m)`$ the space of homogeneous polynomials of degree $`m>0`$ in the variables $`z_i`$ (assigned with the degree $`q_i`$). In other words, a polynomial $`P`$ is in $`S(Q)(m)`$ if $$P(\lambda .(z_0,\mathrm{},z_r))=\lambda ^mP(z_0,\mathrm{},z_r).$$ We may express such a polynomial explicitly: $`P={\displaystyle \underset{(i_0,\mathrm{},i_r)_{Q,m}}{}}a_{i_0\mathrm{}i_r}z_0^{i_0}\mathrm{}z_r^{i_r}`$ (44) where the index set $`_{Q,m}`$ is defined by: $$_{Q,m}=\{(i_0,\mathrm{},i_r)\|\underset{j=0}{\overset{r}{}}q_ji_j=m\}.$$ The sheaf $`𝒪_{𝐏(Q)}(m)`$ is the sheaf over $`𝐏(Q)`$ whose global regular sections are precisely the elements of $`S_Q(m)`$: $`H^0(𝐏(Q),𝒪_{𝐏(Q)}(m))=S_Q(m).`$ (45) For negative integer $`m,m>0`$ the sheaf $`𝒪_{𝐏(Q)}(m)`$ is defined to be the dual of $`𝒪_{𝐏(Q)}(m).`$ Theorem 3.4 (i) For any $`m𝐙,𝒪_{𝐏(Q)}(m)`$ is a reflexive coherent sheaf. (ii) The sheaf $`𝒪_{𝐏(Q)}(m)`$ is locally free if $`m`$ is divisible by each $`q_i(`$hence by the least common multiple$`)`$. (iii) Let $`m_Q`$ be the least common multiple of $`\{q_0,\mathrm{},q_r\}`$ then $`𝒪_{𝐏(Q)}(m_0)`$ is ample. (iv) There exists an interger $`n_0`$ depending only on $`Q`$ such that $`𝒪_{𝐏(Q)}(nm_Q)`$ is very ample for all $`nn_0`$. (v) For any $`\alpha ,\beta 𝐙`$ we have $`𝒪_{𝐏(Q)}(\alpha m_Q)𝒪_{𝐏(Q)}(\beta )𝒪_{𝐏(Q)}(\alpha m_Q+\beta )`$. For any subset $`J\{0,1,\mathrm{},r\}`$ denote by $`m_J`$ the least common multiple of $`\{q_j,jJ\}`$ and define $$m(Q)=\|Q\|+\frac{1}{r}\underset{\nu =2}{\overset{r+1}{}}\frac{\underset{\mathrm{\#}J=\nu }{}m_J}{C_{\nu 2}^{r1}}$$ where $`C_a^b`$ is the usual binomial coefficient and $`\|Q\|=q_0+\mathrm{}+q_r`$. It is known that assertion (iv) holds if $`n>m(Q)`$. In general the line sheaf $`𝒪_{𝐏(Q)}(m)`$ is not invertible if $`m`$ is not an integer multiple of $`m_Q`$. It can be shown that for $`Q=(1,1,2)`$ the sheaf $`𝒪_{𝐏(Q)}(1)`$ is not invertible and hence, neither is $`𝒪_{𝐏(Q)}(1)𝒪_{𝐏(Q)}(1)`$. This also shows that $`𝒪_{𝐏(Q)}(1)𝒪_{𝐏(Q)}(1)\cong ̸𝒪_{𝐏(Q)}(2)`$ as $`𝒪_{𝐏(Q)}(2)`$ is invertible by part (ii) of the preceding Theorem. Theorem 3.5 Let $`Q`$ be a $`(r+1)`$-tuple of positive integers then $`(i)H^i(𝐏(Q),𝒪_𝐏(Q)(p))=\{0\},p𝐙\mathrm{if}i0,r;`$ $`(ii)H^0(𝐏(Q),𝒪_𝐏(Q)(p))=S_Q(p)p𝐙;`$ $`(iii)H^r(𝐏(Q),𝒪_𝐏(Q)(p))S(Q)(p\|Q\|),p𝐙`$ where $`\|Q\|=q_0+\mathrm{}.+q_r`$. Denote by $`Pic(𝐏(Q))`$ and $`Cl(𝐏(Q))`$ the Picard group and respectively the divisor class group. Theorem 3.6 Let $`Q=Q_{\mathrm{norm}}`$ be a normalized $`(r+1)`$-tuple of positive integers then $`(i)Pic(𝐏(Q))𝐙`$ is generated by $`[𝒪_𝐏(Q)(m_Q)];(ii)Cl(𝐏(Q))𝐙`$ is generated by $`[𝒪_𝐏(Q)(1)].`$ Let $`Q`$ be a $`(r+1)`$-tuple of positive integers define for $`k=1,\mathrm{},r`$: $`l_{Q,k}=lcm\{{\displaystyle \frac{q_{i_0}\mathrm{}q_{i_k}}{gcd(q_0,\mathrm{},q_{i_k})}}\|0i_0<\mathrm{}.<i_kr\}.`$ Theorem 3.7 Let $`Q`$ be a $`(r+1)`$-tuple of positive integers then $`H^i(𝐏(Q);𝐙)\{\begin{array}{cc}𝐙,\hfill & \text{if }i\text{ is even,}\hfill \\ 0,\hfill & \text{if }i\text{ is odd.}\hfill \end{array}`$ Moreover, let $`\overline{\rho }_Q:𝐏^r𝐏(Q)`$ be the quotient map as defined by $`(28)`$ then the following diagram commutes, $`H^{2k}(𝐏(Q);𝐙)\stackrel{\overline{\rho }_Q^{}}{}H^{2k}(𝐏^r;𝐙)`$ $``$ $`𝐙\stackrel{l_{Qk}}{}𝐙`$ where the lower map is the multiplication by the number $`l_{Qk}`$. Note that the number $`l_{Qr}`$ is precisely the number of preimages of a point in $`𝐏(Q)`$ under the quotient map $`\overline{\rho }_Q`$. The proof of the preceding Theorem for $`k=r`$ is quite easy. for the general case we refer the readers to \[Ka\]. We shall only be concerned with the case where $`n,k1`$ are positive integers and $`Q=((\underset{n}{\underset{}{1,\mathrm{},1}}),(\underset{n}{\underset{}{2,\mathrm{},2}}),\mathrm{},(\underset{n}{\underset{}{k,\mathrm{},k}})).`$ In this case we shall write $`𝐏_{n,k}`$ for $`𝐏(Q)`$. Note that $`r=`$ dim $`𝐏_{n,k}=nk1`$ In this case the least common multiple of $`Q`$ is $`m_Q=k!`$ and $`l_{Qr}=(k!)^n`$. Let $`\pi :(,h)X`$ be a holomorphic hermitian vector bundle over a compact Kähler manifold $`X`$. Denote by $`()`$ be the ”hyperplane bundle” defined over the projectivized bundle $`𝐏()`$. It is defined as follows: $`\pi ^{}`$ $`prp`$ $`𝐏()\stackrel{\pi }{}X`$ the tautological sub-sheaf is defined by: $$\{((x,[\xi ]),\eta )\pi ^{}\|(x,[\xi ])𝐏(),p([\xi ])=x,[\eta ]=[\xi ]\}$$ and $`_k`$ is defined to be the dual of the tautological line bundle. In other words, since the fiber $`𝐏()`$ over a point $`xX`$ is a projective space, the restriction of $`_kX`$ to $`𝐏(_x)`$ is the hyperplane line bundle $`𝒪_{𝐏^{r1}}(1)`$ (here $`r=`$ rank $``$). We shall often use the notation $`𝒪_{𝐏()}(1)`$ for $`_k`$ and the tensor product $`_k^m`$ by $`𝒪_{𝐏()}(m)`$ for any integer $`m𝐙`$. The following is a classical Theorem of Grothendieck: Theorem 3.8 Let $``$ be a holomorphic vector bundle over a complex manifold $`X`$ then for any $`m,j0`$, the j-th direct image sheaf of the $`m`$-fold tensor product of $`(m)`$ is isomorphic to the $`m`$ fold symmetric product of $`E`$, i.e., $`R_{}^j^m()^m`$ and $$H^j(X,^m𝒮)H^q(X,^m()p^{}𝒮)$$ where $`𝒮`$ is any sheaf on $`X`$. Let $`\pi :J^kXX`$ be the (restricted) $`k`$-jet bundle of a complex manifold $`X`$. Denote by $`_k`$ the ”hyperplane sheaf” defined over the projectivized $`k`$-jet bundle $`𝐏(J^kX)`$. It is defined as follows. Consider the commutative diagram: $`\pi ^{}J^kXJ^kX`$ $`prp`$ $`𝐏(J^kX)\stackrel{\pi }{}X`$ the tautological sub-sheaf is defined by: $$\{((x,[\xi ]),\eta )\pi ^{}J^kX\|(x,[\xi ])𝐏(J^kX),p([\xi ])=x,[\eta ]=[\xi ]\}$$ and $`_k`$ is defined to be the dual the tautological line sheaf. In other words, since the fiber $`𝐏(J_x^kX)`$ over a point $`xX`$ is a weighted projective space of type $`Q=((1,\mathrm{},1);\mathrm{};(k,\mathrm{},k))`$ the restriction of $`_kX`$ to $`𝐏(J_x^kX)`$ is the line sheaf $`𝒪_{𝐏(Q)}(1)`$ as defined in the preceding section. We shall use the notation $`𝒪_{𝐏(J^kX)}(1)`$ for $`_k`$. More generally for any integer $`m,𝒪_{𝐏(J^kX)}(m)`$ is the sheaf on $`𝐏(J^kX)`$ which restricts to the bundle $`𝒪_{𝐏(Q)}(m)`$ along each fiber of the projection map $`p:𝐏(J^kX)X`$. The proof of the preceding Theorem relies on the classical Vanishing Theorem of cohomologies on projective spaces. The analogoue of this for weighted projective spaces is provided by Theorem 3.3 and hence we have (see \[G-G\] and \[K-O\]): Theorem 3.9 Let $`X`$ be a complex manifold and $`𝒮`$ be a sheaf over $`X`$ then for any $`m,j0`$ we have $`R_{}^j𝒪_{𝐏(J^kX)}(m)𝒥_k^mX`$ and $$H^j(X,𝒥_k^mX𝒮)H^j(𝐏(J^kX),𝒪_{𝐏(J^kX)}(m)p^{}𝒮).$$ The following is a consequence of Theorem 3.4: Theorem 3.10 Let $`X`$ be a complex manifold then (i) for any $`m𝐙,𝒪_{𝐏(J^kX)}(m)`$ is a reflexive coherent sheaf; (ii) the sheaf $`𝒪_{𝐏(J^kX)}(m)`$ is locally free if $`m`$ is divisible by each $`q_i(`$hence by the least common multiple $`k!);`$ (iii) for any $`\alpha ,\beta 𝐙,𝒪_{𝐏(J^kX)}(k!\alpha )𝒪_{𝐏(J^kX)}(\beta )𝒪_{𝐏(J^kX)}(k!\alpha +\beta )`$. Due to the fact that $`𝒪_{𝐏(J^kX)}(1)`$ is not locally free and that, in general, $`𝒪_{𝐏(J^kX)}(a)𝒪_{𝐏(J^kX)}(a)𝒪_{𝐏(J^kX)}(a+b)`$ some of the proof of the results that are valid on projectivized vector bundle are not valid even though modifications of the results can be obtained via alternative proofs. We establish some of the results (the counterparts in the case of projectivized vector bundle are well-known) that will be essential in the next section. Lemma 3.11 Let $`X`$ be a complex manifold of dimension $`n`$ and let $`p:𝐏(J^kX)X`$ be the projection map. Then the natural morphism: $$\varphi :p^{}p_{}𝒪_{𝐏(J^kX)}(k!)𝒪_{𝐏(J^kX)}(k!)$$ is surjective and $$\underset{i=0}{\overset{n}{}}(1)^ic_1^{ri}(𝒪_{𝐏(J^kX)}(k!)).p^{}c_i(𝒥_k^{k!}X)=0$$ where $``$ is the kernel of $`\varphi `$ and $`r=nk1`$ is the fiber dimension of $`p`$. Proof. For simplicity we write $`𝒪(k!)`$ for $`𝒪_{𝐏(J^kX)}(k!)`$. By definition the restriction of $`𝒪(k!)`$ to a fiber of the projection map is $`𝒪_{𝐏(Q)}(k!)`$ where $$Q=((1,\mathrm{},1);\mathrm{};(k,\mathrm{},k)).$$ Thus the least common multiple of the indices is $`k!`$ and $`𝒪_{𝐏(Q)}(k!)`$ is ample by Theorem 3.4. This implies that the map $`\varphi `$ is surjective (see for example \[B-S\]). By Theorem 3.9 $`p^{}p_{}𝒪(k!)=p^{}𝒥_k^{k!}X`$ and so the sequence: $$0p^{}𝒥_k^{k!}X𝒪(k!)0$$ is eact. By Whitney’s formula $$\underset{i=0}{\overset{r}{}}p^{}c_i(𝒥_k^{k!}X)=(1+c_1(𝒪(k!)).\underset{i=0}{\overset{r1}{}}c_i()$$ and hence $$p^{}c_i(𝒥_k^{k!}X)=c_1(𝒪(k!)).c_{i1}()+c_i()$$ for $`0ir`$ with $`c_1()=c_r()=0`$ (as rank $`=r1`$). We can eliminate the Chern classes of $``$ by first multiplying the preceding identity by $`c_1^{ri}(𝒪(k!))`$ and then take alternating sum; namely: $`c_1^{r1}(𝒪(k!)).p^{}c_1(𝒥_k^{k!}X)`$ $`=`$ $`c_1^r(𝒪(k!))+c_1^{r1}(𝒪(k!)).c_1()`$ $`c_1^{r2}(𝒪(k!)).p^{}c_2(𝒥_k^{k!}X)`$ $`=`$ $`c_1^{r1}(𝒪(k!)).c_1()+c_1^{r2}(𝒪(k!)).c_2()`$ $`c_1^{r3}(𝒪(k!)).p^{}c_3(𝒥_k^{k!}X)`$ $`=`$ $`c_1^{r1}(𝒪(k!)).c_2()+c_1^{r3}(𝒪(k!)).c_3()`$ $`\mathrm{}`$ $`c_1(𝒪(k!)).p^{}c_{r1}(𝒥_k^{k!}X)`$ $`=`$ $`c_1^2(𝒪(k!)).c_{r2}()+c_1(𝒪(k!)).c_{r1}()`$ $`p^{}c_r(𝒥_k^{k!}X)`$ $`=`$ $`c_1(𝒪(k!)).c_{r1}()`$ and multiply the $`i`$-identity above by $`(1)^r`$ and then taking the sum from $`i=1`$ to $`i=r`$ yields $$\underset{i=1}{\overset{r}{}}(1)^ic_1^{ri}(𝒪(k!)).p^{}c_i(𝒥_k^{k!}X)=c_1^r(𝒪(k!)).$$ Moving the RHS to the LHS yields the identity of the Lemma. QED Note that $`c_i(𝒥_k^{k!}X)=0`$ if $`in=`$ dim $`X`$. Lemma 3.12 Let $`X`$ be a compact complex manifold of complex dimension $`n`$ then for any $`xX`$, $$_{𝐏(J^kX)_x}c_1^{nk1}(𝒪_{𝐏(J^kX)}(k!)\|_{𝐏(J^kX)_x})=(k!)^n$$ where $`𝐏(J^kX)_x`$ is the fiber over $`x`$. Proof. By definition the fiber, $`𝐏(J^kX)_x`$, over any point $`xX`$ of the projection map of $`p:𝐏(J^kX)X`$ is the weighted projective space $$𝐏((\underset{n}{\underset{}{1,\mathrm{},1}}),(\underset{n}{\underset{}{2,\mathrm{},2}}),\mathrm{},(\underset{n}{\underset{}{k,\mathrm{},k}}))$$ of dimension $`nk1`$. By Theorem 3.7 the quotient map $`\overline{\rho }_Q:𝐏^{nk1}𝐏(Q)`$ is a finite morphism with sheet number $`l_{Q,nk1}=(k!)^n`$. The generator of $`H^{2(nk1)}(𝐏^{nk1};𝐙)`$ is represented by the $`(nk1)`$-th power, $`\omega _{FS}^{nk1}`$, of the the Fubini-Study metric $`\omega _{FS}=c_1(𝒪_{𝐏^{nk1}}(1)`$. The Lemma follows readily as we have: $$_{𝐏^{nk1}}c_1^{nk1}(𝒪_{𝐏^{nk1}}(1))=_{𝐏^{nk1}}\omega _{FS}^{nk1}=1.$$ QED Theorem 3.13 Let $`X`$ be a compact complex manifold of complex dimension $`n`$ then the following intersection formulas hold$`:`$ $$c_1^{nk+j1}(𝒪_{𝐏(J^kX)}(k!)).p^{}D_1.\mathrm{}.p^{}D_{nj}=(k!)^n\mathrm{\Delta }_j.D_1.\mathrm{}.D_{nj}$$ for divisors $`D_1,\mathrm{},D_{nj},j=0,1,\mathrm{},n`$ on $`X`$. The numbers $`\mathrm{\Delta }_j`$ is defined by setting $`\mathrm{\Delta }_0=1,\mathrm{\Delta }_1=c_1(𝒥_k^{k!}X)`$ and by the recursive relation$`:`$ $$\mathrm{\Delta }_j=\underset{i=1}{\overset{j}{}}(1)^{i+1}\mathrm{\Delta }_{ji}.c_i(𝒥_k^{k!}X),j2.$$ Proof. Note that dim $`𝐏(J^kX)=n(k+1)1`$ and the fiber dimension, dim $`𝐏(J^kX)_x=nk1`$. Thus, by fiber integration (Lemma 3.12), $$c_1^{nk1}(𝒪_{𝐏(J^kX)}(k!)).p^{}D_1.\mathrm{}.p^{}D_n=(k!)^n\mathrm{\Delta }_0D_1.\mathrm{}.D_n$$ which is the case $`j=0`$. By Lemma 3.11 with $`r=nk1`$, $`{\displaystyle \underset{i=0}{\overset{r}{}}}(1)^ic_1^{ri}(𝒪_{𝐏(J^kX)}(k!)).p^{}c_i(𝒥_k^{k!}X)=0`$ (47) and, multiplying by $`p^{}D_1.\mathrm{}.p^{}D_{n1}`$, we get $`c_1^{nk1}(𝒪_{𝐏(J^kX)}(k!)).p^{}D_1.\mathrm{}.p^{}D_{n1}`$ $`=c_1^{nk2}(𝒪_{𝐏(J^kX)}(k!)).p^{}c_1(𝒥_k^{k!}X).p^{}D_1.\mathrm{}.p^{}D_{n1}`$ as the rest of the terms vanish for dimension reason. Multiplying the above by $`c_1(𝒪_{𝐏(J^kX)}(k!))`$ yields, $`c_1^{nk}(𝒪_{𝐏(J^kX)}(k!)).p^{}D_1.\mathrm{}.p^{}D_{n1}`$ $`=c_1^{nk1}(𝒪_{𝐏(J^kX)}(k!)).p^{}c_1(𝒥_k^{k!}X).p^{}D_1.\mathrm{}.p^{}D_{n1}.`$ Fiber integration shows that the term on the right above equals $$(k!)^nc_1(𝒥_k^{k!}X).D_1.\mathrm{}.D_{n1}=(k!)^n\mathrm{\Delta }_1.D_1.\mathrm{}.D_{n1}.$$ This establish the Theorem for the case $`j=1`$. If we multiply (31) by $`p^{}D_1.\mathrm{}.p^{}D_{n2}`$ we are left with 3 terms (again for dimension reason): $`c_1^{nk1}(𝒪_{𝐏(J^kX)}(k!)).p^{}D_1.\mathrm{}.p^{}D_{n2}`$ $`=c_1^{nk2}(𝒪_{𝐏(J^kX)}(k!)).p^{}c_1(𝒥_k^{k!}X).p^{}D_1.\mathrm{}.p^{}D_{n2}`$ $`c_1^{nk3}(𝒪_{𝐏(J^kX)}(k!)).p^{}c_2(𝒥_k^{k!}X).p^{}D_1.\mathrm{}.p^{}D_{n2}.`$ Now multiply the above by $`c_1^2(𝒪_{𝐏(J^kX)}(k!))`$ then the LHS above is given by $$c_1^{nk+1}(𝒪_{𝐏(J^kX)}(k!)).p^{}D_1.\mathrm{}.p^{}D_{n2}$$ while the first term on the RHS is given by (the case $`j=1`$): $`c_1^{nk}(𝒪_{𝐏(J^kX)}(k!)).p^{}c_1(𝒥_k^{k!}X).p^{}D_1.\mathrm{}.p^{}D_{n2}`$ $`=`$ $`(k!)^n\mathrm{\Delta }_1.c_1(𝒥_k^{k!}X).D_1.\mathrm{}.D_{n2}`$ and the second term on the right is given by (the case $`j=0`$): $`c_1^{nk1}(𝒪_{𝐏(J^kX)}(k!)).p^{}c_2(𝒥_k^{k!}X).p^{}D_1.\mathrm{}.p^{}D_{n2}`$ $`=`$ $`(k!)^n\mathrm{\Delta }_0c_2(𝒥_k^{k!}X).D_1.\mathrm{}.D_{n2}`$ Combining the above yields the case $`j=2`$: $`c_1^{nk+1}(𝒪_{𝐏(J^kX)}(k!)).p^{}D_1.\mathrm{}.p^{}D_{n2}`$ $`=`$ $`(k!)^n\{\mathrm{\Delta }_1.c_1(𝒥_k^{k!}X)\mathrm{\Delta }_0c_2(𝒥_k^{k!}X)\}.D_1.\mathrm{}.D_{n2}`$ $`=`$ $`(k!)^n\mathrm{\Delta }_2.D_1.\mathrm{}.D_{n2}`$ as, by definition, $$\mathrm{\Delta }_2=\mathrm{\Delta }_1.c_1(𝒥_k^{k!}X)\mathrm{\Delta }_0c_2(𝒥_k^{k!}X)=c_1^2(𝒥_k^{k!}X)c_2(𝒥_k^{k!}X).$$ Thus the case $`j=2`$ is also established. Inductively, the prcedure above yields: $`c_1^{nk+j2}(𝒪_{𝐏(J^kX)}(k!)).p^{}D_1.\mathrm{}.p^{}D_{n3}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{j}{}}}(1)^{i1}\mathrm{\Delta }_{ji}.c_i(𝒥_k^{k!}X).D_1.\mathrm{}.D_{n3}`$ $`=`$ $`(k!)^n\mathrm{\Delta }_j.D_1.\mathrm{}.D_{n3}.`$ QED Theorem 3.14 Let $`X`$ be a non-singular projective surface and assume that $`(i)c_1^2(𝒥_k^{k!}X)c_2(𝒥_k^{k!}X)>0`$ and $`(ii)h^2(𝒥_k^{k!m})=O(m^{(n+1)k2)!})`$. Then $`𝒥_k^{k!}X`$ is big. Proof. Let $`𝐏(J^kX)`$ be the projectivized $`k`$-jet bundle. Then $`dim𝐏(J^kX)=(n+1)k1`$. Riemann-Roch applied to the line bundle $`𝒪_{𝐏(J^kX)}(k!)`$ yields $$\chi (𝒪_{𝐏(J^kX)}(k!m))=\frac{c_1^{(n+1)k1}(𝒪_{𝐏(J^kX)}(k!))}{((n+1)k1)!}m^{(n+1)k1}+O(m^{(n+1)k2}).$$ Theorem 3.13 and assumption (i) imply that there exists positive constant $`c>0`$ and positive integer $`m_0^{^{}}`$ such that $`\chi (𝒪_{𝐏(J^kX)}(k!m))`$ $`=`$ $`{\displaystyle \frac{c_1^2(𝒥_k^{k!}X)c_2(𝒥_k^{k!}X)}{((n+1)k1)!}}m^{(n+1)k1}+O(m^{(n+1)k2})`$ $``$ $`cm^{(n+1)k1}`$ for all $`mm_0^{}`$. Theorem 3.8 implies that the same is true for $`𝒥_k^{k!m}`$ i.e. $`\chi (𝒥_k^{k!m})cm^{r+1}`$ and, a priori: $$h^0(𝒥_k^{k!m})+h^2(𝒥_k^{k!m})>cm^{(n+1)k1}$$ for all $`mm_0^{}`$. The Theorem follows now from assumption (ii). § 4 Surfaces of General Type We recall first some well-known results on manifolds of general type. Theorem 4.1 Let $`X`$ be a minimal surface of general type then $`c_1^2(T^{}X)>0,c_2(T^{}X)>0`$ and $`c_1^2(T^{}X)3c_2(T^{}X)`$. Moreover, we have $$5c_1^2(T^{}X)c_2(T^{}X)+360,\mathrm{if}m\mathrm{is}\mathrm{even},$$ $$5c_1^2(T^{}X)c_2(T^{}X)+300,\mathrm{if}m\mathrm{is}\mathrm{odd}.$$ Let $`L_0`$ be a nef line bundle on a non-singular surface $`X`$. A coherent sheaf $`E`$ over $`X`$ is said to be semi-stable (resp. stable) with respect to $`L_0`$ if $`c_1(E).c_1(L_0)0`$ and if, for any coherent subsheaf $`0𝒮`$ of $`E`$, we have: $`\mu _{𝒮,L_0}\stackrel{\mathrm{def}}{=}{\displaystyle \frac{c_1(𝒮).c_1(L_0)}{\mathrm{rank}𝒮}}\mu _{E,L_0}\stackrel{\mathrm{def}}{=}{\displaystyle \frac{c_1(E).c_1(L_0)}{\mathrm{rank}E}}`$ (48) (resp. $`\mu _{𝒮,L_0}<\mu _{E,L_0}).`$ If $`X`$ is of general type then (see Maruyama \[Ma\]) Theorem 4.2 Let $`X`$ be a surface of general type then $`^mT^{}X,^mT^{}X`$ are semi-stable with respect to the canonical bundle $`𝒦_X=detT^{}X`$. Indeed we have: Theorem 4.3 Let $`X`$ be a minimal surface of general type. If $`D`$ is a divisor in $`X`$ such that $`H^0(X,E_k[D])0`$ where $`E_k=(^{i_1}T^{}X\mathrm{}^{i_k}T^{}X),i_1,\mathrm{},i_k`$ being positive integers then $$c_1(E_k).c_1(D)\mu _{E_k}\frac{m}{2}c_1^2(T^{}X)$$ with $`m=i_1+2i_2+\mathrm{}+ki_k`$ Proof. This follows from the calculation of the Chern number $`c_1(E_k)`$ in section 2. The computation there shows that $$\mu _{E_k}\frac{m}{2}c_1^2(T^{}X)$$ with equality if and only if $`k=1`$, i.e., $$\mu _{^mT^{}X}=\frac{c_1(^mT^{}X)}{\mathrm{rank}^mT^{}X}.c_1(T^{}X)=\frac{\frac{m(m+1)}{2}}{m+1}c_1^2(T^{}X)=\frac{m}{2}c_1^2(T^{}X).$$ QED Note that in general, if $`E`$ is a vector bundle of rank $`r`$ then $`\mathrm{rank}^mE={\displaystyle \frac{(m+r1)!}{(r1)!m!}}.`$ (49) The Chern number $`c_1(^mE)`$ is given by (compare section 2) $`c_1(^mE)={\displaystyle \frac{1}{r!}}{\displaystyle \frac{(m+r1)!}{(m1)!}}c_1(E).`$ (50) This is done by induction on the rank of $`E`$. If rank $`E=1`$ then clearly we have $`c_1(^mE)=mc_1(E)`$. If rank $`E=2`$ we may formally split the bundle $`E`$ as direct sum of line bundles, i.e., we have an exact sequence: $$0L_1EL_20$$ so that there is a filtration $$^mE=F_0F_1\mathrm{}F_{m+1}=0$$ with $`F_i/F_{i+1}L_1^iL_2^{mi}`$ and so, by Whitney’s formula: $`c_1(^mE)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{m}{}}}c_1(F_i/F_{i+1})`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{m}{}}}c_1(L_1^iL_2^{mi})`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{m}{}}}ic_1(L_1)+{\displaystyle \underset{i=0}{\overset{m}{}}}(mi)c_1(L_2)`$ $`=`$ $`{\displaystyle \frac{m(m+1)}{2}}(c_1(L_1)+c_1(L_2))`$ $`=`$ $`{\displaystyle \frac{m(m+1)}{2}}c_1(E).`$ If rank $`E=3`$ then we split the bundle into a rank 2 bundle $`A`$ and a line bundle $`L`$, i.e., we have an exact sequence: $$0FEL0$$ so that there is a filtration $$^mE=F_0F_1\mathrm{}F_{m+1}=0$$ with $`F_i/F_{i+1}^iFL^{mi}`$ and the Chern number is given by: $`c_1(^mE)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{m}{}}}c_1(F_i/F_{i+1})`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{m}{}}}\{c_1(^iF)+(\mathrm{rank}^iF)(mi)c_1(L)\}`$ and by induction the RHS above is $$\underset{i=0}{\overset{m}{}}\frac{i(i+1)}{2}c_1(F)+\underset{i=0}{\overset{m}{}}(i+1)(mi)c_1(L)$$ hence $`c_1(^mE)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{m}{}}}{\displaystyle \frac{i(i+1)}{2}}c_1(F)+{\displaystyle \underset{i=0}{\overset{m}{}}}(i+1)(mi)c_1(L)`$ $`=`$ $`{\displaystyle \frac{1}{6}}{\displaystyle \frac{(m+2)!}{(m1)!}}c_1(F)+m{\displaystyle \underset{i=0}{\overset{m}{}}}(i+1)c_1(L){\displaystyle \underset{i=0}{\overset{m}{}}}i(i+1)c_1(L)`$ $`=`$ $`{\displaystyle \frac{1}{6}}{\displaystyle \frac{(m+2)!}{(m1)!}}c_1(F)+{\displaystyle \frac{m(m+1)}{2}}c_1(L){\displaystyle \frac{1}{3}}{\displaystyle \frac{(m+2)!}{(m1)!}}c_1(L)`$ $`=`$ $`{\displaystyle \frac{1}{6}}{\displaystyle \frac{(m+2)!}{(m1)!}}c_1(E).`$ Ne that we have used the formula: $$\underset{i=0}{\overset{m}{}}i(i+1)=\frac{1}{3}\frac{(m+2)!}{(m1)!}.$$ The general case is proved by induction using the formula: $$\underset{i=0}{\overset{m}{}}i(i+1)(i+1)\mathrm{}(i+k)=\frac{1}{k+2}\frac{(m+k+1)!}{(m1)!}.$$ QED Examples of surfaces of general type are provided by complete intersections in $`𝐏^n`$. A Smooth complete intersection of type $`(d_1,\mathrm{},d_{nr}),1rn1`$, in $`𝐏^n`$ is a smooth variety $`X`$ of dimension r which is the transversal intersections of $`(nr)`$ hypersurfaces of degree $`d_1,\mathrm{},d_{nr}`$ respectively. By the adjunction formula, the canonical bundle of a complete intersection $`X`$ of type $`(d_1,\mathrm{},d_{nr})`$ is given by the formula: $$𝒦_X=𝒪_{P^n}(d_1+\mathrm{}+d_{nr}(n+1))\|_X=𝒪_X(d_1+\mathrm{}+d_{nr}(n+1)).$$ The normal bundle $`𝒩_{X\|Y}`$ of a smooth hypersurface $`X`$ in a smooth variety $`Y`$ is given by $$𝒩_{X\|Y}=𝒪_Y(X)\|_X=𝒪_X(X).$$ Thus for a hypersurface $`X_1`$ of degree $`d_1`$ in $`𝐏^n`$, the normal bundle $$𝒩X_1\|P^n=𝒪_{P1}(d_1)\|_{X_1}=𝒪_{X_1}(d_1).$$ Inductively, for a smooth complete intersection $`X`$ of type $`(d_1,\mathrm{},d_{nr})`$ we get $$𝒩_{X\|P^n}=_{1inr}𝒪_X(d_i).$$ To compute Chern classes of $`X`$ we apply the Whitney formula to the exact sequence: $$0TXTP^n\|_X𝒩_{X\|P^n}=_{1inr}𝒪_X(d_i)0.$$ which yields the following formula for the total Chern classes: $$c(TX).c(𝒩_{X\|P^n})=c(T𝐏^n\|_X).$$ Operating symbolically, we get: $$1+c_1(TX)+\mathrm{}+c_r(TX)=(1+\theta )^{n+1}/\underset{1inr}{}(1+d_i\theta )$$ where $$\theta ^r=\underset{1inr}{}d_i.$$ Expanding formally the RHS above yields: $$(1+\theta )^{n+1}=1+C_1^{n+1}\theta +C_2^{n+1}\theta ^2+\mathrm{}+C_r^{n+1}\theta ^r$$ $$(1+d_i\theta )^1=1d_i\theta +(d_i\theta )^2\mathrm{}+(1)^r(d_i\theta )^r,1inr.$$ Define polynomials $`p_q(0qnr)`$ in $`d_1,\mathrm{},d_{nr}`$ by $`p_0(d_1,\mathrm{},d_{nr})=1`$, $$p_q(d_1,\mathrm{},d_{nr})=\underset{1i_1\mathrm{}i_qnr}{}d_{i_1}\mathrm{}.d_{i_q}1qnr.$$ Then for $`0qnr`$, the Chern classes of $`X`$ are given by: $`c_q(TX)={\displaystyle \underset{i=0}{\overset{q}{}}}(1)^iC_{qi}^{n+1}p_i(d_1,\mathrm{},d_{nr})\theta ^q.`$ For hypersurface ($`r=n1`$) the formulas above reduce to $$c_q(TX)=\underset{i=0}{\overset{q}{}}(1)^iC_{qi}^{n+1}d_i\theta ^q,0qn1.$$ For surfaces of complete intersections ($`r=2`$) and the formulas reduce to: $$c_1(TX)=((n+1)\underset{i=1}{\overset{n2}{}}d_i)\theta ,$$ $$c_2(TX)=\{\frac{n(n+1)}{2}(n+1)\underset{i=1}{\overset{n2}{}}d_i+\underset{1ijn2}{}d_id_j\}\theta ^2.$$ In particular, if $`d_1=\mathrm{}=d_{n2}=d`$ then $$c_1(X)=\{(n+1)(n2)d\}\theta ,$$ $$c_2(X)=\{\frac{n(n+1)}{2}(n+1)(n2)d+\frac{(n1)(n2)}{2}d^2\}\theta ^2.$$ For $`n=3`$ then $$c_1(TX)=(4d)\theta ,c_2(TX)=(64d+d^2)\theta ^2;$$ equivalently, for the cotangent bundle, we have: $$c_1(T^{}X)=(d4)\theta ,c_2(T^{}X)=(64d+d^2)\theta ^2.$$ We shall need a vanishing Theorem (see \[G-G\]) which is a consequence of a result of Bogomolov (\[B1\], \[B2\]): Theorem 4.4 Let $`X`$ be a minimal surface of general type and if the geometric genus $`p_g(X)>0`$ then $$H^2(X,𝒥_k^mX)=0$$ if $`k1`$ and $`m>2k`$. Actually, it was asserted in \[G-G\] that the preceding Theorem holds without the assumption that $`p_g(X)>0`$. At the momoent I can only get through the proof with this additional assumption. Thus the condition of Theorem 3.14 is satisfied for a minimal surface of general type and we obtained, Corollary 4.5 Let $`X`$ be a smooth minimal surface of general type with $`p_g(X)>0`$ and if $`c_1^2(𝒥_k^mX)c_2(𝒥_k^mX)>0`$ then there exists $`c>0`$ and $`m_0>0`$ such that $$dimH^0(X,𝒥_k^{k!m}X)cm^{n(k+1)1},$$ if $`k1`$ and $`m>1,`$ i. e., $`𝒥_k^{k!}X`$ is big. Corollary 4.6 Let $`X`$ be a smooth minimal surface of general type with $`p_g(X)>0`$ then $`𝒥_k^{k!}X`$ is big for $`k3`$. Proof. By the calculation in section 2, $$\mathrm{\Delta }(𝒥_3^6X)=c_1^2(𝒥_3^6X)c_2(𝒥_3^6X)=5175c_1^2(T^{}X)175c_2(T^{}X)$$ which is clearly $`>0`$ in view of Theorem 4.1. The Corollary now follows from Corollary 4.5. QED If $`X`$ is a smooth hypersurface the preceding Corollary can be expressed in terms of the degree: Corollary 4.6 Let $`X`$ be a non-singular hyper surface of degree $`d`$ in $`𝐏^3`$. Then $`𝒥_2^2X`$ is big if $`d9`$ and $`𝒥_3^6X`$ is big if $`d5`$. Proof. By the calculation in sectin 2, we have: $$c_1^2(𝒥_2^2X)c_2(𝒥_2^2X)=11c_1^2(T^{}X)5c_2(T^{}X)$$ and as noted before, for a smooth hypersurface $`X`$ in $`𝐏^3`$ of degree $`d`$, the Chern numbers are given by $$c_1(T^{}X)=d4,c_2(T^{}X)=d^24d+6,$$ we conclude that: $$c_1^2(𝒥_2^2X)c_2(𝒥_2^2X)=11d^2(d4)^25d^2(d^24d+6)>0$$ if $`d9`$. Computing similarly we conclude that $`c_1^2(𝒥_3^6X)c_2(𝒥_3^6X)>0`$ if $`d5`$. Moreover, by Noether;s Theorem (i. e., Riemann-Roch): $$1q(X)+p_g(X)=\frac{1}{12}(c_1^2(T^{}X)+c_2(T^{}X))$$ implies that $`p_g(X)>0`$ because the irregularity $`q(X)=0`$. QED We need one last observation to deal with the fact that $`𝒥_k^mX`$ is not semi-stable as can be seen from the calculation in section 2. In fact each of the factors $`^{i_1}T^{}X\mathrm{}^{i_k}T^{}X,i_1+2i_2+\mathrm{}+ki_k=m`$ which is a subsheaf of $`𝒥_k^mX`$ (note that not all of them are) is a destabling subsheaf). However, we also observe that each of these sheaves is semi-stable (by Theorem 4.3). Moreover the ratio: $`c_1(^{i_1}T^{}X\mathrm{}^{i_k}T^{}X)/\mathrm{rank}(^{i_1}T^{}X\mathrm{}^{i_k}T^{}X)m/2`$ thus we have: Theorem 4.7 Let $`X`$ be a complex surface such that $`PicX𝐙`$. If $`H^0(X,𝒥_k^mX[D])\{0\}`$ where $`D`$ is a divisor in $`X`$ then $$c_1([D]).c_1(T^{}X)\frac{m}{2}c_1^2(T^{}X).$$ Proof. This follows from the filtration: $$𝒥_{k1}^mX=_k^0_k^1\mathrm{}_k^{[m/k]}=𝒥_k^mX$$ $`(`$where $`[m/k]`$ is the greatest integer smaller than or equal to $`m/k)`$ such that $$_k^i/_k^{i1}𝒥_{k1}^{mki}X(^iT^{}X).$$ From the exact sequence $$0_k^{[m/k]1}𝒥_k^mX𝒥_{k1}^{mk[m/k]}X(^{[m/k]}T^{}X)0$$ we obtain an exact sequence: $`0H^0(X,_k^{[m/k]1}[D])H^0(X,𝒥_k^mX[D])`$ $``$ $`H^0(X,𝒥_{k1}^{mk[m/k]}X(^{[m/k]}T^{}X)[D])`$ which shows that if $`H^0(X,𝒥_k^mX[D])\{0\}`$ then either $`H^0(X,_k^{[m/k]1}[D])\{0\}`$ or $$H^0(X,𝒥_{k1}^{mk[m/k]}X(^{[m/k]}T^{}X)[D])\{0\}$$ and eventually this means that $$\mathrm{either}H^0(X,^{i_1}T^{}X\mathrm{}^{i_k}T^{}X[D])\{0\}$$ for at least one of the factors $`^{i_1}T^{}X\mathrm{}^{i_k}T^{}X,i_1+2i_2+\mathrm{}+ki_k=m`$. With this the Theorem follows from Theorem 4.3. QED From the computation in section 2 we see that $$\mu (𝒥_k^{k!}X)=\frac{c_1(𝒥_k^{k!}X)}{\mathrm{rank}𝒥_k^{k!}X}<\frac{m}{2}c_1(T^{}X).$$ Thus the estimate is weaker than one would get if it were stable, however this is the best that one can do and this weaker estimate is sufficient for our purpose. We assume from now on that $`(i)X`$ is a minimal surface of general type, (ii) $`p_g(X)>0`$ and $`(iii)Pic(X)𝐙`$. Then $`𝒥_k^{k!}X`$ is big for $`k3`$. This implies that $`𝒥_k^{k!}X[D]`$ is big for any effective ample divisor $`D`$ in $`X`$. Schwarz Lemma implies that the lifting of any holomorphic curves in $`𝐏(J^kX)`$ is contained in the zero set of a non-trivial section of $`𝒪(k!m)p^{}[D])`$. We proceed to consider the subvarieties of $`𝐏(J^kX),k2`$. Let $`Y_1`$ be an irreducible effective horizontal (i. e., not of the form $`p^{}D`$ for some effective divisor $`D`$ in $`X`$) divisor in $`𝐏(J^kX)`$ then: $$[Y_1]=𝒪_{𝐏(J^kX)}(m_1)p^{}[D_1]$$ where $`D_1`$ is a divisor in $`X`$ and $`m_1𝐍`$, we may assume that $`m_1`$ is divisible by $`k!`$ by replacing $`Y_1`$ with $`k!Y_1`$ (so it is non-reduced but set theoretically it has only one irreducible component). Thus we may write $`m_1=k!\alpha _1`$. For simplicity of notations we shall write $`𝒪(j)`$ instead of $`𝒪_{𝐏(J^kX)}(j)`$. Since dim $`𝐏(J^kX)=2(k+1)1=2k+1`$ we get from Theorem 4.7 and Theorem 3.13: $`c_1^{2k}(𝒪(k!)\|_{Y_1})`$ $`=`$ $`c_1^{2k+1}(𝒪(k!)).(c_1(𝒪(k!\alpha _1)p^{}c_1([D_1])`$ $`=`$ $`(k!)^2\{\alpha _1c_1^{2k+1}(𝒪(k!))c_1^{2k}(𝒪(k!)).p^{}c_1([D_1])\}`$ $`=`$ $`(k!)^2\{\alpha _1(c_1^2(𝒥_k^{k!}X)c_2(𝒥_k^{k!}X))c_1(𝒥_k^{k!}X).c_1([D_1])\}`$ $`=`$ $`(k!)^2\{\alpha _1(c_1^2(𝒥_k^{k!}X)c_2(𝒥_k^{k!}X))a(k,k!)c_1(T^{}X).c_1([D_1])\}`$ $``$ $`(k!)^2\{\alpha _1(a(k,k!)^2c_1^2(T^{}X)c_2(𝒥_k^{k!}X)){\displaystyle \frac{\alpha _1k!a(k,k!)}{2}}c_1^2(T^{}X)\}`$ $`=`$ $`(k!)^2\alpha _1\{a(k,k!)(a(k,k!){\displaystyle \frac{k!}{2}})c_1^2(T^{}X)c_2(𝒥_k^{k!}X)\}.`$ This means that $`𝒪(k!)\|_{Y_1}`$ is again big if $$a(k,k!)(a(k,k!)\frac{k!}{2})c_1^2(T^{}X)c_2(𝒥_k^{k!}X)>0.$$ For example if $`k=3,a(k,k!)=101`$ by the calculation in section 2; the preceding inequality yields: $`c_1^{2k}(𝒪(k!)\|_{Y_1})`$ $``$ $`(k!)^2\alpha _1\{101(1013)c_1^2(T^{}X)c_2(𝒥_k^{k!}X)\}`$ $`=`$ $`(k!)^2\alpha _1\{(92925026)c_1^2(T^{}X)c_2(T^{}X)\}`$ $`=`$ $`(k!)^2\alpha _1(4266c_1^2(T^{}X)175c_2(T^{}X))`$ $`>`$ $`0.`$ This means that $`𝒪(k!)\|_{Y_1}`$ is again big. The Schwarz Lemma in the appendix again implies that the lifting of any holomorphic curves in $`𝐏(J^kX)`$ is contained in the zero set of a non-trivial section of $`𝒪(k!m)\|_{Y_1}p\|_{Y_1}^{}[D])`$. Next we consider divisor $`Y_2`$ in $`Y_1`$ which is of the form: $$[Y_2]=(𝒪(m_2)p^{}[D_2])\|_{Y_1}$$ where $`D_2`$ is a divisor in $`X`$ and $`m_2𝐍`$ which we may assume to be divisible by $`k!`$, i.e., $`m_2=\alpha _2k!`$. We remark that for the investigation of degeneration of liftings of a holomorphic curve $`f:𝐂X`$ these are the only type of subvarieties that we have to deal with. We have: $`c_1^{2k1}(𝒪(k!)\|_{Y_2})`$ $`=`$ $`c_1^{2k1}(𝒪(k!)).(c_1(𝒪(k!\alpha _1)p^{}c_1([D_1]).(c_1(𝒪(k!\alpha _2)p^{}c_1([D_2])`$ $`=`$ $`\alpha _1\alpha _2c_1^{2k+1}(𝒪(k!))\alpha _1c_1^{2k}(𝒪(k!)).p^{}c_1([D_2])`$ $`\alpha _2c_1^{2k}(𝒪(k!)).p^{}c_1([D_1])+p^{}c_1([D_1]).p^{}c_1([D_2])`$ $``$ $`(k!)^2\{\alpha _1\alpha _2(c_1^2(𝒥_k^{k!}X)c_2(𝒥_k^{k!}X))\alpha _1c_1(𝒥_k^{k!}X).c_1([D_1])`$ $`\alpha _2\alpha _1c_1(𝒥_k^{k!}X).c_1([D_1])+c_1([D_1]).c_1([D_2])\}`$ $`=`$ $`(k!)^2\{\alpha _1\alpha _2(c_1^2(𝒥_k^{k!}X)c_2(𝒥_k^{k!}X))\alpha _1a(k,k!)c_1(T^{}X).c_1([D_2])`$ $`\alpha _2a(k,k!)c_1(T^{}X).c_1([D_1])+c_1([D_1]).c_1([D_2])\}`$ $``$ $`(k!)^2\{\alpha _1\alpha _2(c_1^2(𝒥_k^{k!}X)c_2(𝒥_k^{k!}X))2{\displaystyle \frac{\alpha _1\alpha _2k!a(k,k!)}{2}}c_1^2(T^{}X)\}`$ $`=`$ $`(k!)^2\alpha _1\alpha _2\{a(k,k!)(a(k,k!)k!)c_1^2(T^{}X)c_2(𝒥_k^{k!}X)\}.`$ Proceeding inductively, we get a sequence of subvarieties $`Y_1Y_2\mathrm{}Y_{2k}`$ where each $`Y_i`$ is of codimension $`i`$ and of the form $$[Y_{i+1}]=(𝒪(m_{i+1})p^{}[D_{i+1}])\|_{Y_i}.$$ A similar calculation shows that: $`c_1^{2ki+1}(𝒪(k!)\|_{Y_i})`$ $``$ $`(k!)^2\alpha _1\mathrm{}\alpha _i\{a(k,k!)(a(k,k!)k!{\displaystyle \frac{i}{2}})c_1^2(T^{}X)c_2(𝒥_k^{k!}X)\},`$ $`i=1,\mathrm{},2k`$. For $`k=3`$ we have: $`a(3,3!)(a(3,3!)3!{\displaystyle \frac{i}{2}})c_1^2(T^{}X)c_2(𝒥_k^{k!}X)`$ $``$ $`101(101(3!)3)c_1^2(T^{}X)c_2(𝒥_k^{k!}X)`$ $`=`$ $`3357c_1^2(T^{}X)175c_2(T^{}X)`$ $`>`$ $`0`$ for all $`i=1,\mathrm{},2k=6`$. Thus we arrive at the following Theorem: Theorem 4.8 Let $`X`$ be a smooth minimal surface of general type such that $`(i)PicX𝐙`$ and $`(ii)p_g(X)>0`$. Then $`X`$ is hyperbolic. Consequently, a generic smooth hypersurface in $`𝐏^3`$ of degree $`d5`$ is hyperbolic. In \[D-E\] certain types of 2-jet differentials $`𝒜`$ were used and the authors establised that $`c_1^2(𝒜)c_2(𝒜)=13c_1^2(T^{}X)9c_2(T^{}X)`$ on any hypersurface $`X`$ of degree $`42`$. This is weaker than what we have, namely $`c_1^2(𝒥_2^2X)c_2(𝒥_2^2X)=11c_1^2(T^{}X)5c_2(T^{}X)`$. Actually I have some trouble using this stronger estimate to get hyperbolicity due to non-semistability (recall that neither $`𝒥_k^m`$ nor $`T_k^{}X`$ is semi-stable) so that $`c_1^2(𝒥_2^2X)c_2(𝒥_2^2X)=11c_1^2(T^{}X)5c_2(T^{}X)`$ is not big enough to reach hyperbolicity. It appears that $`𝒜`$ is not semi-stable either. Appendix A: The Lemma of Logarithmic Derivatives One of the main tool in Nevanlinna Theory is the classical Lemma of Logarithmic Derivatives (abbrev. LLD). For example, LLD implies that, even though there is no pointwise estimate between (the absolute value of) a holomorphic function and (the absolute value of) its derivatives, such estimates do exist in the sense of integral averages (i.e, their characteristic functions bound each other). The purpose of this appendix is to extend the classical LLD to all jet differentials of logarithmic type and in particular all regular jet differentials. The proof is based on the very simple observation that (the absolute value of) any jet differential of logarithmic type is bounded by (the absolute value of) those of the classical type (hence the classical LLD applies). Theorem A1 (Lemma of Logarithmic Derivatives) Let $`X`$ be a projective variety and let $`(i)D`$ be an effective divisor with simple normal crossings or $`(ii)D`$ is the trivial divisor in $`X(`$i.e. the support of $`D`$ is empty or equivalently, the line bundle associate to $`D`$ is trivial$`)`$. Let $`f:𝐂X`$ be a holomorphic map and $`\omega H^0(X,𝒥_k^mX(\mathrm{log}D))(`$resp. $`H^0(X,𝒥_k^mX)`$ in case $`(ii))`$ a jet differential such that $`\omega j^kf`$ is not identically zero, then $$T_{\omega j^kf}(r)=_0^{2\pi }\mathrm{log}^+\|\omega (j^kf(re^{\sqrt{1}\theta }))\|\frac{d\theta }{2\pi }..O(\mathrm{log}T_f(\omega _X;r))+O(\mathrm{log}r).$$ Here $`\omega _X`$ can be taken to be $`c_1()`$ of any ample line bundle $``$ on $`X`$. Proof. We claim that there exists a finite number of rational functions $`t_1,\mathrm{},t_q`$ on $`X`$ such that: > $`()`$ the logarithmic jet differentials $`\{(d^{(j)}t_i/t_i)^{m/j}\|1iq,1jk\}`$ span the fibers of $`𝒥_k^mX(\mathrm{log}D)`$ over every point of $`X`$. Without loss of generality we may assume that $`D`$ is ample; otherwise we may replace $`D`$ be $`D+D^{^{}}`$ so that $`D+D^{^{}}`$ is ample. Observe that if $`s`$ is a function holomorphic on a neighborhood $`U`$ such that $`[s=0]=DU`$ then $`[s^\tau =0]=\tau DU`$ where $`\tau `$ is a rational number. Thus $`d^{(j)}(\mathrm{log}s^\tau )=\tau d^{(j)}(\mathrm{log}s)`$ is still a jet differential with logarithmic singularity along $`DU`$ so the mutiplicity causes no problem. This means that we may assume without loss of generality that $`D`$ is very ample by replacing $`D`$ with $`\tau D`$ for some $`\tau `$ so that $`\tau D`$ is very ample. Let $`uH^0(X,[D])`$ be a section such that $`D=[u=0]`$. At a point $`xD`$ choose a section $`v_1H^0(X,[D])`$ so that $`E_1=[v_1=0]`$ is smooth, $`D+E_1`$ is of simple normal crossings and $`v_1`$ is non-vanishing at $`x`$ (this is possible because the line bundle $`[D]`$ is very ample). The rational function $`t_1=u_1/v_1`$ is regular on the affine open neighborhood $`XE_1`$ of $`x`$ and $`(XE_1)[t_1=0]=(XE_1)D`$. Choose rational functions $`t_2=u_2/v_2,\mathrm{},t_n=u_n/v_n`$ where $`u_i`$ and $`v_i`$ are sections of a very ample bundle $``$ so that $`t_2,\mathrm{},t_n`$ are regular at $`x`$, the divisors $`D_i=[u_i=0],E_i=[v_i=0]`$ are smooth and $`D+D_2+\mathrm{}+D_n+E_1+\mathrm{}+E_n`$ is of simple normal crossings. Moreover, since the bundles involved are very ample the sections can be chosen so that $`dt_1\mathrm{}dt_n`$ is non-vanishing at $`x`$; the complete system of sections provides an embedding, hence at each point there are $`n+1`$ sections with the property that $`n`$ of the the quotients of these $`n+1`$ sections form a local coordinate system on some open neighborhood $`U_x`$ of $`x`$. This implies that $`()`$ is satisfied over $`U_x`$. Since $`D`$ is compact it is covered by a finite number of such open neighborhoods, say $`U_1,\mathrm{},U_p`$ and a finite number of rational functions (constructed as above for each $`U_i`$) on X so that $`()`$ is satisfied on $`_{1ip}U_i`$. Moreover, there exists relatively compact open subsets $`U_i^{}`$ of $`U_i`$ ($`1ip`$) such that $`_{1ip}U_i^{}`$ still covers $`D`$. Next we consider a point $`x`$ in the compact set $`X_{1ip}U_i^{}`$. Repeating the procedure as above we can find rational functions $`s_1=a_1/b_1,\mathrm{},s_n=a_n/b_n`$ where $`a_i`$ and $`b_i`$ are sections of some very ample line $``$ bundle so that $`s_1,\mathrm{},s_n`$ from a holomorphic local coordinate on some open neighborhood $`V_x`$ of $`x`$. Thus $`()`$ is satisfied on $`V_x`$ by the rational functions $`s_1,\mathrm{},s_n`$. Note that we must also choose these sections so that the divisor $`H=[s_1\mathrm{}s_n=0]`$ together with those divisors (finite in number), which had been already constructed above, is still a divisor with simple normal crossings (this is possible by the very ampleness of the line bundle $``$.) Since $`X_{1ip}U_i^{}`$ is compact, it is cover by a finite of such coordinate neighborhoods. The coordinates are rational functions and finite in number and by construction it is clear that the condition $`()`$ is satisfied on $`X_{1ip}U_i^{}`$. Since $`_{1ip}U_i`$ together with $`X_{1ip}U_i^{}`$ covers $`X`$, the condition $`()`$ is satisfied on $`X`$. If $`D`$ is the trivial divisor, then it is enough to use only the second part of the construction above and again $`()`$ is verified with $`𝒥_k^mX(\mathrm{log}D)=𝒥_k^mX`$ To obtain the estimate of the Theorem observe that the function, $$\rho :J^kX(\mathrm{log}D)[0,\mathrm{}]$$ defined by $`\rho (\xi )={\displaystyle \underset{i=1}{\overset{q}{}}}{\displaystyle \underset{j=1}{\overset{k}{}}}\|(d^{(j)}t_i/t_i)^{m/j}(\xi )\|^2,\xi J^kX(\mathrm{log}D)`$ ($`\{t_i\}`$ is the family of rational functions satisfying the condition $`()`$) is continuous in the extended sense; it is continuous, in the ususal sense, outside the fibers over the divisor $`E`$ (the sum of the divisors associated to the rational functions $`\{t_i\}`$; note that $`E`$ contains $`D`$). Over the fiber of each point $`xXE`$, $`\|(d^{(j)}t_i/t_i)^{m/j}(\xi )\|^2`$ is finite for $`\xi J^kX(\mathrm{log}D)_x`$, thus $`\rho `$ is not identically infinite. Moreover, since $$\{(d^{(j)}t_i/t_i)^{m/j}\|1iq,1jk\}$$ span the fiber of $`𝒥_k^mX(\mathrm{log}D)`$ over every point of $`X`$, $`\rho `$ is strictly positive ($`+\mathrm{}`$ allowed) outside the zero-section of $`J_kX(\mathrm{log}D)`$. The quotient $$\|\omega \|^2/\rho :J^kX(\mathrm{log}D)[0,\mathrm{}]$$ does not take on the extended value $`\mathrm{}`$ when restricted to $`J^kX(\mathrm{log}D)\{zerosection\}`$ becuase, as we have just observed, $`\rho `$ is non-vanishing (even though it blows up along the fibers over $`E`$ so that the reciprocal $`1/\rho `$ is zero there) and the singularity of $`\|\omega \|`$ is no worst then that of $`\rho `$ becuase the singularity of $`\omega `$ ocuurs only along $`D`$ (which is contained in $`E`$) and is of log type. Thus the restriction to $`J_kX(\mathrm{log}D)\{zerosection\}`$, $$\|\omega \|^2/\rho :J^kX(\mathrm{log}D)\{zerosection\}[0,\mathrm{})$$ is a continuous non-negative funtion. Moreover, since $`\|\omega \|`$ and $`\rho `$ have the same homogenity: $$\|\omega (\lambda .\xi )\|^2=\|\lambda \|^{2m}\|\omega (\lambda .\xi )\|^2and\rho (\lambda .\xi )=\|\lambda \|^{2m}\rho (\xi )$$ for all $`\lambda 𝐂^{}`$ and $`\xi J^kX(\mathrm{log}D)`$ we see that $`\|\omega \|^2/\rho `$ descends to a well-defined function on $`𝐏(E_{k,D})=(J^kX(\mathrm{log}D)\{zerosection\})/𝐂^{}`$, i.e., $$\|\omega \|^2/\rho :𝐏(E_{k,D})[0,\mathrm{})$$ is a well-defined continuous function and so, by compactness, there exists a constant $`c`$ with the property that $`\|\omega \|^2c\rho .`$ This implies that $`T_{\omega j^kf}(r)`$ $`=`$ $`{\displaystyle _0^{2\pi }}\mathrm{log}^+\|\omega (j^kf(re^{\sqrt{1}\theta }))\|{\displaystyle \frac{d\theta }{2\pi }}`$ $``$ $`{\displaystyle _0^{2\pi }}\mathrm{log}^+\|\rho (j^kf(re^{\sqrt{1}\theta }))\|{\displaystyle \frac{d\theta }{2\pi }}+O(1).`$ Since $`t_i`$ is a rational function on $`X`$ the function $$(d^{(j)}t_i/t_i)^{m/j}(j^kf)=((t_if)^{(j)}/t_if)^{m/j}$$ ($`m`$ is divisible by $`k!`$) is meromorphic on $`𝐂`$ and so, by the definition of $`\rho `$, $$\mathrm{log}^+\|\rho (j^kf)\|O(\underset{1iq,1jk}{\mathrm{max}}\mathrm{log}^+\|(t_if)^{(j)}/t_if\|)+O(1).$$ Now by the classical lemma of logarithmic derivatives for meromorphic functions, $$_0^{2\pi }\mathrm{log}^+\|(t_if)^{(j)}/t_if\|)\frac{d\theta }{2\pi }..O(\mathrm{log}rT_{t_if}(r)).$$ Since $`t_i`$ is a rational function, $$\mathrm{log}T_{t_if}(r)O(\mathrm{log}T_f(\omega _X;r))+O(1)$$ and we arrive at the estimate $`{\displaystyle _0^{2\pi }}\mathrm{log}^+\|\rho (j^kf(re^{\sqrt{1}\theta })\|){\displaystyle \frac{\theta }{2\pi }}`$ $`O({\displaystyle _0^{2\pi }}\mathrm{log}^+\|(t_if)^{(j)}/t_if\|{\displaystyle \frac{d\theta }{2\pi }})+O(1)`$ $`..O(\mathrm{log}T_f(r))+O(\mathrm{log}r).`$ This implies that $$T_{\omega j^kf}(r)..O(\mathrm{log}T_f(r))+O(\mathrm{log}r)$$ as claimed. QED We obtain, as imediate consequence, the following Schwarz’s type Lemma for logarithmic jet differentials. Corollary A2 Let $`X`$ be a projective variety and $`D`$ be an effective divisor $`(`$possibly the trivial divisor$`)`$ with simple normal crossings. Let $`f:𝐂XD`$ be a holomorphic map. Then $$\omega (j^kf)0\mathrm{for}\mathrm{all}\omega H^0(X,𝒥_k^mX(\mathrm{log}D)[H])$$ where $`H`$ is a generic hyperplane section $`(`$and hence any hyperplane section$`)`$. Proof. If $`f`$ is constant then the Corollary holds trivially. So we may assume that $`f`$ is non-constant and suppose that $`\omega j^kf0`$. Moreover, since $`F`$ is non-constant, we may assume without loss of generality that $`\mathrm{log}r=o(T_f(H;r))`$ by replacing $`f`$ with $`f\varphi `$ where $`\varphi `$ is a transcendental function on $`𝐂`$. By Theorem 4.1, we have $$_0^{2\pi }\mathrm{log}^+\|\omega j^kf\|\frac{d\theta }{2\pi }=T_{\omega j^kf}(r)..O(\mathrm{log}rT_f(H;r)).$$ On the other hand, since $`\omega `$ vanishes on $`H`$ and $`H`$ is generic (see (1) or (2) in section 1), we obtain via Jensen’s Formula: $`T_f(H;r)`$ $``$ $`N_f(H;r)+O(\mathrm{log}rT_f(H;r))`$ $`=`$ $`{\displaystyle _0^{2\pi }}\mathrm{log}\|\omega j^kf\|{\displaystyle \frac{d\theta }{2\pi }}+O(\mathrm{log}rT_f(H;r))`$ which, together with the preceding estimate, imply that: $$T_f(H;r)..O(\mathrm{log}rT_f(H;r)).$$ This is impossible hence we must have $`\omega j^kf0`$. If $`H_1=[s_1=0]`$ is any hyperplane section then it is linearly equivalent to a generic hyperplane section $`H=[s=0]`$. If $`\omega `$ vanishes along $`H^{^{}}`$ then $`(s/s_1)\omega `$ vanishes along $`H`$. The preceding discussion implies that $`(s/s_1)\omega (j^kf)0`$. This implies that actually $`\omega (j^kf)0`$ as we may choose a generic section $`H`$ so that the image of $`f`$ is not entirely contained in $`H`$. QED Actually the proof of Theorem A.1 gives a little more. In fact the same proof yields: Theorem A3 Let $`\rho _k`$ be a pseudo singular jet metric on $`J^kX(\mathrm{log}D)`$ with the property that there exists a constant $`c>0`$ such that $`\rho _k\rho `$ where $`\rho `$ is the singular jet metric on $`J^kX(\mathrm{log}D)`$ defined by the family of rational functions $`()(`$see $`(28))`$. Then $$T_{j^kf}(\rho _k;r)=_0^{2\pi }\mathrm{log}^+\|\rho _k(j^kf(re^{\sqrt{1}\theta }))\|\frac{d\theta }{2\pi }..O(\mathrm{log}rT_f(\omega _X;r))$$ for any Kähler mertic $`\omega _X`$ on $`X`$. In particular, if $`\rho _k`$ is a non-singular pseudo metric on $`J^kX`$ then the preceding estimate holds. The Schwarz Lemma can be further extended as follows. Theorem A4 Let $`Y𝐏(J^kX)`$ be a subvariety and suppose that there exists a non-trivial section $`\sigma H^0(Y,𝒪_{𝐏(J^kX)})(m)\|_Yp\|_Y^{}[D]`$ where $`D`$ is a generic ample divisor in $`X`$ and $`p:𝐏(J^kX)X`$ is the projection map. If the image of the lifting $`[j^kf]:𝐂𝐏(J^kX)`$ of a holomorphic curve $`f:𝐂X`$ is contained in $`Y`$ then $`\sigma ([j^kf])0`$. Appendix B Let $`S_n`$ be the symmetric group on $`n`$ elements then the order of $`S_n`$ is $`n!`$. Definition B1 A maximal set of mutually conjugate elements of $`S_n`$ is said to be a class of $`S_n`$. Definition B2 A partition of a natural number $`n`$ is a set of positve integers $`k_1,\mathrm{},k_q`$ such that $`n=k_1+\mathrm{}+k_q`$. A partition can be expressed as $$n=\underset{i=1}{\overset{n}{}}ia_i$$ where the integers $`a_i`$ are non-negative. Theorem B3 The number, denoted $`p(n)`$, of classes of $`S_n`$ is equal to the number of partitions of $`n`$ and also to the number of inequivalent irreducible representations of $`S_n`$. The number $`p(n)`$ is asymptotically approximated by the formula of Hardy-Ramanujan $$p(n)\frac{e^{\pi \sqrt{2n/3}}}{4n\sqrt{3}}.$$ The alternating subgroup $`A_n`$ (i.e. the even permutations) is the commutator subgroup of $`S_n`$ and is obviously of index 2. Thus there are two 1-dimensional representations of $`S_n`$: the trivial representation and the representation $`\mathrm{\Gamma }_\sigma `$ defined by $`P\sigma (P)`$ where $`\sigma `$ is the signature of a permutation $`P`$ (i.e. $`P\pm 1`$ depending on whether $`P`$ is even or odd (i.e., can be expressed as an even or odd number of transpositions: interchanging two of the $`n`$ elements). Lemma B4 Let $`X=X(n)`$ be a set of $`n`$ elements and let $`Y_1,\mathrm{},Y_k`$ be $`k`$ not necessarily distinct subsets of $`X`$. For any subset $`J`$ of the index set $`\{1,\mathrm{},k\}`$, denote by $$n(J)=\mathrm{\#}_{jJ}Y_j;$$ and for $`0ik`$, denote by $$n_0=n,n_i=\underset{\mathrm{\#}J=i}{}n(J),1ik.$$ Then the number of elements not in any of the subsets $`Y_i,i=1,\mathrm{},k`$ is given by the formula $$\mathrm{\#}(X(_{i=1}^kY_i))=nn_1+n_2\mathrm{}+(1)^kn_k=\underset{i=0}{\overset{k}{}}(1)^in_i.$$ Proof. If $`k=2`$ the formula can be expressed as usual: $$\mathrm{\#}(X(Y_1Y_2))=\mathrm{\#}X\mathrm{\#}Y_1\mathrm{\#}Y_2+\mathrm{\#}(Y_1Y_2).$$ One way to prove the Lemma is by induction on $`k`$. Alternatively one can also argue as follows: QED An element $`P`$ of $`𝒮_n`$ is said to be a derangement if $`P(i)i`$ for $`i=1,\mathrm{},n`$. The number of derangements is denoted by $`d_n`$. Then Corollary B5 The number of derangements is given by the formula $$d_n=n!\underset{i=0}{\overset{n}{}}\frac{(1)^i}{i!}.$$ In particular we see that asymptotically $`d_ne^1`$. Proof. Apply Lemma 4 with $`X=𝒮_n`$ and $`Y_i=\{P𝒮_n\|P(i)=i\},i=1,\mathrm{},n.`$ Alternatively, the formula can be obtained by considering the power series: $$e^x\underset{i=0}{\overset{\mathrm{}}{}}d_i\frac{x^i}{i!}=\underset{i=0}{\overset{\mathrm{}}{}}(\underset{j=0}{\overset{i}{}}\frac{i!}{j!(ij)!}d_{ij})\frac{x^i}{i!}=\underset{i=0}{\overset{\mathrm{}}{}}x^i=\frac{1}{1x}.$$ Thus we have $$\underset{i=0}{\overset{\mathrm{}}{}}d_i\frac{x^i}{i!}=e^x(1x)^1$$ which yields the formual of the corollary. QED The formula of the Corollary can also be obtained via the recursive formula: $$d_nnd_{n1}+(1)^n.$$ Corollary B6 The number of surjections from a set $`A`$ of $`n`$ elements to a set $`B`$ of $`k`$ elements is given by the formula $$\underset{i=0}{\overset{k}{}}(1)^i\frac{k!}{i!(ki)!}(ki)^n.$$ Proof. Apply Lemma 4 to the set $`X`$ of all maps from $`A`$ to $`B`$ and $`Y_i,=1,\mathrm{},k`$ be the subset consisting of those maps such that $`i`$ is not in the image. QED Note that Corollary 6 implies trivially that $`{\displaystyle \underset{i=0}{\overset{k}{}}}(1)^i{\displaystyle \frac{k!}{i!(ki)!}}(ki)^n=\{\begin{array}{cc}n!\hfill & \text{if }k=n,\hfill \\ 0\hfill & \text{if }k>n\text{.}\hfill \end{array}`$ There is a more general formula which can be proved in a similar fashion: $`{\displaystyle \underset{i=0}{\overset{n}{}}}(1)^i{\displaystyle \frac{n!}{i!(ni)!}}{\displaystyle \frac{(m+ni)!}{(m+nk)!(ki)!}}=\{\begin{array}{cc}m!/k!(mk)!\hfill & \text{if }mk,\hfill \\ 0\hfill & \text{if }m<k\text{.}\hfill \end{array}`$ Theorem B7 The number of non-negative integer solutions of the equation $$x_1+\mathrm{}+x_k=n$$ is $`(n+k1)!/(k1)!n!.`$ On the other hand the number of positive integer solutions is $`(n1)!/(k1)!(nk)!`$. Proof. So we have to find the number of ways to put $`n`$ black (otherwise identical) balls in $`k`$ slots. If we insert white balls in between the slots we end up with a total of $`n+k1`$ balls $`k1`$ of them white. This is the same as choosing $`k1`$ balls from a total of $`n+k1`$ balls and the first assertion follows. The second assertion follows from the first by making the substitution $`y_i=x_i1`$. resulting in the equation $$y_1+\mathrm{}+y_k=nk.$$ QED The number $`(n+k1)!/(k1)!n!`$ is the coefficient of $`x^n`$ in the expansion of the function $`{\displaystyle \frac{1}{(1x)^k}}={\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}c_nx^n.`$ (53) Let $`c_{n,k}`$ be the number of elements of $`𝒮_n`$ consisting of exactly $`k`$ cycles. Theorem B8 With the notations above we have $$c_{n,k}=(n1)c_{n1,k}+c_{n1,k1}$$ and these numbers are the coefficients of the expansion of the function $`x(x+1)\mathrm{}(x+n1)`$ $$x(x+1)\mathrm{}(x+n1)=\underset{k=0}{\overset{n}{}}c_{n,k}x^k$$ and also $$\frac{x!}{(xn)!}=\underset{k=0}{\overset{n}{}}(1)^{nk}c_{n,k}x^k.$$ Moreover these numbers are the coefficients of the expansion of the function $$\mathrm{log}(1+x)^k=k!\underset{n=k}{\overset{\mathrm{}}{}}c_{n,k}\frac{x^n}{n!}.$$ Proof. The recursive relation follows from the observation that there are exactly $`n1`$ different ways to get a permutation on $`n`$ elements consisting of exactly $`k`$ cycles from a permutation on $`n1`$ elements consisting of exactly $`k`$ cycles. These account for the first term on the right of the recursive formula. Next we observe that there is exactly one way to get a permutation on $`n`$ elements consisting of exactly $`k`$ cycles from a permutation on $`n1`$ elements consisting of exactly $`k1`$ cycles and these account for the second term in the formula. The rest of the Theorem follows from the observation that if we write $$x(x+1)\mathrm{}(x+n1)=\underset{k=0}{\overset{n}{}}a_{n,k}x^k$$ then the coefficients satisfy the same recursive formula as $`c_{n,k}`$: $$a_{n,k}=(n1)a_{n1,k}+a_{n1,k1}.$$ The last assetion follows by observing that $$(1+x)^t=e^{t\mathrm{log}(1+x)}=\underset{k=0}{\overset{\mathrm{}}{}}\frac{1}{k!}t^k(\mathrm{log}(1+x))^k.$$ On the other hand, we have $`(1+x)^t`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{t!}{n!(tn)!}}x^n`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{x^n}{n!}}{\displaystyle \underset{j=0}{\overset{n}{}}}c_{n,j}t^j`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}t^j{\displaystyle \underset{n=j}{\overset{\mathrm{}}{}}}c_{n,j}{\displaystyle \frac{x^n}{n!}}.`$ QED Denote by $`P(n,k)`$ the set of all partitions of a set of $`n`$ elements into $`k`$ non-empty subsets and let $$p_{n,k}=\mathrm{\#}P(n,k).$$ Theorem B9 With the notations above we have $$p_{n,k}=kp_{n1,k}+p_{n1,k1}$$ and these numbers are the coefficients of the expansion of the function $$x^n=\underset{k=0}{\overset{n}{}}p_{n,k}\frac{x!}{(xk)!}$$ and also as the coefficients of the power series expansion $$(e^x1)^k=k!\underset{n=k}{\overset{\mathrm{}}{}}p_{n,k}\frac{x^n}{n!}.$$ Proof. The recursive formula follows from the observation that a partition of $`n`$ elements into $`k`$ subsets can be obtained from a partition of $`n1`$ elements into $`k`$ subsets by insertingthe element $`n`$ into any one of the $`k`$ subsets. Alternatively one can also get a partition of $`n`$ elements into $`k`$ subsets from a partition of $`n1`$ elements into $`k1`$ subsets by simply adding one more subset consisting of just the element $`n`$. For a positive integer $`x`$ there are exactly $`x^n`$ maps from the set $`\{1,\mathrm{},n\}`$ of $`n`$ elements to the set $`\{1,\mathrm{},x\}`$. On the other hand, by definition of the number $`p_{n,k}`$ we have the relation: $`k!p_{n,k}=\mathrm{\#}`$ of surjections from a set of $`n`$ element onto a set of $`k`$ elements. Hence for any subset $`Y`$ of $`k`$ elements of $`\{1,\mathrm{},x\}`$ there are $`k!p_{n,k}`$ surjections from $`\{1,\mathrm{},n\}`$ onto the set $`Y`$. Since the number of subsets of $`k`$ elements of $`\{1,\mathrm{},x\}`$ is $`x!/k!(xk)!`$ we get $$x^n=\underset{k=0}{\overset{n}{}}\frac{x!}{k!(xk)!}k!p_{n,k}=\underset{k=0}{\overset{n}{}}\frac{x!}{(xk)!}p_{n,k}.$$ By Corollary 4 we have: $`k!p_{n,k}={\displaystyle \underset{i=0}{\overset{k}{}}}(1)^i{\displaystyle \frac{k!}{i!(ki)!}}(ki)^n={\displaystyle \underset{i=1}{\overset{k}{}}}(1)^{ki}{\displaystyle \frac{k!}{i!(ki)!}}i^n.`$ (54) If $`k=1`$ then $`p_{n,1}=1`$ as there is only one such partition. The usual expansion of the exponential function yields $$e^x1=\underset{n=1}{\overset{\mathrm{}}{}}\frac{x^n}{n!}.$$ The case of general $`k`$ can be verified by induction by differentiating the power series $$F_k(x)=\underset{n=k}{\overset{\mathrm{}}{}}s_{n,k}\frac{x^n}{n!}$$ resulting in $`F_k^{^{}}(x)`$ $`=`$ $`{\displaystyle \underset{n=k}{\overset{\mathrm{}}{}}}p_{n,k}{\displaystyle \frac{x^{n1}}{(n1)!}}`$ $`=`$ $`{\displaystyle \underset{n=k}{\overset{\mathrm{}}{}}}(kp_{n1,k}+p_{n1,k1}){\displaystyle \frac{x^{n1}}{(n1)!}}`$ $`=`$ $`kF_k(x)+F_{k1}(x).`$ Bu induction hypothesis we have: $$F_{k1}(x)=\frac{1}{(k1)!}(e^x1)^{k1}$$ hence the function $`F_k`$ satisfies the differential equation $$F_k^{^{}}(x)=kF_k(x)+\frac{1}{(k1)!}(e^x1)^{k1}.$$ It is clear that $$F_k(x)=\frac{1}{k!}(e^x1)^k$$ is a solution and is indeed the uniqe solution satisfying $`p_{k,k}=1`$. QED Theorem B10 The number of partitions of $`n`$ $$p(n)=\underset{k=1}{\overset{n}{}}\frac{k^n}{k!}.$$ Denote by $`p_k(n)`$ the number of solutions of the equation $`x_1+\mathrm{}+x_k=n`$ (55) with the condition that $`1x_kx_{k1}\mathrm{}x_1`$. This number is obviously equal to the number of solutions of the equation $`y_1+\mathrm{}+y_k=nk`$ (56) with the condition that $`0y_ky_{k1}\mathrm{}y_1`$. If there are exactly $`i`$ of the integers $`y_i`$ which are positive then these are the solutions of $`x_1+\mathrm{}+x_i=nk`$ and so there are $`p_i(nk)`$ of such solutions; consequently we have: Theorem B11 With the notations above we have $$p_k(n)=\underset{i=1}{\overset{k}{}}p_i(nk).$$ Consider the case $`k=3`$ then the number of solutions of $$x_1+x_2+x_3=n$$ such that $`0x_3x_2x_1`$ is the same as $`p_3(n+3)`$. Let $`y_1=x_1x_20,y_2=x_2x_30,y_3=x_30`$ then this is also the number of solutions of the equation $$y_1+2y_2+3y_3=n$$ with the condition that $`y_i0`$. Thus the number $`p_3(n+3)`$ is the coefficient of $`x^n`$ in the expansion of the function (compare ()) $$(1x)^1(1x^2)^1(1x^3)^1=\underset{n=0}{\overset{\mathrm{}}{}}p_3(n+3)x^n.$$ We have the factorization $$(1x^3)=(1x)(1\theta x)(1\theta ^2x)$$ where $`\theta `$ is a $`3`$-rd root of unity, hence $`(1x)^1(1x^2)^1(1x^3)^1`$ $`=(1x)^3(1+x)^1(1\theta x)^1(1\theta ^2x)^1`$ $`={\displaystyle \frac{1}{6(1x)^3}}+{\displaystyle \frac{1}{4(1x)^2}}+{\displaystyle \frac{17}{72(1x)}}+{\displaystyle \frac{1}{8(1+x)}}+`$ $`+{\displaystyle \frac{1}{9(1\theta x)}}+{\displaystyle \frac{1}{9(1\theta ^2x)}}`$ and we get from the expansion of each of the term of the partial fraction decomposition that $$p_3(n+3)=\frac{(n+3)^2}{12}\frac{7}{72}+\frac{(1)^n}{8}+\frac{\theta ^n+\theta ^{2n}}{9}.$$ We infer that $$\|p_3(n+3)\frac{(n+3)^2}{12}\|<\frac{1}{2}$$ or equivalently that $$\|p_3(n)\frac{n^2}{12}\|<\frac{1}{2}.$$ The following identity is easily established by induction: Theorem B12 The number $`p_k(n)`$ satisfies the following recursive relation: $`p_k(n)=p_{k1}(n1)+p_k(nk)`$. Obviously we have $`p_1(n)=n`$ and $`p_2(n)=n/2`$ or $`(n1)/2`$ according to $`n`$ being even or odd. Thus Theorem 10 yields $`p_3(n)=p_2(n1)+p_3(n3)`$, $`p_4(n)=p_3(n1)+p_4(n4),`$ $`p_5(n)=p_4(n1)+p_5(n5)`$ and we get for example $$p_1(6)=1,p_2(6)=3,p_6(6)=1$$ $$p_3(6)=p_2(5)+p_3(3)=3,$$ $$p_4(6)=p_3(5)=p_2(4)=2,$$ $$p_5(6)=p_4(5)=p_3(4)=p_2(3)=1$$ hence $$p(6)=\underset{k=1}{\overset{6}{}}p_k(6)=1+3+3+2+1+1=11.$$ The total partition length $`L(n)`$ of a positive integer $`n`$ is defined to be $`L(n)={\displaystyle \underset{k=1}{\overset{n}{}}}kp_k(n).`$ (57) For example if $`n=6`$ then $`L(6)=1+6+9+8+5+6=35.`$ For $`n=7`$ we have $$p_1(7)=1,p_2(7)=3,p_7(7)=1$$ $$p_3(7)=p_2(6)+p_3(4)=p_2(6)+p_2(3)=4,$$ $$p_4(7)=p_3(6)=3,$$ $$p_5(7)=p_4(6)=2,$$ $$p_6(7)=p_5(6)=1$$ hence $$p(7)=\underset{k=1}{\overset{7}{}}p_k(7)=1+3+4+3+2+1+1=15$$ and the total partition length $$L(7)=1+6+12+12+10+6+7=54.$$ For general $`k`$ one has the following asymptotic formula: Theorem B13 For $`n\mathrm{}`$ the number $`p_k(n)`$ is asymptotically given by: $$p_k(n)\frac{n^{k1}}{(k1)!k!}.$$ Proof. The number $`p_k(n)`$ is defined to be the number of solutions of $`x_1+\mathrm{}+x_k=n`$ with the condition that $`1x_kx_{k1}\mathrm{}x_1`$. If we drop this last condition then the $`k!`$ permutations of a solution is also a solution of $`x_1+\mathrm{}+x_k=n`$. However since $`x_i`$ may equal $`x_j`$ for $`ij`$ hence we have the inequality: $$C_{k1}^{n1}=\frac{(n1)!}{(k1)!(nk)!}k!p_k(n).$$ On the other hand, if we set $`y_i=x_i+(ki)`$ and if $`x_1,\mathrm{},x_k`$ is a solution with $`1x_kx_{k1}\mathrm{}x_1`$ then the $`y_i`$’s are distinct and is a solution of the equation: $$y_1+\mathrm{}+y_k=n+\frac{k(k1)}{2}.$$ From this we obtain a reverse inequality: $$k!p_k(n)C_{k1}^{n+(k(k1)/2)1}=\frac{\{n+(k(k1)/2)1\}!}{(k1)!\{n+k(k1)/2)k\}!}.$$ The Theorem follows immediately from these two estimates. QED References \[A-S\] Azukawa, K. and Suzuki, M.: Some examples of algebraic geometry and hyperbolic manifolds. Rocky Mountain J. 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warning/0003/nucl-th0003014.html	ar5iv	text	# 1 Introduction and Overview ## 1 Introduction and Overview ### 1.1 Introduction The study of highly excited and dense hadronic matter by means of ultra-relativistic nuclear collisions is a relatively novel area of research at the border between nuclear and particle physics. As such it, is in a rapid experimental and theoretical evolution. The primary goal of research in this field is the creation and investigation of elementary (particle) matter under extreme density and temperature conditions. The existence of a novel non-nuclear high temperature phase of elementary matter is an unavoidable consequence of the current knowledge about the strong nuclear interaction, rooted in the theory of strong interactions, the quantum field theory of quarks and gluons called quantum chromodynamics (QCD).<sup>?,?</sup> Discovery and study of quark-gluon plasma (QGP), a ‘deconfined’ state consisting of mobile, color-charged quarks and gluons, is the objective of the relativistic heavy ion experimental research program underway at the Relativistic Heavy Ion Collider (RHIC) at Brookhaven National Laboratory (BNL), New York, and at the Super Proton Synchrotron (SPS) accelerator at the European Organization for Nuclear Research (CERN), Geneva.<sup>?</sup> In a recent half day long workshop, and in an accompanying press release in February 2000 the CERN laboratory has formally announced that it views the collective evidence available today in their seven relativistic nuclear collision experiments as being conclusive: “…A common assessment of the collected data leads us to conclude that we now have compelling evidence that a new state of matter has indeed been created, at energy densities which had never been reached over appreciable volumes in laboratory experiments before and which exceed by more than a factor 20 that of normal nuclear matter. The new state of matter found in heavy ion collisions at the SPS features many of the characteristics of the theoretically predicted quark-gluon plasma. …”.<sup>?</sup> This research program has been developed over the past two decades in order to study the properties of elementary matter at conditions similar to those seen in the very early universe 30$`\mu s`$ after the big-bang, before the temperature decreases to about $`T=150`$ MeV$`1.710^{12}`$ K. It has been many times argued that it is possible to achieve a laboratory recreation of this condition in a small-bang relativistic nuclear collision. The question is at what collision energy the transition to a color deconfined QGP phase first occurs. Early suggestion has been that this could occur at an intrinsic available energy per participating nucleon as low as 4 and 8 times the nucleon mass, corresponding to the range 30 and 120$`A`$ GeV per nucleon beam energy in fixed target experiments.<sup>?</sup> The conditions in the early universe and those created in nuclear collision experiments differ somewhat: whereas the primordial quark-gluon plasma survived for about 30 $`\mu `$s in the big bang, the comparable conditions in nuclear collisions are not expected to last for more than $`10^{22}`$ s due to the rapid explosion of the hot matter “fireball”. Moreover, in the matter created in heavy ion collisions quarks are expected to outnumber antiquarks noticeably due to the baryon content of the colliding nuclei, whereas the net relative excess of the quarks over antiquarks in the universe was less than $`10^9`$. Numerical simulations of QCD suggest that the nature of the transformation between the hadronic and quark-gluon phases can change drastically as the values of the parameters of the theory are varied.<sup>?,?,?</sup> Recent analytical studies of the phase properties of QCD have supported the conclusion that the dependence on the net baryon density (baryochemical potential) is especially interesting.<sup>?,?</sup> Perhaps the most fundamentally important observable in this context is the latent heat associated with the breaking of color bonds among quarks, leading to the deconfinement of quarks. An experimental determination of this quantity and its dependence on beam energy would be of great scientific interest. Strange particle signatures for the formation and evolution of the deconfined quark-gluon phase of elementary matter form a significant cornerstone of experimental QGP discovery. This subject has been developed quite intensely for the past 20 years.<sup>?,?,?,?,?,?</sup> The enabling difference in physics between confined and deconfined matter concerning strange particle signatures is rather simple: $``$ In the QGP phase the particle density is high enough and the strange flavor production energy threshold low enough to assure that a high abundance of strangeness can actually be produced on the time scale available<sup>?,?,?,?,?</sup> while in the confined matter phase this has been shown not to be the case,<sup>?,?</sup> as long as one wants to remain consistent with other experimental results. $``$ Population at, and even in excess, of chemical equilibrium of hadron phase space occupancies occurs only when the entropy rich QGP phase disintegrates rapidly and explosively into hadrons.<sup>?,?,?,?</sup> There are several important and when viewed together, uniquely QGP characteristic predictions regarding strangeness, expected to occur should deconfinement set in. Specifically, the three pillars on which the QGP hypothesis stands when seen by means of strangeness flavor observables are: 1) matter-antimatter symmetry as seen for directly emitted strange baryon and antibaryon particles in the $`m_{}`$-spectral shape and strange quark fugacity; 2) (multi)strange baryon and antibaryon enhancement increasing with strangeness content; 3) enhancement of the (specific) strangeness flavor yield per reaction participant (baryon), by a factor 1.5–3 at SPS conditions, the value depending on what is used as baseline, and if one looks alone at the central rapidity region, where this effect is strongest, or considers the global strangeness yield, including the kinematic domains of projectile and target fragmentation. All three predictions have recently been confirmed at the current SPS energy range 158 GeV per nucleon for Lead (Pb), and some also for 200 GeV per nucleon Sulphur (S) induced reactions. Examples and stepping stones are in particular: 1) The WA97 collaboration reported a detailed study of transverse mass strange baryon and antibaryon spectra which show a highly unusual symmetry between strange baryon and antibaryon sector.<sup>?</sup> 2) A detailed analysis of Pb–Pb results by the WA97 collaboration has demonstrated, comparing p–p, p–A with A–A results, a strong enhancement in the pertinent (multi) strange baryon and antibaryon yields, increasing with strangeness content.<sup>?,?,?,?</sup> The results of the NA49 collaboration are consistent with these findings.<sup>?</sup> The WA85 collaboration also finds an enhancement of multi-strange baryons and antibaryons, increasing with strangeness content in S-induced reactions.<sup>?</sup> 3a) Strangeness enhancement at mid-rapidity has been observed in S-induced reactions by the experiments NA35,<sup>?</sup> WA85 and WA94,<sup>?</sup> and NA44.<sup>?</sup> In the larger Pb-Pb reaction system strangeness enhancements are reported by WA97,<sup>?,?,?</sup> NA49,<sup>?</sup> and NA44.<sup>?</sup> Results of the experiment NA52 suggest further that the onset of strangeness enhancement occurs rather suddenly as the centrality of the collisions and thus the size of participating matter rises above baryon number $`B`$ = 40–50 .<sup>?</sup> 3b) Global Strangeness enhancement has been observed both in S-induced and in Pb-Pb reactions by the NA35,<sup>?</sup> and NA49,<sup>?</sup> experiments. In this article, we also rely in many aspects of this discussion indirectly on other experimental results of NA35 and NA49 collaborations, which offer a global view on particle production pattern considering the large kinematic acceptance.<sup>?,?</sup> We will not discuss or use here other experimental discoveries which have contributed to the CERN announcement, which are not related to strangeness, such as $`J/\mathrm{\Psi }`$ suppression, dilepton and direct photon production. Readers interested in these topics should consult the list of experimental results.<sup>?</sup> In view of many often intricate but, when analyzed, convincing experimental findings about strange particle production, the purpose of this article is to present a comprehensive and selfconsistent view on the understanding of the evidence comprised in strange hadron production for the formation of quark-gluon plasma at CERN, and to discuss resulting expectations how this observable will perform at RHIC. We address in this review: iii) in section 2 the status of a analysis of the experimental data obtained at SPS,<sup>?</sup> iii) in section 3 the implications of these results for the understanding of the dense iii) phase formed in these reactions, iii) in section 4 an adaptation of the dynamical theory of strangeness production iii) in QGP to RHIC conditions, iv) in section 5 an application of these results to obtain predictions for hyperon iii) yields from QGP at RHIC, iv) in section 6 highlights of the results presented here and we draw our conclusions. ### 1.2 Overview We introduce in section 2, the Fermi-2000 model,<sup>?,?</sup> a straightforward elaboration of the original Fermi proposal,<sup>?</sup> that final state strongly interacting particles are produced with a probability commensurate to the size of the accessible phase space. In this approach, the hadron phase space is characterized to the required degree of accuracy by six parameters which have a clear physical meaning and can be in computed, and/or rather easily understood qualitatively, when their values have been determined by an analysis of the experimental data. The physical picture underlying the use of the statistical Fermi model in the 21st century is the sudden, explosive disintegration of a high temperature hadronic matter fireball, apparently consisting of deconfined quark-gluon matter. Since its proposal 50 years ago, the Fermi model has been subject to considerable scrutiny and adaptation, with Hagedorn’s ‘boiling’ hadronic matter being the most important stepping stone.<sup>?</sup> The following were the relevant recent steps in the development of the statistical particle production description required to analyze the strange particle experimental results: 1. Considering the interest and considerable theoretical effort vested in understanding strange quark production mechanisms and the study of chemical equilibration processes,<sup>?,?</sup> it was a natural refinement to introduce an expression of chemical non-equilibrium in the number of strange quark pairs, $`\gamma _s1`$,<sup>?</sup> into the analysis of strange hadrons. 2. In the analysis of the entropy content in S–W/Pb interactions,<sup>?</sup> we found entropy excess related to excess of meson abundance. This prompted us to explore possible nonequilibrium yield of mesons compared to that of baryons.<sup>?</sup> 3. Since the dense hadron fireball is subject to explosive disintegration, final state hadrons emerge from rapidly outward moving volume cells. In order to describe quantitatively spectra of hadronic particles, and their yields in restricted domains of phase space, such ‘collective’ matter flow motion needs to be modeled.<sup>?,?,?</sup> 4. Chemical (particle abundance changing) processes occurring at the time of hadronization do not generally lead to an equilibrium chemical yield of light quarks, and the chemical non-equilibrium is more pronounced if hadronization is a sudden process on the time scale of chemical quark equilibration. While this effect has been well accepted for strange quarks, as these need to be produced in microscopic processes, the need to introduce the light quark pair abundance parameter, $`\gamma _q1`$, was recognized rather late,<sup>?,?</sup> and is not yet widely accepted. A data analysis we perform, allowing for all these effects, does not simply yield a set of ‘best’ parameters; rather: iii) it offers a complete characterization of the phase space of hadrons and its occupancy, allowing one to extrapolate reliably the particle yields into kinematic domains not accessible at present; iii) it allows one to study and understand the magnitude of model parameters, so that we can safely to extrapolate their values to other reaction conditions, e.g., from SPS to RHIC as will be done here; iii) it allows one to evaluate the physical properties of the phase space characterized by these parameters, which provides very precise information about the physical properties of the particle source. This in turn leads to the understanding of the nature of the dense fireball created in the heavy ion collision. In order to pursue these aims, we need to reach considerable precision in the description of experimental results. Our approach and objectives elaborate significantly on the now commonly accepted observation that all hadronic particles produced in strong interaction processesqualitatively satisfy statistical model predictions, as discovered and discussed in great detail by Rolf Hagedorn more than 30 years ago.<sup>?</sup> The fact that the statistical model indeed ‘works’ does not cease to amaze and impress.<sup>?,?</sup> At times this even provokes the hypothesis that ‘thermal’ abundances of hadrons could arise in some mysterious and unknown way,<sup>?</sup> and thus one could proceed to predict ‘thermal’ yields of hadrons, apparently believing that the statistical ‘thermal’ model substitutes for conventional particle production mechanisms. The recent proposal by Bialas to consider the fluctuation of string tension is indeed suggesting how such an explanation could arise in p–p reactions.<sup>?</sup> However, for reactions of large nuclei, the statistical description implies and exploits the result of repeated occurrence of microscopic collision reactions, and the associated approach to an equilibrium distribution shape, and independently, also approach to chemical particle abundance equilibrium. This is independent of the above described possibility that in elementary reaction systems equilibration is possibly a consequence of microscopic properties of strong interactions. Our discussion of A–A reactions thus aims at an improvement of the statistical description beyond the ‘thermal’ model, such that we can deal with standard deviation errors comparing theory and experiment, as is common in the field of particle and nuclear physics. The way we set up the Fermi-2000 model represents the microscopic processes that are occurring, and of course there are limits to this description. For example, primary high energy initial interactions can produce heavy quark flavor which could not arise from the generally softer interactions occurring in the kinetically equilibrated system (an example is expected production of charm at RHIC<sup>?,?</sup>). More generally, an acceptable failure of a statistical description is the one which under-predicts the yield of rarely produced particles. For this reason, it is necessary to scrutinize the validity of the statistical description, at least for the most rarely produced hadrons. Pertinent results of a complete analysis of the Pb–Pb system are presented in section 3. We describe abundances and spectra of hadronic particles observed by both the wide acceptance NA49-experiment and the central rapidity (multi)strange (anti)baryon WA97-experiment. Our method of analysis shows that both these families of results obtained with widely different methods are consistent and it allows us to reach the required precision in their description. This objective could be reached only after we have introduced light quark chemical non-equilibrium and allowed that the all strange $`\mathrm{\Omega }(sss)`$ and $`\overline{\mathrm{\Omega }}(\overline{s}\overline{s}\overline{s})`$ hadrons are enhanced beyond their statistical phase space yield, a point we will discuss in greater detail in subsection 3.4 below. These developments occurred after the last comprehensive review of the subject appeared,<sup>?</sup> and after the extensive SPS-Pb-beam experimental results became available. The introduction of light quark chemical nonequilibrium has had a significant impact on the determination of the physical properties of the hadron emitting source, as it allows for a considerable reduction of the chemical freeze-out temperature: specifically we have determined that $`T_f=143\pm 5`$ MeV, we have also included here an estimate of the systematic error, for the temperature at which practically all strange hadrons are formed. This result is nearly 30 MeV below values that one might infer otherwise in a qualitative study of the statistical hadron yields.<sup>?,?,?</sup> This relatively low freeze-out temperature is consistent with the result that the chemical freeze-out parameters determine correctly the shape of hadron $`m_{}`$-spectra, which suggests that after the deconfined QGP source has dissociated the resulting hadrons are practically free-streaming and thus that thermal and chemical freeze-out do not differ much if at all. This particle production scheme is called sudden hadronization.<sup>?,?</sup> In terms of a microscopic model it occurs when hadronic particles are produced either in: a) an evaporation process from a hot expanding surface, or b) a sudden global hadronization process of exploding, possibly super-cooled deconfined matter. Experimental evidence supporting the picture of sudden QGP hadronization is most directly derived from the baryon-antibaryon transverse mass $`m_{}`$-spectral symmetry. Another piece of evidence for a sudden hadronization is the chemical overabundance of light quark pairs in hadronization and the associated maximization of entropy density in hadron phase space as will be discussed in subsection 2.3. Not to be forgotten is the original observation about experimental data that has led to the data interpretation in terms of the sudden hadronization model<sup>?</sup>: the strange quark fugacity as measured by emitted strange hadrons implies a source with freely moving strange $`s`$ and antistrange $`\overline{s}`$ quarks such that $`s\overline{s}=0`$, a point we will discuss in subsection 2.2. In subsection 3.3, we study in depth the phenomenon of strangeness enhancement and show that the rather precise analysis results we obtain are in excellent agreement with the theoretically computed strangeness yields, assuming formation of the QGP phase. We also show that, in the case of Pb–Pb, the explosive disintegration of the dense QGP fireball leads to an overpopulation of the strangeness phase space abundance, and show that theoretical results are again in good agreement with the results of the data analysis. The initial temperatures for this agreement to occur are in the range $`260<T_{\mathrm{ch}}<320`$, which values we obtain from models of collision dynamics. In section 4, we develop the theoretical method which leads to the finding of overpopulated strangeness phase space discussed in subsection 3.3. This, at a first sight surprising result, occurs due to early freeze-out of strangeness abundance in a rapid explosive evolution of the QGP fireball. We find a similar non-equilibrium result, in section 5, for RHIC condition in presence of transverse expansion which increases the speed at which QGP dilutes. We explore the dynamics of the phase space occupancy, rather than particle density, which allows us to eliminate much of the dependence on the dynamical flow effects by incorporating in the dynamics considered the hypothesis of entropy conserving matter flow and evolution. Exploiting the experience with SPS data analysis, we are able to consider, in section 5, strange particle production at RHIC. We obtain an unexpected particle abundance pattern: during the hadronization of the baryon-poor RHIC-QGP phase there is considerable advantage for strangeness flavor to stick to baryons and antibaryons. This can be easily understood realizing that production of strange (anti)baryons is favored over production of kaons by the energy balance, i.e.: $`E(\mathrm{\Lambda }+\pi )<E(\text{N+K})`$. Moreover, at RHIC there are a high number of strange quarks per baryon available, and in this strangeness bath just a few (anti)baryons will manage to emerge without strangeness content. We thus expect and find in detailed study in section 5 that hyperon production dominates baryon production, i.e., most baryons and antibaryons produced will be strange. We consider this result to be a unique consequence of the sudden QGP hadronization scenario observed at SPS and hope and expect that hyperon dominance should, when observed at RHIC, be generally accepted as proof of formation of the deconfined phase of nuclear matter. This phenomenon shows how much more pronounced will be the physics of strangeness in QGP at RHIC, as compared to the ten times lower SPS energy range. ## 2 Contemporary Fermi Model of Hadron Production ### 2.1 Phase space and parameters The relative number of final state hadronic particles freezing out from, e.g., a thermal quark-gluon source is obtained noting that the fugacity $`f_i`$ of the $`i`$-th emitted composite hadronic particle containing $`k`$-components is derived from fugacities $`\lambda _k`$ and phase space occupancies $`\gamma _k`$: $$N_ie^{E_i/T_f}f_i=e^{E_i/T_f}\underset{ki}{}\gamma _k\lambda _k.$$ (1) We study chemical properties of light quarks $`u,d`$ jointly, denoting these by a single index $`q`$ and also consider chemical properties of strange quarks $`s`$ . Thus as seen in Eq. (1), we study particle production in terms of five statistical parameters. In addition there is at least one matter flow velocity parameter. The six parameters which characterize the accessible phase space of hadronic particles made of light quarks ‘q’ and strange quarks ‘s’ and their natural values, assuming a QGP source, are in turn: 1) $`\lambda _s`$: The value of strange quark fugacity $`\lambda _s`$ cab be obtained from the requirement that strangeness balances, $`N_sN_{\overline{s}}=0,`$ which for a source in which all $`s,\overline{s}`$ quarks are unbound and thus have symmetric phase space, implies $`\lambda _s=1`$ . However, the Coulomb distortion of the strange quark phase space plays an important role in the understanding of this constraint for Pb–Pb collisions,<sup>?</sup> leading to the Coulomb-deformed value $`\lambda _s1.1`$ , as discussed in next subsection. 2) $`\gamma _s`$: The strange quark phase space occupancy $`\gamma _s`$ can be computed, and will be studied in this review in detail within the framework of kinetic theory.<sup>?,?</sup> For a rapidly expanding system the production processes will lead to an over-saturated phase space with $`\gamma _s>1`$ . The difference between the two different types of chemical parameters $`\lambda _i`$ and $`\gamma _i`$ is that the phase space occupancy factor $`\gamma _i`$ regulates the number of pairs of flavor ‘$`i`$’, and hence applies in the same manner to particles and antiparticles, while fugacity $`\lambda _i`$ applies only to particles, while $`\lambda _i^1`$ is the antiparticle fugacity. 3) $`\lambda _q`$: The light quark fugacity $`\lambda _q`$, or equivalently, the baryochemical potential: $`\mu _B=3T_f\mathrm{ln}\lambda _q,`$ (2) regulate the baryon density of the fireball and hadron freeze out. This density can vary dependent on the energy and size of colliding nuclei, and thus the value of $`\lambda _q`$ is not easily predicted. However, since we know the energy per baryon content in the incoming nuclei, if we assume that the deposition of the baryon number and energy (‘stopping’) is similar, we know the energy per baryon content in the fireball.<sup>?</sup> This qualitative knowledge can be used in a study of equations of state applicable to the dense fireball to establish a constraint. 4) $`\gamma _q`$: The equilibrium phase space occupancy of light quarks $`\gamma _q`$ is expected to significantly exceed unity to accommodate the excess entropy content in the plasma phase.<sup>?,?</sup> There is an upper limit: $$\gamma _q<\gamma _q^ce^{m_\pi /2T},$$ (3) which arises if all pions produced are simultaneously present, forming a Bose gas. We will address this effect in subsection 2.3. 5) $`T_f`$: The freeze-out temperature $`T_f`$ is expected to be not much different from the Hagedorn temperature $`T_H160`$ MeV,<sup>?</sup> which characterized particle production in proton-proton reactions. 6) $`v_c`$: The collective expansion velocity $`v_c`$ is expected to remain below the relativistic sound velocity<sup>?</sup>: $$v_c1/\sqrt{3}.$$ (4) When the source emitting the free streaming particles is undergoing local collective flow motion, spectra of particles emitted are described by replacing the Boltzmann factor in Eq. (1) by: $`e^{E_i/T}`$ $``$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle 𝑑\mathrm{\Omega }_v\gamma _c(1\stackrel{}{v}_\mathrm{c}\stackrel{}{p}_i/E_i)e^{\frac{\gamma _cE_i}{T}\left(1\stackrel{}{v}_\mathrm{c}\stackrel{}{p}_i/E_i\right)}},`$ $`\gamma _c`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1\stackrel{}{v}_\mathrm{c}^{\mathrm{\hspace{0.17em}2}}}}},`$ (5) a result which can be intuitively obtained by a Lorentz transformation between an observer on the surface of the fireball, and one at rest in the laboratory frame. Formal derivation of this and more elaborated results requires a considerably more precise framework.<sup>?</sup> The resulting yields of final state hadronic particles are most conveniently characterized taking the Laplace transform of the accessible phase space. This approach generates a function which in its mathematical properties is identical to the partition function. For example for the open strangeness sector we find, for the case $`v_\mathrm{c}=0`$: $`\mathrm{ln}𝒵_s={\displaystyle \frac{VT^3}{2\pi ^2}}`$ $`\{(\lambda _s\lambda _q^1+\lambda _s^1\lambda _q)\gamma _s\gamma _qF_K+(\lambda _s\lambda _q^2+\lambda _s^1\lambda _q^2)\gamma _s\gamma _q^2F_Y`$ (6) $`+(\lambda _s^2\lambda _q+\lambda _s^2\lambda _q^1)\gamma _s^2\gamma _qF_\mathrm{\Xi }+(\lambda _s^3+\lambda _s^3)\gamma _s^3F_\mathrm{\Omega }\}.`$ The integrated momentum phase space factors $`F_i`$ for kaons $`i=K`$, single strange hyperons $`i=Y`$, doubly strange cascades $`i=\mathrm{\Xi }`$ and triply strange omegas $`i=\mathrm{\Omega }`$ are: $$F_i=\underset{j}{}g_{i_j}W(m_{i_j}/T),W(x)=x^2K_2(x).$$ (7) $`g_{i_j}`$ is the statistical degeneracy of each contributing hadron resonance ‘$`j`$’ of the kind ‘$`i`$’ in the $`_j`$, which comprise all known strange hadron resonances. $`K_2`$ is the modified Bessel function which arises from the relativistic phase space integral of the thermal particle distribution $`f(\stackrel{}{p})e^{\sqrt{m^2+p^2}/T}`$. It is important to keep in mind that: a) Eq. (6) does not require formation of a phase comprising a gas of hadrons, but is not inconsistent with such a step in evolution of the fireball; in that sense it is not a partition function, but just a look-alike object arising from the Laplace transform of the accessible phase space, and b) the final particle abundances measured in an experiment are obtained after all unstable hadronic resonances ‘$`j`$’ are allowed to disintegrate, contributing to the yields of stable hadrons; c) the unnormalized particle multiplicities arising are obtained differentiating Eq. (6) with respect to particle fugacity. The relative particle yields are simply given by ratios of corresponding chemical factors, weighted with the size of the momentum phase space accepted by the experiment. For particles showing the same spectral shape comparison of normalization of $`m_{}`$ spectra suffices, e.g., Ref. <sup>?</sup>: $$\frac{\mathrm{\Xi }^{}(dss)}{\mathrm{\Lambda }(dds)}\|_m_{}=\frac{g_\mathrm{\Xi }\gamma _d\gamma _s^2\lambda _d\lambda _s^2}{g_\mathrm{\Lambda }\gamma _d^2\gamma _s\lambda _d^2\lambda _s}.$$ (8) $`g_i`$ are the spin statistical factors of the states considered. Similarly: $$\frac{\overline{\mathrm{\Xi }^{}(dss)}}{\overline{\mathrm{\Lambda }(dds)}}\|_m_{}=\frac{g_\mathrm{\Xi }\gamma _d\gamma _s^2\lambda _d^1\lambda _s^2}{g_\mathrm{\Lambda }\gamma _d^2\gamma _s\lambda _d^2\lambda _s^1}.$$ (9) When acceptance is limited to central rapidity, and significant flow is present considerable effort must be made to introduce appropriate phase space weights. d) In some experimental data it is important to distinguish the two light quark flavors as is in fact the case in the two above examples. This can be incorporated considering how the average light quark fugacity varies between both light quark species,<sup>?</sup> and assuming that the phase space occupancies are equal. We consider, for SPS energy range, the radial flow model, which is without doubt the simplest of the reasonable and expected matter flow cases possible, in view of the behavior of global observables seen in these experiments. As the results below show, this suffices to assess the impact of collective flow on the data analysis originally developed to be as little as possible sensitive to collective matter flow, even when particle yields in highly restricted regions of $`m_{},y`$ are considered. The collective source flow can completely change the shape of momentum distribution of particle produced, though of course it leaves unchanged the total particle yield, which is the integral sum of particle multiplicity over the entire phase space of the flow spectrum. However, particles of different mass experience differing flow effects when $`m_{},y`$ acceptance cuts are present. Moreover, particles can freeze-out at slightly different conditions. In order to limit the influence of the practically unknown collective flow structure on particle yields in limited domains of the accessible phase space, we study compatible particle ratios: these are yield ratios obtained in a restricted domain of $`m_{},y`$, for particles of similar mass and believed to have a similar interaction strength with the matter background. We now will address in turn two special topics, which slightly contradict expectations, and thus require more attention. Firstly, we review the the properties of the strange quark fugacity $`\lambda _s`$ , which is sensitive to the possible asymmetry between strange and antistrange quarks in the source. The importance of this parameter is that it potentially helps distinguish the confined from deconfined phase: while in the baryon-rich confined phase the requirement of strangeness conservation implies that $`\lambda _s>1`$ , in the deconfined phase the symmetry between phase space of strange and antistrange quarks implies $`\lambda _s1`$ . Following this, we address in more detail case of pions,<sup>?</sup> which is exceptional since we will be considering rather large values of $`\gamma _q>1.5`$ . As we shall see, the pion gas emerging from the QGP phase is strongly influenced by Bose correlation effects; in fact it is close to satisfying the Bose condensation condition. ### 2.2 Coulomb force It has been recognized for a long time that the Coulomb force can be of considerable importance in the study of relativistic heavy ion collisions. It plays an important role in the HBT analysis of the structure of the particle source.<sup>?</sup> We show that the analysis of chemical properties at freeze-out is also subject to this perturbing force, and in consideration of the precision reached in the study of particle ratios, one has to keep this effect in mind. We consider a Fermi gas of strange and antistrange quarks, allowing that the Coulomb potential $`V`$ established by the excess charge of the colliding nuclei distorts significantly the phase space. Within a relativistic Thomas-Fermi phase space occupancy model,<sup>?</sup> and allowing for finite temperature in QGP we have<sup>?</sup>: $$N_sN_{\overline{s}}=\underset{R_\mathrm{f}}{}g_s\frac{d^3rd^3p}{(2\pi )^3}\left[\frac{1}{1+\gamma _s^1\lambda _s^1e^{(E(p)\frac{1}{3}V(r))/T}}\frac{1}{1+\gamma _s^1\lambda _se^{(E(p)+\frac{1}{3}V(r))/T}}\right],$$ (10) which clearly cannot vanish for $`V0`$ in the limit $`\lambda _s1`$. In Eq. (10) the subscript $`R_\mathrm{f}`$ on the spatial integral reminds us that only the classically allowed region within the fireball is covered in the integration over the level density; $`E=\sqrt{m^2+\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}}}`$, and for a uniform charge distribution within a radius $`R_\mathrm{f}`$ of charge $`Z_\mathrm{f}`$: $$V=\{\begin{array}{cc}\frac{3}{2}\frac{Z_\mathrm{f}e^2}{R_\mathrm{f}}\left[1\frac{1}{3}\left(\frac{r}{R_\mathrm{f}}\right)^2\right],\hfill & \text{for}r<R_\mathrm{f};\hfill \\ \frac{Z_\mathrm{f}e^2}{r},\hfill & \text{for}r>R_\mathrm{f}.\hfill \end{array}$$ (11) One obtains a rather precise result for the range of parameters of interest to us using the Boltzmann approximation: $$N_sN_{\overline{s}}=\gamma _s\left\{g_s\frac{d^3p}{(2\pi )^3}e^{E/T}\right\}_{R_\mathrm{f}}d^3r\left[\lambda _se^{\frac{V}{3T}}\lambda _s^1e^{\frac{V}{3T}}\right].$$ (12) The Boltzmann limit allows us also to verify and confirm the signs: the Coulomb potential is negative for the negatively charged $`s`$-quarks with the magnitude of the charge, $`1/3`$, made explicit in the potential terms in all expressions above. We thus have<sup>?</sup>: $$\stackrel{~}{\lambda }_s\lambda _s\lambda _\mathrm{Q}^{1/3}=1,\lambda _\mathrm{Q}\frac{_{R_\mathrm{f}}d^3re^{\frac{V}{T}}}{_{R_\mathrm{f}}d^3r}.$$ (13) $`\lambda _\mathrm{Q}<1`$ expresses the Coulomb deformation of strange quark phase space. $`\lambda _\mathrm{Q}`$ is not a fugacity that can be adjusted to satisfy a chemical condition, since consideration of $`\lambda _i,i=u,d,s`$, exhausts all available chemical balance conditions for the abundances of hadronic particles. The subscript $`R_f`$, in Eq. (13), reminds us that the classically allowed region within the dense matter fireball is included in the integration over the level density. Choosing $`R_\mathrm{f}=8`$ fm, $`T=140`$ MeV, $`m_s=200`$ MeV, noting that the value of $`\gamma _s`$ is practically irrelevant as this factor cancels in Boltzmann approximation, see Eq. (12), we find for $`Z_\mathrm{f}=150`$ that the value $`\lambda _s=1.10`$ corresponds to $`R_\mathrm{f}=7.9`$ fm. The Coulomb effect is thus relevant in central Pb–Pb interactions, while for S–Au/W/Pb reactions, similar analysis leads to a value $`\lambda _s=1.01`$, little different from the value $`\lambda _s=1`$ expected in the absence of the Coulomb phase space deformation. Another way to understand the varying importance of the Coulomb effect is to note that while the Coulomb potential acquires in the Pb–Pb case a magnitude comparable to the quark chemical potential, it remains small on this scale for S–Au/W/Pb reactions. ### 2.3 Super-dense pion gas and chemical non-equilibrium For pions composed of a light quark-antiquark pair, the chemical fugacity is $`\gamma _q^2`$, see Eq. (1). Thus the pion momentum space distribution has the Bose shape: $$f_\pi (E)=\frac{1}{\gamma _q^2e^{E_\pi /T}1},E_\pi =\sqrt{m_\pi ^2+p^2}.$$ (14) The range of values for $`\gamma _q`$ is bounded from above by the Bose singularity. When $`\gamma _q\gamma _q^c`$, Eq. (3), the lowest energy state (in the continuum limit with $`p0`$) will acquire macroscopic occupation and a pion condensate is formed. Formation of such a condensate ‘consumes’ energy without consuming entropy of the primordial high entropy QGP phase. On the other hand, as we shall see presently, when $`\gamma _q\gamma _q^c`$ the entropy content of the pion gas initially grows! Thus while the development, directly from the QGP phase, of a pion condensate is not likely, the sudden hadronization of entropy rich QGP should lead to the limiting value $`\gamma _q\gamma _q^c`$, in order to more efficiently connect the entropy rich deconfined and the confined phases. An interesting feature of such a mechanism of phase transition is that the chemical non-equilibrium reduces and potentially eliminates any discontinuity in the phase transition, which thus, in the experiment, will appear more like a phase transformation without critical fluctuations, even if theory implies a 1st order phase transition for statistical equilibrium system. In Fig. 2.3, we show the physical properties of a pion gas as function of $`\gamma _q`$ for a gas temperature $`T=142`$ MeV.<sup>?</sup> We see (solid line) that a large range of entropy density can be accommodated by varying the parameter $`\gamma _q`$. It is important to remember that in the hadronization of a quark-gluon phase it is relatively easy to accommodate energy density, simply by producing a few heavy hadrons. However, such particles being in fact non-relativistic at the temperature considered, are not effective carriers of pressure and entropy. However, as we see now in Fig. 2.3, the super-dense pion gas is just the missing element to allow a rapid hadronization process, since the entropy density is nearly twice as high at $`\gamma _q\gamma _q^c`$ than at $`\gamma _q=1`$. Without this phenomenon one has to introduce a mechanism that allows the parameter $`VT^3`$ to grow, thus expanding either the volume $`V`$ due to formation of so called mixed phase or invoking rise of $`T`$ in so called reheating. The specific properties of the super-dense pion gas are shown Fig. 2.3. In Fig. 2.3a, we relate the properties to the chemical equilibrium value $`\gamma _q=1`$ and we also show that the Boltzmann approximation is not qualitatively wrong, as long as $`\gamma _q<\gamma _q^c`$ . In Fig. 2.3b, we see the relative change in energy per pion, (inverse of) entropy per pion, and energy per entropy. Interestingly, we note that the entropy per pion drops as $`\gamma _q`$ increases, and at the condensation point $`\gamma _q=\gamma _q^c`$, we can add pions without increase in entropy. We further note that a hadronizing gas will consume, at higher $`\gamma _q`$, less energy per particle, and that the energy per entropy is nearly constant. It is important to remember that if the hadronization process were adiabatic, allowing a full equilibrium relaxation, naturally $`\gamma _i1`$ would arise: as is implicitly well known the value $`\gamma _i1`$ maximizes the entropy for a particle gas at fixed total energy, corresponding to the chemical equilibrium.<sup>?</sup> This result is easily found considering Boltzmann pion gas, and recalling in some detail the definition of entropy, since the standard equilibrium expressions do not apply: $`S_{\mathrm{B},\mathrm{F}}=`$ $`{\displaystyle \frac{d^3pd^3x}{(2\pi \mathrm{})^3}\left[\pm (1\pm f(x,p))\mathrm{ln}(1\pm f(x,p))f(x,p)\mathrm{ln}f(x,p)\right]},`$ (15) $``$ $`V{\displaystyle \frac{d^3p}{(2\pi \mathrm{})^3}f(p)\mathrm{ln}[e/f(p)]},`$ (16) where aside of ‘B,F’ (Bose, Fermi) also the Boltzmann limit for a homogeneous spatial distribution is shown explicitly. Evaluating in Boltzmann limit the particle number and energy, we find that the factor $`\gamma _q^2`$ becomes a normalization factor which describes the average occupancy of the phase space relative to the equilibrium value, and for entropy, we also find a logarithmic term: $`N=`$ $`\gamma _q^2N\|_{\mathrm{eq}}aV\gamma _q^2T^3,`$ (17) $`E=`$ $`\gamma _q^2E\|_{\mathrm{eq}}3aV\gamma _q^2T^4,`$ (18) $`S=`$ $`\gamma _q^2S\|_{\mathrm{eq}}+\mathrm{ln}\left(\gamma _q^2\right)\gamma _q^2N\|_{\mathrm{eq}}4aV\gamma _q^2T^3+\mathrm{ln}\left(\gamma _q^2\right)aV\gamma _q^2T^3.`$ (19) For massless pions, $`a=g/\pi ^2`$, with pion degeneracy $`g=3`$ . Setting $`E=`$ Const., we eliminate $`T`$ and find that the entropy as function of $`\gamma _q`$ varies according to: $$S\|_{E=\text{Const.}}\gamma _q^{\frac{1}{2}}(4\mathrm{ln}\gamma _q^2),$$ (20) which has a very weak maximum at $`\gamma _q=1`$; note that at $`\gamma =1.4`$, the entropy is at 98.3% of the value at $`\gamma =1`$. At this point it is important to realize that the chemical equilibrium is much better defined for the more familiar case of a fixed temperature bath (and not an isolated fixed energy fireball discussed above): consider the free energy $``$ of the (non-interacting) relativistic gas at fixed temperature $`T`$. Since $`=ETS`$, we combine Eqs. (1719) and obtain in the Boltzmann limit, $$^l=aVT^4\gamma _q^2\left[1+\mathrm{ln}\left(\gamma _q^2\right)\right],$$ (21) which has a minimum for chemical equilibrium value $`\gamma _q=1`$. However, one now finds that a change by factor 1.4 in $`\gamma _q`$, at fixed $`T`$ leads to a change by 35% in the value of the free energy and even a greater change in entropy. We conclude that, in adiabatic condition, the fireball would evolve to the maximum entropy equilibrium case $`\gamma _i=1,i=q,s`$, the gain in entropy for an isolated system is in this limit very minimal. Thus for a system with rapidly evolving volume, we will in general find more effective paths to increase entropy, other than the establishment of the absolute chemical equilibrium. Hence the values we report $`\gamma _i1,i=q,s`$, are consistent with the present day understanding of explosive evolution of the hadronic matter fireball. In a systematic study of the relevance of different physical parameters, the chemical non-equilibrium at hadron freeze-out has been shown to be a required ingredient in order to arrive at a precise interpretation of the experimental results on particle ratios $`R^j`$ obtained at CERN. This is best seen considering the results for the statistical parameters obtained for the S–Au/W/Pb collisions,<sup>?</sup> and the associated total statistical error, $$\chi _\mathrm{T}^2\frac{_j(R_{\mathrm{th}}^jR_{\mathrm{exp}}^j)^2}{(\mathrm{\Delta }R_{\mathrm{exp}}^j)^2},$$ (22) which are presented in table 2.3. We clearly see the gain in physical significance that is accomplished as chemical non-equilibrium is allowed for by releasing the fixed value $`\gamma _i=1`$, first for strange and next, light quarks. We also observe that allowing for $`\lambda _s1`$ does not lead to an improvement in statistical significance, since the data is compatible with this value expected for the deconfined QGP. Similar systematic study has also been completed for the Pb–Pb system,<sup>?</sup> reconfirming the need to use $`\gamma _i1`$ in the data analysis. The errors in the results shown in table 2.3, and in results that follow below, are one standard deviation errors arising from the propagation of the experimental measurement error. However, these errors are meaningful only when the theoretical model describes the data well, as is the case for last entry line in table 2.3 when we allow light quark chemical nonequilibrium, $`\gamma _q1`$ . ## 3 Strange Hadron Data Analysis ### 3.1 Particle yields The available compatible particle yield ratios (excluding presently $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$, see subsection 3.4) are listed in table 3.1, top section from the experiment WA97, for $`p_{}>0.7`$ GeV, within a narrow $`\mathrm{\Delta }y=0.5`$ central rapidity window. Further below are shown results from the large acceptance experiment NA49, extrapolated by the collaboration to full $`4\pi `$ phase space coverage. We first fit 11 experimental results shown in table 3.1, and then turn to include also the $`m_{}`$-slope into this consideration, and thus have 12 data points. The total error $`\chi _\mathrm{T}^2`$ for the four result columns is shown at the bottom of this table along with the number of data points ‘$`N`$’, parameters ‘$`p`$’ used and (algebraic) redundancies ‘$`r`$’ connecting the experimental results. For $`r0`$ it is more appropriate to quote the total $`\chi _\mathrm{T}^2`$, with a initial qualitative statistical relevance condition $`\chi _\mathrm{T}^2/(Np)<1`$. The first theoretical columns refer to results without collective velocity $`v_c`$ (subscript $`0`$) the three other were allowing for $`v_c`$ (subscript $`v_c`$). In column three, superscript ‘sb’ means that $`\lambda _s`$ is fixed by strangeness balance and, in column four, superscript ‘sc’ means that $`\gamma _q=\gamma _q^c=e^{m_\pi /2T_f}`$, that is $`\gamma _q`$ is fixed by its upper limit, the pion condensation point. All results shown account for slightly higher value of the ratio $`h^{}/B`$ recently reported<sup>?</sup>; $`B`$ is here the number of baryon participants and $`h^{}=\pi ^{}+K^{}+\overline{p}`$ is the yield of stable negative hadrons comprising as indicated pions, kaons and antiprotons. First we note that all columns in table 3.1 represent physically acceptable result for the Pb–Pb collision system: a) presence of collective flow (three last columns) leads to very similar compatible particle ratios, even though improvement of $`\chi _\mathrm{T}`$ occurs when $`v_c0`$ is allowed for; b) the highest confidence result is found just when the light quark phase space occupancy assumes a value at below the pion condensation point; c) strangeness conservation (enforced in second last column) is naturally present, enforcing it does not change in any way the results for particle multiplicities. Allowing radial flow not only improves the capability to describe the data, but it allows us to study $`m_{}`$ particle spectra, which offer another independent measure of flow, and confirm the value of $`v_c`$ — when considering $`v_c`$ along with $`T_{}`$, the inverse slope of the $`m_{}`$ spectra, we have one parameter and several spectral inverse slopes of particles considered. However, we will in the first instance assume that we have just one additional data point and we proceeded as follows: for a given pair of values $`T_f`$ and $`v_\mathrm{c}`$ we evaluate the resulting $`m_{}`$ particle spectrum and analyze it using the spectral shape and kinematic cuts employed by the experimental groups. Once we find values of $`T_\mathrm{f}`$ and $`v_\mathrm{c}`$, we study again the inverse slopes of individual particle spectra and obtain an acceptable agreement with the experimental $`T_{}^j`$ as shown in left section of table 3.1 . We have considered in the same framework the S-induced reactions, and the right section of table 3.1 shows also a good agreement with the WA85 experimental data.<sup>?</sup> We have updated the experimental Pb–Pb results shown in table 3.1 with the current high precision results.<sup>?</sup> However, the theoretical results shown were obtained earlier for slightly different results with larger error bars, and we hope to reevaluate these results in the near future. To model these slopes theoretically, one needs to remember that the vast majority of $`\mathrm{\Lambda }`$ and $`\overline{\mathrm{\Lambda }}`$ is a decay product of $`\mathrm{\Sigma }^0`$ and $`\overline{\mathrm{\Sigma }^0}`$, $`\mathrm{\Lambda }^{}`$ and $`\overline{\mathrm{\Lambda }^{}}`$ and $`\mathrm{\Xi }`$ and $`\overline{\mathrm{\Xi }}`$. Consequently, given the precision of the (inverse) slopes presented, in order to model the $`\mathrm{\Lambda }`$ and $`\overline{\mathrm{\Lambda }}`$ spectra one will need to consider the effect of hadron cascading, which introduces uncertainty arising from a dependence on unmeasured yields. However, given the current availability of quite precise $`\mathrm{\Xi }`$ and $`\overline{\mathrm{\Xi }}`$ slopes, and the fact that these particles are rarely decay products of other hadronic resonances, we will in future use these slopes as the spectral data point input in the data analysis studies. As result, we anticipate a slight reduction in the collective velocity within the errors shown below. ### 3.2 Chemical freeze-out properties The six parameters ($`T_f,v_c,\lambda _q,\lambda _s,\gamma _q,\gamma _s`$) describing the particle abundances are shown in the top section of table 3.2. We also show in the last column the best result for S-induced reactions, where the target has been W/Au/Pb.<sup>?</sup> All results shown in table 3.2 have convincing statistical confidence level. For the S-induced reactions the number of redundancies $`r`$ shown in the heading of the table 3.2 is large, since the same data comprising different kinematic cuts has been included in the analysis. Within error, the freeze-out temperature $`T_\mathrm{f}143\pm 3`$ MeV, seen in table 3.2, is the same for both the S- and Pb-induced reactions, even though the chemical phase space occupancies differ greatly. Such a behavior is expected in view of the similarity of the energy content in the collision in both reaction systems, but greatly differing collision geometry. We find that the variation in the shape of particle $`m_{}`$-spectra is fully explained by a change in the collective velocity, which rises from $`v_c^\mathrm{S}=0.49\pm 0.02`$ to $`v_c^{\mathrm{Pb}}=0.54\pm 0.041/\sqrt{3}=0.577`$. The value of light quark fugacity $`\lambda _q`$ implies that baryochemical potential is $`\mu _B^{\mathrm{Pb}}=203\pm 5>\mu _B^\mathrm{S}=178\pm 5`$ MeV. As in S-induced reactions where $`\lambda _s=1`$, now in Pb-induced reactions, a value $`\lambda _s^{\mathrm{Pb}}1.1`$ characteristic for a source of freely movable strange quarks with balancing strangeness, i.e., $`\stackrel{~}{\lambda }_s=1`$, is obtained, see Eq. (13). Further evidence for low chemical freeze-out temperature is contained in the $`m_{}`$-particle spectra considered in subsection 3.1. Our approach offers a natural understanding of the equality of the $`m_{}`$-slopes of the strange baryons and antibaryons considered which arises because within the sudden hadronization model both these particles emerge free-streaming from QGP. In the hadron based microscopic simulations this behavior of $`m_{}`$-slopes of baryons and antibaryons arises from fine-tuning of the particle-dependent freeze-out times.<sup>?</sup> On the other hand, in such a microscopic study one finds in view of the small reaction cross sections that $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ could freeze out somewhat sooner than the other hadrons, and thus would have a softer spectrum as also confirmed in direct hadronization simulations.<sup>?</sup> We will return to this point just below. The reader should keep in mind that since we find a rather low chemical freeze-out temperature, and can explain the $`m_{}`$ spectra well based on this value, the implied kinetic (collision) freeze-out temperature must be rather similar to the chemical freeze-out. In the bottom section of table 3.2, we also see the energy and entropy content per baryon. The energy per baryon seen in the emitted hadrons is nearly equal to the available specific energy of the collision (8.6 GeV for Pb–Pb, 8.8–9 GeV for S–Au/W/Pb). This implies that the fraction of energy deposited in the central fireball must be nearly (within 10%) the same as the fraction of baryon number. The small reduction of the specific entropy in Pb–Pb compared to the lighter S–Au/W/Pb system maybe driven by the greater baryon stopping in the larger system, also seen in the smaller energy per baryon content. Both collision systems freeze out at the same energy per unit of entropy, $$E/S=0.185\text{GeV}.$$ There is a loose relation of this universality in the chemical freeze-out condition with the suggestion made recently that particle freeze-out occurs at a fixed energy per baryon for all physical systems,<sup>?</sup> considering that the entropy content is related to particle multiplicity. The overall high specific entropy content we find agrees well with the entropy content evaluation we made earlier for the S–Pb case.<sup>?</sup> The high entropy content is observed in the final hadron state in terms of enhanced pion yield. Thus the ratio of $`K^+/\pi ^+`$ is combines these two effects and is not a good indicator of new physics, even though this relatively simple observable continues to attract attention.<sup>?</sup> It would have been more useful if systematic studies of strangeness production and enhancement were to offer as result of their analysis the strangeness yields per participating baryon number. The large values of $`\gamma _q>1`$, seen in table 3.2, imply as discussed earlier that there is phase space over-abundance of light quarks, which receives contribution from, e.g., gluon fragmentation at QGP breakup. $`\gamma _q`$ assumes in the data analysis a value near to where pions could begin to condense,<sup>?</sup> Eq. (3). This result is consistent with the expectations for hadronization of an entropy rich quark gluon plasma, as we discussed above in subsection 2.3. We found by studying the ratio $`h^{}/B`$ separately from other experimental results that the value of $`\gamma _q\gamma _q^c`$ is fixed consistently and independently both, by the negative hadron ($`h^{}`$), and the strange hadron yields. The unphysical range $`\gamma _q>\gamma _q^c1.63`$ can arise (see column Pb$`\|_v^{\text{s}b}`$) since, up to this point, we had used only a first quantum (Bose/Fermi) correction. However, when Bose distribution for pions is implemented, which requires the constraint $`\gamma _q\gamma _q^c`$, we obtain practically the same results, as shown in the third column of table 3.2. ### 3.3 Strangeness enhancement We show, in the bottom section of table 3.2, the specific strangeness content, $`s_f/B`$ along with specific strangeness asymmetry $`(\overline{s}_fs_f)/B`$ seen in the hadronic particles emitted. In the data analysis the requirement that the number of $`s`$ and $`\overline{s}`$ quarks in hadrons is equal is in general not enforced. We see that in lower portion of table 3.2 that this result is found automatically for the symmetric Pb–Pb collision system. However, a 3.5 s.d. effect is seen in the asymmetrical S–Au/W/Pb system Though the errors which we derive from the experimental data are small, there could be in this asymmetric system a considerable systematic experimental error due to data extrapolations made in presence of a significant spectator matter component, coupled with theoretical error from the varying CM-rapidity. On the other hand, the consistency of the Pb–Pb and S–Au/W/Pb results suggest that this asymmetry is possibly a real effect, thus the unseen balance of strangeness could be hidden in a residual (strange) quark matter nugget, which is escaping detection. Such strangeletts could in principle form, since in the hadronization of the S–Au/W/Pb deconfined system the hadron phase space is asymmetric, which leads to strangeness distillation.<sup>?,?,?,?</sup> One of the important quantitative results of this analysis is shown in the bottom section of table 3.2: the high yield of strangeness per baryon, $`s_f/B0.7`$ . We now proceed to verify if this yield is in agreement with the pedictions made over the past 20 years. Perhaps more by chance than design, this analysis result is in agreement with the first calculations of strangeness production employing perturbative QCD,<sup>?</sup> where the value $`N_s/B=n_s/\nu =0.7`$ is reached for the plasma temperature of 300 MeV as shown there in Fig. 3. Since, considerably more refined methods have been developed,<sup>?</sup> and these are in excellent agreement with results of the analysis of experimental results. In view of the high precision reached in this data analysis, we have recomputed the theoretical yield taking for the QCD parameters values generally accepted today: $`\alpha _\mathrm{s}(M_Z)=0.118`$ and $`m_\mathrm{s}(1GeV)=200`$ MeV, correpsonding to $`m_\mathrm{s}(M_Z)=90`$ MeV. In table 5, we summarize for three collision systems we consider S–Au/W/Pb, Ag–Ag, Pb–Pb the key input parameters used in computing the result for $`N_s/B`$ shown below in Fig. 3.3. The first entry line gives the central collision particpant numbers for the three systems considered. Next, in table 5, we see the initial temperature $`T_{\mathrm{ch}}`$ which the evaluation of strangeness production requires as input. $`T_{\mathrm{ch}}`$ is the temperature at the time when light quarks and gluons reach equilibrium. To obtain this value, we compute the collisional pressure and set it equal to thermal pressure at the time the fireball begins to expand.<sup>?,?</sup> To do this we need the (momentum, energy) stopping fractions $`\eta `$ here taken from NA35/NA49 experimental results,<sup>?</sup> (except for interpolation for Ag–Ag, the dotted line in Fig. 3.3). The last line in table 5 addresses the expansion dynamics we use: we employ the observed freeze-out expansion velocity $`v_c`$ as given in the top section of table 3.2. We assume that each local volume expands its size scale $`R`$ at this local velocity, and we consider the process to be entropy conserving, hence we use $`R^3T^3=`$Const. to obtain the time dependence of local fireball temperature. We obbtain the result for $`N_s/B`$ shown in Fig. 3.3, as function of the specific energy available in the fireball $`E/B`$, for the three collision systems S-Au/W/Pb (short-dashed line), Ag–Ag (long dashed) and Pb–Pb (solid line). Since we compute the intial temperature from the collision energy our approach allows us to extrapolate as function of $`E/B`$, assuming that the stopping fraction for the collisional pressure is known. When we keep the stopping fraction constant and as given at the 160–200$`A`$ GeV collision energy, we find the results shown in Fig. 3.3. However, a constant stopping underestimates the intial temperature at lower collision energy, where we would expect higher stopping, and it overestimates the initial temperature at higher collision energies, where we would expect smaller stopping, thus we believe that the slope of the result we present in Fig. 3.3 is too steep. We will be able to improve on this result after the behavior of stopping as function of collision energy has been understood. In Fig. 3.3, the solid square is the result of the analysis for S-Au/W/Pb system, and open square for Pb–Pb as shown in the lower section of table 3.2. We note that the reason that the available energy $`E/B`$ in the fireball is the dominant parameter controlling strangeness yield is the cancellation of effect of higher initial temperature in the larger, more stopping systems, by the faster explosion of such a system, which leaves less time for strangeness production. We note that even though we did not analyze here the S–S system, for which case we would need to adapt the method to allow significant longitudonal flow, it is understood that the available fireball energy and strangeness content per baryon is higher in S–S 200$`A`$ GeV interactions,<sup>?</sup> consistent with the results shown in Fig. 3.3. This high strangeness yield corresponds to (above) equilibrium abundance phase space occupancy in hadronization. In the top section of table 3.2, the ratio $`\gamma _s/\gamma _q0.8`$, which corresponds (approximately) to the parameter $`\gamma _s`$ when $`\gamma _q=1`$ has been assumed. We observe that $`\gamma _s^{\mathrm{Pb}}>1`$. This strangeness over-saturation effect could arise from the effect of gluon fragmentation combined with early chemical equilibration in the QGP, $`\gamma _s(t<t_f)1`$. The ensuing rapid expansion preserves this high strangeness yield, and thus we find the result $`\gamma _s>1`$ , as we reported in Ref. <sup>?</sup>. This high phase space occupancy is one of the requirements for the enhancement of multi-strange (anti)baryon production, which is an important hadronic signal of QGP phenomena.<sup>?</sup> We compare this result of data analysis, in quantitative manner, with the theoretical computation of $`\gamma _s`$ which is easily obtained from the above study of total strangeness production, as we only need to divide the total momentary strangeness yield by the expected equilibrium abundance, for which we choose to consider ideal gas of strange quarks with QCD running mass $`m_s(\mu =5.5T)`$ . The factor 5.5 converts the value of $`T`$ into the appropriate scale $`\mu `$ of energy at which the kinetic equilibrium distribution is formed, and we note that $`m_s(T=182\mathrm{MeV})=200`$ MeV. The effect of QCD running influences the agreement between theory and experiment at the level of 10–15%. The result is shown in figure 3.3, right as function of time $`t`$ for the 160–200$`A`$ GeV collision systems and left as function of temperature $`T`$. Horizontal dotted line refers to equilibrium phase space occupancy, and the vertical line indicates expected freeze-out condition at $`T_f=143`$ MeV. Solid dots in figure 3.3a) show where this freeze-out temperature occurs as function of time $`t`$. The analysis point uncertainty in freeze-out time is obtained assuming that, in isentropic evolution, size scale $`R`$ and temperature satisfy $`RT=`$ Const., and thus with $`v_c=dR/dt\|_f`$ we find: $$\mathrm{\Delta }t=\frac{R_f}{v_cT}\mathrm{\Delta }T.$$ (23) The chemical freeze-out occurs at about 10 fm/$`c`$ and 13 fm/$`c`$, after onset of the collision, allowing for about 1 fm/$`c`$ initial time $`\tau _{\mathrm{ch}}`$ for the two systems S–Au/W/Pb and Pb–Pb respectively. Considering that the expansion velocity has been $`0.5c`$ and $`0.64c`$ respectively we obtain an estimate of the chemical freeze-out radius $`R_f^\mathrm{S}5`$fm and $`R_f^{\mathrm{Pb}}8`$fm, the latter value is in excellent agreement with the discussion of the Coulomb effect presented in section 2.2. ### 3.4 The Omega riddle The QGP formation and sudden disintegration model we have described above has natural limitations. When we attempt to describe within this approach most rarely produced particles, there is the potential for under-prediction of experimental results, which could receive contributions from other more effective production mechanisms. In this context, the most rarely produced hadron is the triply strange $`\mathrm{\Omega }(sss)`$ and $`\overline{\mathrm{\Omega }}(\overline{s}\overline{s}\overline{s})`$ which are the heaviest stable hadrons, $`M_\mathrm{\Omega }=1672`$ MeV. The phase space for $`\mathrm{\Omega }`$ is more than 10 times smaller than that for $`\mathrm{\Xi }`$ at the conditions of chemical freeze-out we have obtained. $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ production pattern can thus be altered by processes not implemented in the one stage fireball model used to analyze the data. When we attempted to describe along with the other hadrons the yields of $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ within the single stage freeze-out model, we indeed have discovered considerable loss of physical significance.<sup>?</sup> Already for the S–Au/W/Pb case, we have found that a more reliable description of the data arises if we did not consider the qualitative Omega yields available.<sup>?</sup> For the parameters as reported, we find in the Pb–Pb reactions that we under-predict the $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ yields by about factor 2. The experimental results are shown in the first three columns of table 6 and the theoretical yield computed using the Fermi-2000 model with parameters fixed by other particle abundances are shown in columns 4 and 5: we see that the presence of radial flow ($`v=v_c`$) has a minimal impact on the relative yields, compared to the case without radial flow ($`v=0`$). To put this result into proper perspective, consider that we find within the sudden hadronization of QGP with uncorrelated strange quarks in the deconfined phase an enhancement of $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ yields ‘only’ by factor 10 as compared to what is expected from extrapolation of p–A reactions. However, the experiment reports an enhancement by factor 15–20. Such a ‘failure’ is in fact confirming the early expectations that $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ yields are the best signature of deconfinement, considering the possibility of strange quark clustering.<sup>?</sup> In fact it is a bit surprising how well this early prediction works, and this requires further study to understand more precisely what exactly this means. Several groups have noted, studying the microscopic evolution of $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$, that due to low reaction cross section they decouple from hadron background somewhat sooner than all the other hadrons.<sup>?,?</sup> An early chemical freeze-out would impact statistical yields of $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ greatly. To augment the $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ yields by factor $`k`$, it is sufficient to take an incrementally $`\delta T`$ higher freeze-out temperature, as determined from study of the $`\mathrm{\Omega }`$ phase space: $$\delta TT\frac{\mathrm{ln}k}{M_\mathrm{\Omega }/T}.$$ (24) Thus in order to increase the yields by a factor 2 the $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ freeze-out would need to occur at $`T_\mathrm{\Omega }=150`$ MeV rather than at $`T_f=143`$ MeV. Since the temperature drops as the explosion of the fireball develops, this higher freeze-out temperature means an earlier in time freeze-out. Even if the required staging in time of hadron production is apparently small, a consistent picture requires fine-tuning and it seems unnatural, considering that all the other particles are perfectly consistent with just one sudden freeze-out condition. Pursuing other alternatives, we note that $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ enhancement is caused by strangeness pre-clustering in the deconfined phase which would enhance multistrange hadrons, but most prominently and noticeable enhance the phase space suppressed $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$. In this context, it is interesting to note that the missing yield is not symmetric: as seen in table 6 we miss in relative terms more $`\mathrm{\Omega }`$ than $`\overline{\mathrm{\Omega }}`$ . Interestingly, the missing yield is exactly proportional to the yield of $`\mathrm{\Xi }`$ and $`\overline{\mathrm{\Xi }}`$ and the best description of all particle yields, including all $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ is arrived at describing what is missing as proportional (11 %) to the $`\mathrm{\Xi }`$ and $`\overline{\mathrm{\Xi }}`$ yield, this is shown in the last column ‘11 %’ of table 6. It is now easy to propose a model that would lead just to this result: there are colored di-strange quarks clusters at hadronization and when their color strings break $`\mathrm{\Xi }`$ and $`\mathrm{\Omega }`$ are produced. This imprints a ‘shadow’ of $`\mathrm{\Xi }`$ and $`\overline{\mathrm{\Xi }}`$ in the $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$-abundance. While this works for $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$, we find that this mechanism is not compatible with the other particle abundances, in other words a similar ‘shadow’ of $`\mathrm{\Lambda }`$ and $`\overline{\mathrm{\Lambda }}`$ in the $`\mathrm{\Xi }`$ and $`\overline{\mathrm{\Xi }}`$ channel seems unacceptable. Thus this mechanism would work only if pairing of strange quarks would be significant near to phase transition. Current models of ‘color super conductivity’ support such a clustering mechanism for additional $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ enhancement, though detailed studies are still in progress.<sup>?</sup> We have also explored the possibility that unknown $`\mathrm{\Omega }^{}`$ and $`\overline{\mathrm{\Omega }^{}}`$ resonances contribute to the $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ yield, but we were not able to find a good set of parameters for these hypothetical resonances. Moreover, this hypothesis implies a baryon–antibaryon symmetric contribution in the sense that both $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ yields are multiplied by the same factor. However, the missing yield is clearly also baryon–antibaryon asymmetric — thus despite several ad-hoc parameters the model description remains poor. We note that earlier statistical descriptions of $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ yields have not been sensitive to the problems we described.<sup>?,?</sup> In fact as long as the parameter $`\gamma _q`$ is not considered, it is not possible to describe the experimental data at the level of precision that would allow recognition of the $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ yield as a problem for the statistical Fermi phase space model. ## 4 Kinetic Strangeness Production In some computational details, the methods to describe strangeness production differ.<sup>?,?,?,?</sup> This leads to different expectations regarding chemical equilibration of quark flavor at RHIC energies, with some authors finding marginal at best chemical equilibration. We therefore develop in more detail the computational approach which is consistent with the SPS-energy scale results discussed in previous sections.<sup>?</sup> One important difference to the earlier work is that the two loop level running of QCD parameters for both coupling strength $`\alpha _s`$ and strange quark mass $`m_s`$ is used. $`\alpha _{M_Z}=0.118`$ is assumed as determined at the $`\mu =M_{Z^0}`$ energy scale. Another improvement is that an entropy conserving explosive flow of matter is incorporated directly into the dynamical equations describing the evolution of strangeness phase space occupancy. This approach is entailing significant cancellations in the dynamical equations and the only model dependence on matter flow which remains is the relationship between the local temperature and local proper time. In consequence, a relatively simple and physically transparent model for the evolution of the phase space occupancy $`\gamma _s`$ of strange quarks in the expanding QGP can be studied. We use two assumptions of relevance for the results we obtain: $``$ the kinetic (momentum distribution) equilibrium is reached faster than the chemical (abundance) equilibrium<sup>?,?</sup>; $``$ gluons equilibrate chemically significantly faster than strangeness.<sup>?</sup> The first assumption allows us to study only the chemical abundances, rather than the full momentum distribution, which simplifies greatly the structure of the master equations; the second assumption allows us to consider the evolution of the strangeness population only after an initial time $`\tau _{\mathrm{ch}}`$ period has passed: $`\tau _{\mathrm{ch}}`$ is the time required for the development to near chemical equilibrium of the gluon population, and the corresponding temperature $`T_{\mathrm{ch}}`$ is the initial condition we need to compute the evolution of strangeness. Aside of $`T_{\mathrm{ch}}`$, the strange quark mass $`m_s`$ introduces the greatest uncertainty that enters strangeness yield calculations based on resumed perturbative QCD rates.<sup>?</sup> The overpopulation of the strangeness phase space, seen before in section 3 in SPS data, arises in particular for $`T_{\mathrm{ch}}>250`$ MeV and values of strange quark mass $`m_s`$(1GeV)$`200\pm 20`$ MeV. In view of these assumptions the phase space distribution $`f_s`$ can be characterized by a local temperature $`T(\stackrel{}{x},t)`$ of a (Boltzmann) equilibrium distribution $`f_s^{\mathrm{}}`$ , with normalization set by a phase space occupancy factor: $$f_s(\stackrel{}{p},\stackrel{}{x};t))\gamma _s(T)f_s^{\mathrm{}}(\stackrel{}{p};T).$$ (25) Eq. (25) invokes in the momentum independence of $`\gamma _s`$ the first assumption. More generally, the factor $`\gamma _i,i=g,q,s,c`$, allows the local density of gluons, light quarks, strange quarks and charmed quarks, respectively to evolve independently of the local momentum shape. With variables $`(t,\stackrel{}{x})`$ referring to an observer in the laboratory frame, the chemical evolution can be described by the strange quark current non-conservation arising from strange quark pair production described by a Boltzmann collision term: $`_\mu j_s^\mu {\displaystyle \frac{\rho _s}{t}}+{\displaystyle \frac{\stackrel{}{v}\rho _s}{\stackrel{}{x}}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\rho _g^2(t)\sigma v_T^{ggs\overline{s}}`$ (26) $`+`$ $`\rho _q(t)\rho _{\overline{q}}(t)\sigma v_T^{q\overline{q}s\overline{s}}\rho _s(t)\rho _{\overline{\mathrm{s}}}(t)\sigma v_T^{s\overline{s}gg,q\overline{q}}.`$ The factor 1/2 avoids double counting of gluon pairs. The implicit sums over spin, color and any other discreet quantum numbers are combined in the particle density $`\rho =_{s,c,\mathrm{}}d^3pf`$, and we have also introduced the momentum averaged production/annihilation thermal reactivities (also called ‘rate coefficients’): $$\sigma v_{\mathrm{rel}}_T\frac{d^3p_1d^3p_2\sigma _{12}v_{12}f(\stackrel{}{p}_1,T)f(\stackrel{}{p}_2,T)}{d^3p_1d^3p_2f(\stackrel{}{p}_1,T)f(\stackrel{}{p}_2,T)}.$$ (27) $`f(\stackrel{}{p}_i,T)`$ are the relativistic Boltzmann/Jüttner distributions of two colliding particles of momentum $`p_i`$, $`i=1,2`$. The current conservation used above in the laboratory ‘Eulerian’ formulation can also be written with reference to the individual particle dynamics in the so called ‘Lagrangian’ description: consider $`\rho _s`$ as the inverse of the small volume available to each particle. Such a volume is defined in the local frame of reference for which the local flow vector vanishes $`\stackrel{}{v}(\stackrel{}{x},t)\|_{\text{local}}=0`$. The considered volume $`\delta V_l`$ being occupied by small number of particles $`\delta N`$ (e.g., $`\delta N=1`$), we have: $$\delta N_s\rho _s\delta V_l.$$ (28) The left hand side (LHS) of Eq. (26) can be now written as: $$\frac{\rho _s}{t}+\frac{\stackrel{}{v}\rho _s}{\stackrel{}{x}}\frac{1}{\delta V_l}\frac{d\delta N_s}{dt}=\frac{d\rho _s}{dt}+\rho _s\frac{1}{\delta V_l}\frac{d\delta V_l}{dt}.$$ (29) Since $`\delta N`$ and $`\delta V_ldt`$ are L(orentz)-invariant, the actual choice of the frame of reference in which the right hand side (RHS) of Eq. (29) is studied is irrelevant and we drop henceforth the subscript $`l`$. We can further adapt Eq. (29) to the dynamics we pursue: we introduce $`\rho _s^{\mathrm{}}(T)`$ as the (local) chemical equilibrium abundance of strange quarks, thus $`\rho =\gamma _s\rho _s^{\mathrm{}}`$. We evaluate the equilibrium abundance $`\delta N_s^{\mathrm{}}=\delta V\rho _s^{\mathrm{}}(T)`$ integrating the Boltzmann distribution: $$\delta N_s^{\mathrm{}}=[\delta VT^3]\frac{3}{\pi ^2}z^2K_2(z),z=\frac{m_s}{T}.$$ (30) We will below use: $`d[z^\nu K_\nu (z)]/dz=z^\nu K_{\nu 1}`$, where $`K_\nu `$ is the modified Bessel function of order $`\nu `$. The first factor on the RHS in Eq. (30) is a constant in time should the evolution of matter after the initial pre-thermal time period $`\tau _0`$ be entropy conserving,<sup>?</sup> and thus $`\delta VT^3=\delta V_0T_0^3=`$ Const. . We now substitute in Eq. (29) and obtain $$\frac{\rho _s}{t}+\frac{\stackrel{}{v}\rho _s}{\stackrel{}{x}}=\dot{T}\rho _s^{\mathrm{}}\left(\frac{d\gamma _s}{dT}+\frac{\gamma _s}{T}z\frac{K_1(z)}{K_2(z)}\right),$$ (31) where $`\dot{T}=dT/dt`$. Note that, in Eq. (31), only a part of the usual flow-dilution term is left, since we implemented the adiabatic volume expansion, and study the evolution of the phase space occupancy in lieu of particle density. The dynamics of the local temperature is the only quantity we need to model. We now return to study the collision terms seen on the RHS of Eq. (26). A related quantity is the (L-invariant) production rate $`A^{1234}`$ of particles per unit time and space, defined usually with respect to chemically equilibrated distributions: $$A^{1234}\frac{1}{1+\delta _{1,2}}\rho _1^{\mathrm{}}\rho _2^{\mathrm{}}\sigma _sv_{12}_T^{1234}.$$ (32) The factor $`1/(1+\delta _{1,2})`$ is introduced to compensate double-counting of identical particle pairs. In terms of the L-invariant $`A`$ , Eq. (26) takes the form: $`\dot{T}\rho _s^{\mathrm{}}\left({\displaystyle \frac{d\gamma _s}{dT}}+{\displaystyle \frac{\gamma _s}{T}}z{\displaystyle \frac{K_1(z)}{K_2(z)}}\right)=\gamma _g^2(\tau )A^{ggs\overline{s}}+`$ $`+\gamma _q(\tau )\gamma _{\overline{q}}(\tau )A^{q\overline{q}s\overline{s}}\gamma _s(\tau )\gamma _{\overline{s}}(\tau )(A^{s\overline{s}gg}+A^{s\overline{s}q\overline{q}}).`$ (33) Only weak interactions convert quark flavors, thus, on hadronic time scale, we have $`\gamma _{s,q}(\tau )=\gamma _{\overline{s},\overline{q}}(\tau )`$. Moreover, detailed balance, arising from the time reversal symmetry of the microscopic reactions, assures that the invariant rates for forward/backward reactions are the same, specifically $$A^{1234}=A^{3412},$$ (34) and thus: $`\dot{T}\rho _s^{\mathrm{}}\left({\displaystyle \frac{d\gamma _s}{dT}}+{\displaystyle \frac{\gamma _s}{T}}z{\displaystyle \frac{K_1(z)}{K_2(z)}}\right)`$ $`=`$ $`\gamma _g^2(\tau )A^{ggs\overline{s}}\left[1{\displaystyle \frac{\gamma _s^2(\tau )}{\gamma _g^2(\tau )}}\right]`$ (35) $`+\gamma _q^2(\tau )A^{q\overline{q}s\overline{s}}\left[1{\displaystyle \frac{\gamma _s^2(\tau )}{\gamma _q^2(\tau )}}\right].`$ When all $`\gamma _i1`$, the Boltzmann collision term vanishes, we have reached equilibrium. As discussed, the gluon chemical equilibrium is thought to be reached at high temperatures well before the strangeness equilibrates chemically, and thus we assume this in what follows, and the initial conditions we will study refer to the time at which gluons are chemically equilibrated. Setting $`\lambda _g=1`$ (and without a significant further consequence for what follows, since gluons dominate the production rate, also $`\lambda _q=1`$), we obtain after a straightforward manipulation the dynamical equation describing the evolution of the local phase space occupancy of strangeness: $$2\tau _s\dot{T}\left(\frac{d\gamma _s}{dT}+\frac{\gamma _s}{T}z\frac{K_1(z)}{K_2(z)}\right)=1\gamma _s^2.$$ (36) Here, we defined the relaxation time $`\tau _s`$ of chemical (strangeness) equilibration as the ratio of the equilibrium density that is being approached, with the rate at which this occurs: $$\tau _s\frac{1}{2}\frac{\rho _s^{\mathrm{}}}{(A^{ggs\overline{s}}+A^{q\overline{q}s\overline{s}}+\mathrm{})}.$$ (37) The factor 1/2 is introduced by convention in order for the quantity $`\tau _s`$ to describe the exponential approach to equilibrium. Eq. (36) is the final analytical result describing the evolution of phase space occupancy. Since one generally expects that $`\gamma _s1`$ in a monotonic fashion as function of time, it is important to appreciate that this equation allows the range $`\gamma _s>1`$: when $`T`$ drops below $`m_s`$, and $`1/\tau _s`$ becomes small, the dilution term (2nd term on LHS) in Eq. (36) dominates the evolution of $`\gamma _s`$ . In simple terms, the high abundance of strangeness produced at high temperature over-populates the available phase space at lower temperature, when the equilibration rate cannot keep up with the expansion cooling. This behavior of $`\gamma _s`$ has been shown for the SPS conditions allowing explosive transverse expansion in subsection 3.3. Since we assume that the dynamics of transverse expansion of QGP is similar at RHIC as at SPS, we obtain similar behavior for $`\gamma _s`$ in section 5 below. $`\tau _s(T)`$ , Eq. (37), has been evaluated using pQCD cross section and employing next to leading order running of both the strange quark mass and QCD-coupling constant $`\alpha _s`$.<sup>?</sup> We believe that this method produces a result for $`\alpha _s`$ that can be trusted down to just below 1 GeV energy scale which is here relevant. We employ results obtained with $`\alpha _s(M_{Z^0})=0.118`$ and $`m_s(\text{1GeV})=200`$ MeV; we have shown results with $`m_s(\text{1GeV})=220`$ MeV earlier.<sup>?</sup> There is some systematic uncertainty due to the appearance of the strange quark mass as a fixed rather than running value in both, the chemical equilibrium density $`\rho _s^{\mathrm{}}`$ in Eq. (37), and in the dilution term in Eq. (36). We use the value $`m_s(\text{1\hspace{0.17em}GeV})`$, with the 1 GeV energy scale chosen to correspond to typical interaction scale in the QGP at temperatures under consideration. ## 5 Expectations for Strange Hadron Production at RHIC We now combine all recent advances in theoretical models of strangeness production and data interpretation at SPS energies with the objective of making reliable predictions for the RHIC energy range.<sup>?</sup> First we address the question how much strangeness can be expected at RHIC. The numerical study of Eq. (36) becomes possible as soon as we define the temporal evolution of the temperature for RHIC conditions. We expect that a global cylindrical expansion should describe the dynamics: aside of the longitudinal flow, we allow the cylinder surface to expand given the internal thermal pressure. SPS experience suggests that the transverse matter flow will not exceed the sound velocity of relativistic matter $`v_{}c/\sqrt{3}`$. We recall that for a pure longitudinal expansion local entropy density scales according to $`ST^31/\tau `$.<sup>?</sup> It is likely that the transverse flow of matter will accelerate the drop in entropy density. We thus consider the following temporal evolution function of the temperature: $$T(\tau )=T_0\left[\frac{1}{(1+\tau 2c/d)(1+\tau v_{}/R_{})^2}\right]^{1/3}.$$ (38) We take the thickness of the initial collision region at $`T_0=0.5`$ GeV to be $`d(T_0=0.5)/2=0.75`$ fm, and the transverse dimension in nearly central Au–Au collisions to be $`R_{}=4.5`$ fm. The time at which thermal initial conditions are reached is assumed to be $`\tau _0=1`$fm/$`c`$. When we vary $`T_0`$, the temperature at which the gluon equilibrium is reached, we also scale the longitudinal dimension according to: $$d(T_0)=(0.5\text{ GeV}/T_0)^31.5\text{ fm}.$$ (39) This assures that when comparing the different evolutions of $`\gamma _s`$ we are looking at an initial system that has the same entropy content by adjusting its initial volume $`V_0`$. The reason we vary the initial temperature $`T_0`$ down to 300 MeV, maintaining the initial entropy content is to understand how the assumption about the chemical equilibrium of gluons, reached by definition at $`T_0`$, impacts strangeness evolution. In fact when considering decreasing $`T_0`$ (and thus increasing $`V_0`$), the thermal production is turned on at a later time in the history of the collision. The numerical integration of Eq. (36) is started at $`\tau _0`$, and a range of initial temperatures $`300T_0600`$, varying in steps of 50 MeV. The high limit of the temperature we explore exceeds somewhat the ‘hot glue scenario’,<sup>?</sup> while the lower limit of $`T_0`$ corresponds to the more conservative estimates of possible initial conditions.<sup>?</sup> Since the initial $`p`$$`p`$ collisions also produce strangeness, we take as an estimate of initial abundance a common initial value $`\gamma _s(T_0)=0.15`$. The time evolution in the plasma phase is followed up to the break-up of QGP. This condition we establish in view of results of the analysis for SPS presented in section 3. We recall that SPS-analysis showed that the system dependent baryon and antibaryon $`m_{}`$-slopes of particle spectra are result of differences in collective flow in the deconfined QGP source at freeze-out. In consequence there is universality of physical properties of hadron chemical freeze-out between different SPS systems. This value is nearly applicable to RHIC conditions, as can be seen extrapolating the phase boundary curve to the small baryochemical potentials. The QGP break-up temperature $`T_f^{\text{SPS}}143\pm 5`$ MeV will see just a minor upward change, and we adopt here the value $`T_f^{\text{RHIC}}150\pm 5`$ MeV. With the freeze-out condition fixed, one would think that the major remaining uncertainty comes from the initial gluon equilibration temperature $`T_0`$, and we now study how different values of $`T_0`$ influence the final state phase space occupancy. We integrate numerically Eq. (36) and present $`\gamma _s`$ as function of both time $`t`$ in Fig. 5a, and temperature $`T`$ in Fig. 5b, up to the expected QGP breakup at $`T_f^{\text{RHIC}}150\pm 5`$ MeV. We see that: $``$ widely different initial conditions (with similar initial entropy content) lead to rather similar chemical conditions at chemical freeze-out of strangeness, $``$ despite a series of conservative assumptions, we find, not only, that strangeness equilibrates, but indeed that the dilution effect allows an overpopulation of the strange quark phase space. For a wide range of initial conditions, we obtain a narrow band $`1.15>\gamma _s(T_f)>1`$ . We will in the following, taking into account some contribution from hadronization of gluons in strange/antistrange quarks, adopt what the value $`\gamma _s(T_f)=1.25`$. We now consider how this relatively large value of $`\gamma _s`$, characteristic for the underlying QGP formation and evolution of strangeness, impacts the strange baryon and anti-baryon observable emerging in hadronization. Remembering that major changes compared to SPS should occur in rapidity spectra of mesons, baryons and antibaryons, we will apply the same hadronization model that worked in the analysis of the SPS data. This hypothesis can be falsified easily, since based and compared to the Pb–Pb 158$`A`$ GeV results it implies: a) shape identity of all RHIC $`m_{}`$ and $`y`$ spectra of antibaryons $`\overline{p},\overline{\mathrm{\Lambda }},\overline{\mathrm{\Xi }},`$ since there is no difference in their production mechanism, and the form of the spectra is determined in a similar way to SPS energy range by the local temperature and flow velocity vector; b) the $`m_{}`$-inverse-slopes of these antibaryons should be very similar to the result obtained at CERN for Pb–Pb 158$`A`$ GeV, since the expected 3% increase in the freeze-out temperature is accompanied by a comparable increase in collective transverse flow. The abundances of particles produced from QGP within the sudden freeze-out model are controlled by several parameters we addressed earlier: the light quark fugacity $`1<\lambda _q<1.1`$ , value is limited by the expected small ratio between baryons and mesons (baryon-poor plasma) when the energy per baryon is above 100 GeV, strangeness fugacity $`\lambda _s1`$ which value for locally neutral plasma assures that $`s\overline{s}=0`$; the light quark phase space occupancy $`\gamma _q1.5`$, overabundance value due to gluon fragmentation. Given these narrow ranges of chemical parameters and the freeze-out temperature $`T_f=150`$ MeV, we compute the expected particle production at break-up. In general, we cannot expect that the absolute numbers of particles we find are correct, as we have not modeled the important effect of flow in the laboratory frame of reference. However, ratios of hadrons subject to similar flow effects (compatible hadrons) can be independent of the detailed final state dynamics, as the results seen at SPS suggest, and we will look at such ratios more closely. Taking $`\gamma _q=1.5\begin{array}{c}+0.10\\ 0.25\end{array}`$, we choose the value of $`\lambda _q`$, see the header of table 5, for which the energy per baryon ($`E/B`$) is similar to the collision condition (100 GeV), which leads to the range $`\lambda _q=1.03\pm 0.005`$. We evaluate for these examples aside of $`E/B`$, the strangeness per baryon $`s/B`$ and entropy per baryon $`S/B`$ as shown in the top section of the table 5. We do not enforce $`s\overline{s}=0`$ exactly, but since baryon asymmetry is small, strangeness is balanced to better than 2% in the parameter range considered. In the bottom portion of table 5, we present the compatible particle abundance ratios, computed according to the procedure developed in section 2. We have given, aside of the baryon and antibaryon relative yields, also the relative kaon yield, which is also well determined within this approach. The meaning of these results can be better appreciated when we assume in an example the central rapidity density of direct protons is $`dp/dy\|_{\text{cent.}}=25`$. In table 5, we present the resulting (anti)baryon abundances. The net baryon density $`db/dy16\pm 3`$, there is baryon number transparency. We see that (anti)hyperons are indeed more abundant than non-strange (anti)baryons. Taking into account the disintegration of strange baryons, we are finding a much greater number of observed protons $`dp/dy\|_{\text{cent.}}^{\text{obs.}}65\pm 5`$ in the central rapidity region. It is important when quoting results from table 5 to recall that: 1) we have chosen arbitrarily the overall normalization in table 5 , only particle ratios were computed, and 2) the rapidity baryon density relation to rapidity proton density is a consequence of the assumed value of $`\lambda _q`$, which we chose to get $`E/B100`$ GeV per participant. The most interesting result seen in table 5 , the hyperon-dominance of the baryon yields at RHIC, does not depend on detailed model hypothesis. We have explored another set of parameters in our first and preliminary report on this matter,<sup>?</sup> finding this result. Another interesting property of the hadronizing hot RHIC matter, as seen in table 5, is that strangeness yield per participant is expected to be 13–23 times greater than seen at present at SPS energies, where we have 0.75 strange quark pairs per baryon. As seen in table 5, the baryon rapidity density is in this examples similar to the proton rapidity density. ## 6 Summary and Conclusions We believe that this study of SPS strangeness results decisively shows interesting new physics. We see considerable convergence of the results around properties of suddenly hadronizing QGP.<sup>?</sup> The key results we obtained in the Fermi-2000 model data analysis are: 1) the same hadronization temperature $`T`$=142–144 MeV for very different collision systems with different hadron spectra; 2) QGP expected results for the source phase space properties: $`\stackrel{~}{\lambda }_s=1`$ for both S- and Pb-collisions, implying $`\lambda _s^{\mathrm{Pb}}1.1`$ ; 3) $`\gamma _s^{\mathrm{Pb}}>1`$, indicating that a high strangeness yield was reached before freeze-out; 4) $`\gamma _q>1`$ as would be expected from a high entropy phase and the associated value $`S/B43\pm 3`$ ; 5) the yield of strangeness per baryon $`\overline{s}/B0.7`$ just as predicted by gluon fusion in thermal QGP, a point we studied in detail in section 3.3; 6) the transverse expansion velocity for Pb–Pb: $`v_c^{\mathrm{Pb}}1/\sqrt{3}`$, just below the sound velocity of quark matter. The universality of the physical properties at chemical freeze-out for S- and Pb-induced reactions points to a common nature of the primordial source of hadronic particles. The difference in spectra between the two collision systems considered arises, in this analysis, due to the difference in the degree of chemical equilibration of light and strange quarks, expected for systems of differing size and lifespan, and a difference in the collective surface explosion velocity, $`v_c^\mathrm{S}0.5<v_c^{\mathrm{Pb}}1/\sqrt{3}`$ , which for larger system is higher, having more time to develop. Considering how small the experimental WA97 spectral slope errors shown in table 3.1 are presently, there is now overwhelming evidence that the production mechanism of both $`\mathrm{\Lambda }`$ and $`\overline{\mathrm{\Lambda }}`$ is the same, which observation is very probably also true for both $`\mathrm{\Xi }`$ and $`\overline{\mathrm{\Xi }}`$. This symmetry between matter–antimatter production is an important cornerstone of the claim that the strange antibaryon data can only be interpreted in terms of direct emission from a deconfined and thus matter-antimatter symmetric quark matter. We note that the QGP break-up temperature we find, $`T_f=143`$ MeV, corresponds to an energy density $`\epsilon =𝒪`$(0.5) GeV/fm<sup>3</sup>.<sup>?</sup> Among other interesting results which also verify the consistency of the experimental data understanding within the Fermi-2000 model, we recall: $``$ the exact balancing of strangeness $`\overline{s}s=0`$ also in the final hadronic particles in the symmetric Pb–Pb case; $``$ the increase of the baryochemical potential $`\mu _B^{\mathrm{Pb}}=203\pm 5>\mu _B^\mathrm{S}=178\pm 5`$ MeV as the collision system grows; $``$ the energy per baryon near to the value expected if energy and baryon number deposition in the fireball are similar; $``$ hadronization into pions at $`\gamma _q\gamma _q^c=e^{m_\pi /2T_f}1.6`$ seen in Pb–Pb reactions, which is an effective way to convert excess of entropy in the plasma into hadrons, without need for reheating, or a mixed phase; the finding of the maximum allowable $`\gamma _q`$ is intrinsically consistent with the notion of an explosively disintegrating QGP phase. A reassuring feature of the Pb–Pb analysis related to chemical equilibration has been described in subsection 2.3: we find a pion yield which maximizes the entropy density of hadronic particles produced.<sup>?</sup> This detailed technical result explains how sudden hadronization can occur: in general the deconfined state with broken color bonds and thus the high entropy density has to find an exit into the hadronic world, maintaining or increasing the total entropy and preferentially also the local entropy contained within a small, comoving volume cell. Our analysis of experimental results suggests that this is accomplished by generating an over-saturated pion phase space, in which the entropy density rises to values as high as are believed to occur in QGP at hadronization. Chemical equilibrium hadronization requires the formation of a mixed plasma-hadron gas phase and is generally believed to require a relatively long time, followed by kinetic reequilibration. In our opinion such a hadronization model is now inconsistent with the experimental strange baryon and antibaryon data on yields and spectra. The reader should note that such technical differences between different groups about the dynamics of the evolution of the hadron fireball after the deconfined phase has hadronized, do not impact the primary agreement about the deconfined nature of the high density source of hadronic particles. We believe that omission to consider chemical non-equilibrium in the study of freeze-out conditions employing the analysis of spectral shape (flow) and also pion correlation (HBT) effect is the source of the difference of results here presented with some other recent work.<sup>?,?</sup> To understand the source of this difference it is important to realize that there is a considerable influence on the shape of pion spectra by the light quark chemical non-equilibrium which the data analysis presented includes: the cocktail of resonance decays contributing to pion spectra is altered, and moreover, there is spectral deformation at low $`m_{}`$ due to pion correlation effects caused by the overpopulated phase space.<sup>?</sup> We note that results presented also differ somewhat from the WA98 experiment analysis addressing solely $`\pi ^0`$ spectra, and which again assumes pion chemical equilibrium.<sup>?</sup> In consequence, the $`\pi ^0`$-freeze-out conditions as seen in Ref. <sup>?</sup>, Table 1 are different from those determined here. On the other hand, another recent hadron spectral shape analysis,<sup>?</sup> which did not introduce low $`m_{}`$ pion spectra into consideration obtains a chemical freeze-out conditions nearly identical to those we discussed. While the precise understanding of hadronization condition is required for a measurement of physical properties of QGP including the latent heat, the differences discussed are of little if any consequence concerning the fundamental issue, the question if deconfinement is achieved. The sudden hadronization of entropy rich QGP leads to value $`\gamma _q\gamma _q^c`$, in order to connect the entropy rich deconfined and the confined phases more efficiently. The dominant pion contribution to the entropy density (and pressure) is nearly twice as high at $`\gamma _q\gamma _q^c`$ than at $`\gamma _q=1`$. Without this phenomenon one has to introduce a mechanism that allows the parameter $`VT^3`$ to grow, thus expanding either the volume $`V`$ due to formation of the mixed phase or invoking a rise of $`T`$ in the reheating. The range of values for $`\gamma _q`$ is bounded from above by the Bose distribution singularity $`\gamma _q\gamma _q^c`$, but a pion condensate is not formed since it ‘consumes’ energy without consuming entropy of the primordial high entropy QGP phase. An interesting feature of such a mechanism of phase transition which side-steps the need to form a mixed phase or reheating is that the chemical non-equilibrium reduces and potentially eliminates any discontinuity in the phase transition. This being the case, experimental searches will not find the critical fluctuations expected for a discontinuous phase change, even if theory implies a 1st order phase transition for the statistical equilibrium system. This is in agreement with the failure of NA49 experiment to find precritical fluctuations in event-by-event analysis.<sup>?</sup> The only not fully quantitatively described particle yields are $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$: for the parameters we find, the Fermi-2000 model applied to Pb–Pb reactions under-predicts this smallest of all hadronic abundances by about factor 2 . This means that we expect, within the sudden hadronization of QGP with uncorrelated strange quarks in the deconfined phase, only an enhancement of $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ yields by a factor 8-10 as compared to what is expected from extrapolation of p–A reactions. Since the experiment reports an enhancement by factor 15–20, we need to think again. This ‘failure’ of Fermi-2000 model is in fact confirming the early expectation that $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ yields are the best signature of deconfinement,<sup>?</sup> we just must in future address the question what exactly this tells us about QGP structure. We have argued nearly 20 years ago that strangeness pairing in the color anti-triplet channel $`(ss)_{\overline{3}}`$ in the QGP source would enhance $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Omega }}`$ yields,<sup>?</sup> a point that is of some topical interest today in context of color superdconductivity studies.<sup>?</sup> Despite this unexpected additional enhancement, we firmly conclude in view of all diverse evidence that (multi)strange hadronic particles seen at CERN-SPS are emerging from a deconfined QGP phase of hadronic matter and do not undergo a re-equilibration after they have been produced. This finding has encouraged us to consider within the same computational scheme the production of strange hadrons at RHIC conditions. First, we have shown that one can expect strangeness chemical equilibration in nuclear collisions at RHIC if the deconfined QGP is formed. There will, as at SPS, be overpopulation effect associated with the early strangeness abundance freeze-out before hadronization. Most importantly for signatures of new physics at RHIC, we found that strange (anti)baryon abundances will be greater than the yields of non-strange baryons (protons, neutrons). Consequently, the rapidity distributions of (anti)protons are arising from decays of (anti)hyperons. We are not aware that microscopic model studies reported in the literature about RHIC conditions which have noted this remarkable hyperon dominance result, see, e.g., Ref. <sup>?</sup>. The reader could wonder why is this unusual phenomenon not happening at SPS energies described in section 3? At SPS there is still an appreciable relative baryon abundance among all hadrons (about 15%) and the strangeness yield is at SPS energies only at a level similar to the baryon yield. Thus while abundant (anti)hyperon formation begins to set in, there are still many non-strange (anti)baryons produced. With increasing per baryon energy the yield of strange quark pairs per baryon rises, and at the same time the relative abundance of baryons among all hadrons diminishes. As result, at RHIC energies, we have predicted that hyperons and/or antihyperons are the dominant population fraction among all baryons and/or antibaryons. We thus believe that the preponderance of hyperons as the dominant (anti)baryon population at RHIC energies can be uniquely correlated with the formation and sudden hadronization of deconfined QGP phase. In a nutshell: we find that strangeness and (anti)baryon QGP signatures are conclusively proving formation of deconfined quark matter phase at SPS energies, and that these signatures of new physics are much more distinct at higher RHIC energies. Acknowledgment: We thank the editor of Int. J. Mod. Phys. E, Ernest Henley, and Keith Dienes for valuable comments and suggestions. References
warning/0003/hep-th0003290.html	ar5iv	text	# Alberta-Thy-02-00 The Dimensional-Reduction Anomaly in Spherically Symmetric Spacetimes ## 1 Introduction Spacetimes with continuous symmetries play an important role in quantum field theory for the reason that symmetries often allow one to reduce greatly the computational difficulty of a given problem, making practical calculations feasible. Of particular interest is the separation of variables approach to solving partial differential equations in geometries with a high degree of symmetry. For example, expanding a field propagating in a spherically symmetric geometry in terms of spherical harmonics and substituting into the field equation reduces the system from that of a single field $`\widehat{\mathrm{\Phi }}(t,r,\theta ,\varphi )`$ in four dimensions to a collection of effective two-dimensional fields $`\widehat{\phi }_{\mathrm{}}(t,r)`$, one for each spherical harmonic $`Y_\mathrm{}m(\theta ,\varphi )`$. In principle, after solving the simpler two-dimensional problems, one can obtain quantities like the stress tensor or the effective action for the original four-dimensional field theory by summing the corresponding results for the two-dimensional field theories over all modes. In a previous paper , it was noted that separation of variables can break down when applied to quantum field theory, so that summing over the dimensionally reduced results no longer yields the corresponding quantity in four dimensions. This occurs because in quantum field theory, quantities of physical interest, such as the effective action, stress energy tensor, and square of the field operator, are divergent and must be renormalized. While the bare field can be dimensionally reduced into the sum of lower-dimensional fields, the divergent parts which are to be subtracted in four dimensions generally do not equal the sum of the corresponding divergent terms from the two-dimensional theories. As a result, one obtains an incorrect answer if one calculates a renormalized quantity in four dimensions by summing over modes of the corresponding renormalized quantities in two dimensions. This failure of dimensional reduction when applied to quantum field theories is called the dimensional-reduction anomaly. In this paper we calculate the dimensional-reduction anomaly which occurs when a scalar field propagating in a spherically symmetric four-dimensional spacetime is decomposed into spherical harmonics and treated as a collection of two-dimensional fields. This case may be of particular importance to recent attempts to calculate the stress tensor and Hawking radiation in black-hole spacetimes using two-dimensional dilaton gravity models . Since four-dimensional renormalized quantities will not in general equal the sum of the corresponding two-dimensional quantities, it may well be necessary to take into account the contribution of the dimensional-reduction anomaly in order to reproduce the correct results for four dimensions (see also ). We begin in Section 2 with a brief discussion of dimensional reduction in spherically symmetric spacetimes. In Section 3 we examine the simple case of flat space, both to illustrate the basic idea behind the anomaly and to lay the necessary computational groundwork. In Sections 4 and 5 we extend our calculations to a general spherically symmetric four-dimensional space, and calculate the dimensional-reduction anomalies in $`\widehat{\mathrm{\Phi }}^2`$ and the effective action. We conclude with a brief discussion of the possible implications of the anomaly. We work in Euclidean signature, using dimensionless units where $`G=c=\mathrm{}=1`$ and the sign conventions of for the definition of the curvature. ## 2 Spherical Decompositions In this section we briefly consider the dimensional reduction of a quantum field in a four-dimensional spherically symmetric space, and show how it may be reduced to a collection of two-dimensional fields. The line element for such a space may be written as $$ds^2=g_{\mu \nu }(X^\tau )dX^\mu dX^\nu =h_{ab}(x^c)dx^adx^b+\rho ^2\mathrm{e}^{2\varphi (x^c)}\omega _{ij}(y^k)dy^idy^j,$$ (2.1) where $`X^\alpha =(x^a,y^i)`$, $`h_{ab}`$ is an arbitrary two-dimensional metric, $`\omega _{ij}`$ is the metric of a two-sphere, $`\rho `$ is a constant with the dimensions of length, and $`\varphi `$ is known as the dilaton. The radius of a two-sphere of fixed $`x^a`$ is given by $`r=\rho \mathrm{e}^{\varphi (x^a)}`$. Consider a massive scalar field propagating on the space (2.1) and obeying the field equation $$F\widehat{\mathrm{\Phi }}(X)\left(\mathrm{}m^2V\right)\widehat{\mathrm{\Phi }}(X)=0,$$ (2.2) where the potential $`V`$ is also spherically symmetric. The corresponding Green function is a solution of the equation $$FG(X,X^{})=\delta (X,X^{}).$$ (2.3) Knowledge of the Green function for a given quantum state allows one to calculate other expectation values of interest, such as the square of the field operator, $`\widehat{\mathrm{\Phi }}^2`$, and the stress tensor, $`\widehat{T}_{\mu \nu }`$. Now consider what happens if we decompose $`\widehat{\mathrm{\Phi }}`$ in terms of spherical harmonics $`Y_\mathrm{}m(y^i)`$ as follows: $$\widehat{\mathrm{\Phi }}(X)=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\widehat{\phi }_{\mathrm{}}(x^a)\frac{Y_\mathrm{}m(y^i)}{r}.$$ (2.4) Substitution into (2.2) shows that $`\widehat{\phi }_{\mathrm{}}`$ behaves as a field propagating in the two-dimensional space with line element $$ds^2=h_{ab}(x^c)dx^adx^b,$$ (2.5) and satisfying the field equation $$_{\mathrm{}}\widehat{\phi }_{\mathrm{}}(x)\left(\mathrm{\Delta }m^2V_{\mathrm{}}\right)\widehat{\phi }_{\mathrm{}}(x)=0.$$ (2.6) Here $`\mathrm{\Delta }`$ is the d’Alembertian operator for the two-dimensional metric $`h_{ab}`$, and the induced potential $`V_{\mathrm{}}`$ is given by $$V_{\mathrm{}}=V+\frac{\mathrm{}(\mathrm{}+1)}{r^2}\mathrm{\Delta }\varphi +(\varphi )^2.$$ (2.7) The corresponding two-dimensional Green functions $`𝒢_{\mathrm{}}`$ satisfy $$_{\mathrm{}}𝒢_{\mathrm{}}(x,x^{})=\delta (x,x^{}),$$ (2.8) and are related to $`G`$ via $$G(X,X^{})=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{(2\mathrm{}+1)}{4\pi rr^{}}P_{\mathrm{}}(\mathrm{cos}\lambda )𝒢_{\mathrm{}}(x,x^{}),$$ (2.9) where $`P_{\mathrm{}}`$ is a Legendre polynomial and $`\lambda `$ is the angular separation of $`X`$, $`X^{}`$. Since the square of the field operator is given by the coincidence limit of the Green function, equation (2.9) implies that the four-dimensional $`\widehat{\mathrm{\Phi }}^2`$ can be obtained by solving the two-dimensional theory for $`\widehat{\phi }_{\mathrm{}}^2`$: $$\widehat{\mathrm{\Phi }}^2=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{(2\mathrm{}+1)}{4\pi r^2}\widehat{\phi }_{\mathrm{}}^2.$$ (2.10) The Green function, however, diverges in the coincidence limit, and must be renormalized to yield a finite $`\widehat{\mathrm{\Phi }}^2`$. Denoting the renormalized and divergent parts by the subscripts ‘ren’ and ‘div’ respectively, we have $`\widehat{\mathrm{\Phi }}^2_{\text{ren}}`$ $`=`$ $`\underset{X^{}X}{lim}G_{\text{ren}}(X,X^{})=\underset{X^{}X}{lim}\left[G(X,X^{})G_{\text{div}}(X,X^{})\right],`$ (2.11) $`\widehat{\phi }_{\mathrm{}}^2_{\text{ren}}`$ $`=`$ $`\underset{x^{}x}{lim}𝒢_{\mathrm{}\|\text{ren}}(x,x^{})=\underset{x^{}x}{lim}\left[𝒢_{\mathrm{}}(x,x^{})𝒢_{\mathrm{}\|\text{div}}(x,x^{})\right].`$ (2.12) While the bare quantities $`G`$ and $`𝒢_{\mathrm{}}`$ are related by the mode-decomposition relation (2.9), we shall find that the divergent parts $`G_{\text{div}}`$ and $`𝒢_{\mathrm{}\|\text{div}}`$ are not. As a result, the renormalized theories in two and four dimensions are related not by (2.9, 2.10) but rather by $$G_{\text{ren}}(X,X^{})=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{(2\mathrm{}+1)}{4\pi rr^{}}P_{\mathrm{}}(\mathrm{cos}\lambda )\left[𝒢_{\mathrm{}\|\text{ren}}(x,x^{})+\mathrm{\Delta }𝒢_{\mathrm{}}(x,x^{})\right],$$ (2.13) $$\widehat{\mathrm{\Phi }}^2_{\text{ren}}(X)=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{(2\mathrm{}+1)}{4\pi r^2}\left[\widehat{\phi }_{\mathrm{}}^2_{\text{ren}}(x)+\mathrm{\Delta }\widehat{\phi }_{\mathrm{}}^2(x)\right],$$ (2.14) where the anomalous terms are easily shown to be $`\mathrm{\Delta }\widehat{\phi }_{\mathrm{}}^2(x)`$ $`=`$ $`\underset{x^{}x}{lim}\mathrm{\Delta }𝒢_{\mathrm{}}(x,x^{})`$ (2.15) $``$ $`\underset{x^{}x}{lim}\left[𝒢_{\mathrm{}\|\text{div}}(x,x^{})2\pi rr^{}{\displaystyle _1^1}d(\mathrm{cos}\lambda )P_{\mathrm{}}(\mathrm{cos}\lambda )G_{\text{div}}(X,X^{})\right].`$ One can show that similar formulae hold for other renormalized quantities, such as the effective action $`W`$ and the stress tensor: $$W_{\text{ren}}=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}(2\mathrm{}+1)\left[𝒲_{\mathrm{}\|\text{ren}}+\mathrm{\Delta }𝒲_{\mathrm{}}\right];$$ (2.16) $$\widehat{T}_{\mu \nu }_{\text{ren}}=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{(2\mathrm{}+1)}{4\pi r^2}\left[\widehat{𝒯}_{\mu \nu }_{\mathrm{}\|\text{ren}}+\mathrm{\Delta }\widehat{𝒯}_{\mu \nu }_{\mathrm{}}\right].$$ (2.17) In each case the anomaly is the difference between the divergent subtraction terms for the dimensionally reduced theory and the mode-decomposed subtraction terms for the original four-dimensional theory. Equations (2.132.15) demonstrate that the renormalized value of a field quantity is generally not equal to the sum of the same renormalized quantities for the dimensionally-reduced theory. Rather, a quantity like $`\widehat{\mathrm{\Phi }}^2`$ can be obtained from dimensional reduction only if the contribution $`\widehat{\phi }_{\mathrm{}}^2`$ for each mode $`\mathrm{}`$ is modified by an extra anomalous term. This failure of dimensional reduction under renormalization is the dimensional-reduction anomaly. The remainder of this paper is devoted to explicit calculations of the anomalies in $`\widehat{\mathrm{\Phi }}^2`$ and $`W`$ for the important case of spherically symmetric geometries, described by (2.1). For further general discussion of the dimensional-reduction anomaly and the related multiplicative anomaly, the reader is referred to and . At this point some conventions on notation are in order. We shall need to be able to distinguish quantities like Green functions defined in different dimensions. ‘Ordinary’ letters such as $`G`$, $`W`$ are used for the original four-dimensional theory, while calligraphic letters such as $`𝒢`$, $`𝒲`$ refer to dimensionally reduced quantities. The anomalous difference between $`A`$, $`𝒜`$ is denoted $`\mathrm{\Delta }𝒜`$. All curvatures will be with respect to $`h`$ unless explicitly labelled otherwise; for example, $`R=R[h]`$ and $`{}_{}{}^{4}R=R[g]`$. As for differential operators, we shall understand $`\mathrm{}`$ to represent the d’Alembertian with respect to $`g`$, while $`\mathrm{\Delta }`$ is the d’Alembertian calculated using the metric $`h`$. Single covariant derivatives will be denoted by $``$; there will be no need to distinguish the metric used. For the dilaton $`\varphi `$ we shall understand $`\varphi _a`$, $`\varphi _{ab}`$, etc. to denote multiple two-dimensional covariant derivatives of $`\varphi `$ calculated using the metric $`h`$. ## 3 The Dimensional-Reduction Anomaly in Flat Space The simplest example of the dimensional reduction anomaly occurs in the spherical decomposition of a scalar field in flat space, and was originally considered in . We reproduce here the main formulae, as we shall require them for the generalization to curved space, and because some of the notation we use is different from that of . Let us assume that the potential $`V`$ vanishes inside the region of interest, and is spherically symmetric outside. In this case, the Green function for a given state is renormalized by subtracting the Green function for the Euclidean vacuum. In four dimensions the latter is $$G_{\text{div}}(X,X^{})=\frac{m}{4\pi ^2\sqrt{2\sigma }}K_1(m\sqrt{2\sigma }),$$ (3.1) where $`\sigma `$ is one-half the square of the geodesic distance between $`X`$ and $`X^{}`$, and $`K_1`$ is a modified Bessel function. In spherical coordinates $`X^\mu =(t,r,\theta ,\eta )`$<sup>1</sup><sup>1</sup>1We denote the azimuthal coordinate by $`\eta `$ rather than $`\varphi `$ for obvious reasons. the line element is $$ds^2=dt^2+dr^2+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\eta ^2\right),$$ (3.2) and $$2\sigma =(tt^{})^2+(rr^{})^2+2rr^{}\left(1\mathrm{cos}\lambda \right),$$ (3.3) where $`\lambda `$ is the angle between $`X`$ and $`X^{}`$, given by $$\mathrm{cos}\lambda =\mathrm{cos}\theta \mathrm{cos}\theta ^{}+\mathrm{sin}\theta \mathrm{sin}\theta ^{}\mathrm{cos}(\eta \eta ^{}).$$ (3.4) If before renormalizing we first decompose the field $`\widehat{\mathrm{\Phi }}`$ into spherical harmonics as in (2.4), we will be left with an effective field $`\widehat{\phi }_{\mathrm{}}`$ propagating on the two-dimensional space with line element $`ds^2=dt^2+dr^2`$, where $`t(\mathrm{},\mathrm{})`$ and $`r[0,\mathrm{})`$. In this case $`\sigma `$ is given by $`\frac{1}{2}[(tt^{})^2+(rr^{})^2]`$, and the theory is renormalized by subtracting the two-dimensional vacuum Green function $$𝒢_{\mathrm{}\|\text{div}}(x,x^{})=\frac{1}{2\pi }K_0(m\sqrt{(tt^{})^2+(rr^{})^2}).$$ (3.5) To compare the renormalization of the two- and four-dimensional theories and establish the existence of the dimensional-reduction anomaly, we decompose $`G_{\text{div}}`$ into spherical harmonics. Defining the mode decomposition by $$G_{\text{div}}(X,X^{})=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{(2\mathrm{}+1)}{4\pi rr^{}}P_{\mathrm{}}(\mathrm{cos}\lambda )G_{\text{div}\|\mathrm{}}(x,x^{})$$ (3.6) in accordance with (2.9), we have $$G_{\text{div}\|\mathrm{}}(x,x^{})=2\pi rr^{}_1^1d(\mathrm{cos}\lambda )P_{\mathrm{}}(\mathrm{cos}\lambda )G_{\text{div}}(X,X^{}).$$ (3.7) Inserting (3.1) into (3.7) and using the well-known integral representation for $`K_\nu `$, $$_0^{\mathrm{}}𝑑xx^{1\nu }\mathrm{exp}\left\{x\frac{\alpha ^2}{4x}\right\}=2\left(\frac{2}{\alpha }\right)^\nu K_\nu (\alpha ),$$ (3.8) the integral $$_1^1𝑑zP_{\mathrm{}}(z)\mathrm{e}^{p(1z)}=(1)^{\mathrm{}}\mathrm{e}^p\sqrt{\frac{2\pi }{p}}I_{\mathrm{}+1/2}(p),$$ (3.9) where $`I_{\mathrm{}+1/2}`$ is a modified Bessel function (see e.g. , vol.2, eq.2.17.5.2), and the representation $$I_{\mathrm{}+1/2}(p)=\frac{1}{\sqrt{2\pi p}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(\mathrm{}+k)!}{k!(\mathrm{}k)!}\frac{1}{(2p)^k}\left[(1)^k\mathrm{e}^p(1)^{\mathrm{}}e^p\right],$$ (3.10) (see, for example, 8.467 of ), we obtain $`G_{\text{div}\|\mathrm{}}(x,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\mathrm{}+k)!}{k!(\mathrm{}k)!}}[(1)^k{\displaystyle \frac{[(tt^{})^2+(rr^{})^2]^{k/2}}{(2mrr^{})^k}}K_k(m\sqrt{(tt^{})^2+(rr^{})^2})`$ (3.11) $`(1)^{\mathrm{}}{\displaystyle \frac{[(tt^{})^2+(r+r^{})^2]^{k/2}}{(2mrr^{})^k}}K_k(m\sqrt{(tt^{})^2+(r+r^{})^2})].`$ Let us compare this result for the mode-decomposed subtraction terms from four dimensions with the subtraction term for the two-dimensional theory, (3.5). While $`𝒢_{\mathrm{}\|\text{div}}`$ is the free-field Green function in two dimensions, it is not difficult to verify that $`G_{\text{div}\|\mathrm{}}`$ is the Green function for a field propagating in the centrifugal barrier potential<sup>2</sup><sup>2</sup>2See (2.7). It is easy to verify that $`\mathrm{\Delta }\varphi (\varphi )^2=0`$ for the line element (3.2). $$V_{\mathrm{}}=\frac{\mathrm{}(\mathrm{}+1)}{r^2}$$ (3.12) which obeys Dirichlet boundary conditions at $`r=0`$. Naive renormalization in two dimensions requires subtracting $`𝒢_{\mathrm{}\|\text{div}}`$. We see, however, that $`G_{\text{div}\|\mathrm{}}`$ is the quantity that should be subtracted to yield the correct results for the renormalized four-dimensional theory<sup>3</sup><sup>3</sup>3Examples of the correct procedure of renormalizing using the mode-decomposed subtraction terms from four dimensions in spherically symmetric spaces can be found in .. If one was to ignore the anomaly and calculate $`\widehat{\mathrm{\Phi }}^2_{\text{ren}}`$ using the two-dimensional $`\widehat{\phi }_{\mathrm{}}^2_{\text{ren}}`$ as in (2.10), one would obtain incorrect results, such as nonvanishing expectation values for the vacuum state. Instead, using (2.14, 2.15, 3.5, 3.11) one finds that the renormalized theories in two and four dimensions are related by $$\widehat{\mathrm{\Phi }}^2_{\text{ren}}(X)=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{(2\mathrm{}+1)}{4\pi r^2}\left[\widehat{\phi }_{\mathrm{}}^2_{\text{ren}}(x)+\mathrm{\Delta }\widehat{\phi }_{\mathrm{}}^2(x)\right],$$ where the anomaly is $`\mathrm{\Delta }\widehat{\phi }_{\mathrm{}}^2`$ $`=`$ $`\underset{x^{}x}{lim}\left[𝒢_{\mathrm{}\|\text{div}}G_{\text{div}\|\mathrm{}}\right]`$ (3.13) $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\mathrm{}+k)!}{(\mathrm{}k)!}}{\displaystyle \frac{1}{k}}{\displaystyle \frac{(1)^{k+1}}{(mr)^{2k}}}+{\displaystyle \frac{(1)^{\mathrm{}}}{2\pi }}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\mathrm{}+k)!}{k!(\mathrm{}k)!}}{\displaystyle \frac{K_k(2mr)}{(mr)^k}}.`$ For example, for the first two modes the anomalies are $`\mathrm{\Delta }\widehat{\phi }_{\mathrm{}=0}^2`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}K_0(2mr),`$ $`\mathrm{\Delta }\widehat{\phi }_{\mathrm{}=1}^2`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\left[{\displaystyle \frac{1}{(mr)^2}}K_0(2mr){\displaystyle \frac{2}{(mr)}}K_1(2mr)\right].`$ Note that the anomaly diverges logarithmically as $`r0`$, and vanishes as $`r\mathrm{}`$. ## 4 The Dimensional-Reduction Anomaly in $`\widehat{𝚽}^\mathrm{𝟐}`$ In the previous section we considered the dimensional-reduction anomaly in $`\widehat{\mathrm{\Phi }}^2`$ arising from the spherical decomposition of a scalar field in flat space. In this section we shall extend those calculations to general four-dimensional spherically symmetric spaces. Our system consists of a massive scalar field with arbitrary coupling to the four-dimensional scalar curvature, described by (2.2) with $$V=\xi ^4R.$$ (4.1) We assume that the spacetime of interest is given by the line element (2.1), which in standard spherical coordinates $`y^i=(\theta ,\eta )`$ becomes $$ds^2=h_{ab}(x^c)dx^adx^b+r^2(d\theta ^2+\mathrm{sin}^2\theta d\eta ^2),$$ (4.2) with $`r=\rho \mathrm{e}^{\varphi (x^c)}`$. As we saw in Section 2, under the dimensional reduction (2.4), the quantum field $`\widehat{\mathrm{\Phi }}`$ reduces to a collection of effective fields $`\widehat{\phi }_{\mathrm{}}`$ on the two-dimensional space (2.5) with metric $`h_{ab}`$, satisfying the field equation (2.6) with induced potential $$V_{\mathrm{}}=\xi ^4R+\frac{\mathrm{}(\mathrm{}+1)}{r^2}\mathrm{\Delta }\varphi +(\varphi )^2.$$ (4.3) We wish to compute the anomaly associated with renormalizing this dimensionally reduced theory versus (2.2). A standard approach to renormalization in curved space is via the heat kernel. For the system (2.2), the heat kernel $`K(X,X^{}\|s)`$ is a solution of the equation $$FK(X,X^{}\|s)=\frac{d}{ds}K(X,X^{}\|s)$$ (4.4) with boundary condition $`K(X,X^{}\|s=0)=\delta (X,X^{})`$. Once the heat kernel is known for a given state, both the Green function and the effective action may be obtained using $$G(X,X^{})=_0^{\mathrm{}}𝑑sK(X,X^{}\|s),$$ (4.5) $$W=\frac{1}{2}_0^{\mathrm{}}\frac{ds}{s}d^4X\sqrt{g}K(X,X\|s).$$ (4.6) Analogous formulae hold for the dimensionally reduced theory with operator $`_{\mathrm{}}`$, heat kernel $`𝒦_{\mathrm{}}(x,x^{}\|s)`$, Green function $`𝒢_{\mathrm{}}(x,x^{})`$ and effective action $`𝒲_{\mathrm{}}`$. The advantage of the heat kernel formulation is that the divergences in both the Green function and the effective action come from the $`s0`$ limit of the $`s`$ integral, and the small-$`s`$ behavior of the heat kernel is known for arbitrary curved spaces of any dimension. In particular, in four dimensions, $$K(X,X^{}\|s)=\frac{D^{\frac{1}{2}}}{(4\pi s)^2}\mathrm{exp}\left\{m^2s\frac{\sigma }{2s}\right\}\underset{i=0}{\overset{\mathrm{}}{}}a_i(X,X^{})s^i.$$ (4.7) Here again $`\sigma =\sigma (X,X^{})`$ is one-half of the square of the geodesic distance between the points $`X`$ and $`X^{}`$, while $`D=D(X,X^{})`$ is the Van Vleck determinant, $$D(X,X^{})=\frac{1}{\sqrt{g(X)}\sqrt{g(X^{})}}det\left[\frac{}{X^\mu }\frac{}{X_{}^{}{}_{}{}^{\nu }}\sigma (X,X^{})\right].$$ (4.8) The $`a_n`$ are the Schwinger-DeWitt coefficients for the operator $`F`$ of (2.2). In the coincidence limit $`X^{}X`$ the first few of these are $`a_0^{\mathrm{}\xi ^4R}`$ $`=`$ $`1,`$ (4.9) $`a_1^{\mathrm{}\xi ^4R}`$ $`=`$ $`\left({\displaystyle \frac{1}{6}}\xi \right)^4R,`$ (4.10) $`a_2^{\mathrm{}\xi ^4R}`$ $`=`$ $`{\displaystyle \frac{1}{180}}\left[{}_{}{}^{4}R_{\alpha \beta \gamma \delta }^{4}R^{\alpha \beta \gamma \delta }^4R_{\alpha \beta }^4R^{\alpha \beta }+\mathrm{}^4R\right]`$ (4.11) $`+{\displaystyle \frac{1}{6}}({\displaystyle \frac{1}{6}}\xi )\mathrm{}^4R+{\displaystyle \frac{1}{2}}({\displaystyle \frac{1}{6}}\xi )^2(^4R)^2.`$ For the two-dimensional operator $`_{\mathrm{}}`$ of (2.6) we only need the Schwinger-DeWitt expansion of the heat kernel in the coincidence limit. This is $$𝒦_{\mathrm{}}(x,x\|s)=\frac{1}{4\pi s}\mathrm{exp}\left\{m^2s\right\}\left[1+s\left(\frac{1}{6}RV_{\mathrm{}}\right)+\mathrm{}\right].$$ (4.12) Considering (4.5, 4.6), it is clear that in four dimensions the divergences in $`G`$ ($`W`$) arise from the first two<sup>4</sup><sup>4</sup>4 Comparing to (3.8), one sees that the integral representation used for $`G`$ in the previous section was just the heat kernel representation (4.5) with (4.7-4.10). (three) terms in the Schwinger-DeWitt expansion for $`K`$, while in two dimensions we need consider only the first term (first two terms) in $`𝒦_{\mathrm{}}`$. The anomaly in $`\widehat{\mathrm{\Phi }}^2`$ or $`W`$ can then be calculated by mode-decomposing the appropriate terms from $`K`$, comparing to the heat kernel $`𝒦_{\mathrm{}}`$ for the dimensionally-reduced theory, and finally integrating the difference over $`s`$ according to (4.5) or (4.6). Let us begin with the anomaly in $`\widehat{\mathrm{\Phi }}^2`$. The divergent part of the Green function in four dimensions is given by $$G_{\text{div}}(X,X^{})=_0^{\mathrm{}}𝑑sK_{\text{div}}(X,X^{}\|s),$$ (4.13) where $`K_{\text{div}}`$ consists of the first two terms of (4.7): $$K_{\text{div}}(X,X^{}\|s)=\frac{1}{(4\pi s)^2}\mathrm{exp}\left\{m^2s\frac{\sigma }{2s}\right\}\left[\mathrm{}_0^{\mathrm{}\xi ^4R}(X,X^{})+s\mathrm{}_1^{\mathrm{}\xi ^4R}(X,X^{})\right].$$ (4.14) Here we use the convenient notation $$\mathrm{}_n^{\mathrm{}\xi ^4R}(X,X^{})D^{\frac{1}{2}}(X,X^{})a_n^{\mathrm{}\xi ^4R}(X,X^{}).$$ (4.15) In principle, the anomaly in $`\widehat{\mathrm{\Phi }}^2`$ is straightforward to calculate. We mode-decompose $`K_{\text{div}}`$ in terms of Legendre polynomials in the usual manner: $`K_{\text{div}}(X,X^{}\|s)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2\mathrm{}+1)}{4\pi rr^{}}}P_{\mathrm{}}(\mathrm{cos}\lambda )K_{\text{div}\|\mathrm{}}(x,x^{}\|s);`$ (4.16) $`K_{\text{div}\|\mathrm{}}(x,x^{}\|s)`$ $`=`$ $`2\pi rr^{}{\displaystyle _1^1}d(\mathrm{cos}\lambda )P_{\mathrm{}}(\mathrm{cos}\lambda )K_{\text{div}}(X,X^{}\|s).`$ (4.17) The anomaly in $`\widehat{\mathrm{\Phi }}^2`$ is then just the coincidence limit of the difference between the subtraction terms in two dimensions and those mode-decomposed from four dimensions, integrated over $`s`$: $$\mathrm{\Delta }\widehat{\phi }_{\mathrm{}}^2(x)=_0^{\mathrm{}}𝑑s\left[𝒦_{\mathrm{}\|\text{div}}(x,x\|s)K_{\text{div}\|\mathrm{}}(x,x\|s)\right].$$ (4.18) We encounter a difficulty, however, when we try to perform the mode decomposition. For a general space, $`\sigma `$ and the $`a_n^{\mathrm{}\xi ^4R}`$ are known only for infinitesimal separations<sup>5</sup><sup>5</sup>5 In terms of momentum integrals, finite separations correspond to the low-frequency regime, where the renormalization terms are not fixed by the divergences in the theory. of $`X`$ and $`X^{}`$, while evaluation of the mode-decomposition integral (4.17) requires knowing $`\sigma `$ and the $`a_n^{\mathrm{}\xi ^4R}`$ for finite separations of $`X`$, $`X^{}`$ on the two-sphere. We proceed by determining an approximate $`K_{\text{div}}`$ for finite separation based on the following criteria: 1. Our approximate $`K_{\text{div}}`$ must reduce to the known value in the flat-space limit. 2. Our approximate $`K_{\text{div}}`$ must respect the periodicity of the two-spheres (i.e., it must be periodic in the angular separation $`\lambda `$ with period $`2\pi `$). In a previous case in which the mode decompositions were performed over noncompact spaces, the following procedure was found to work quite well. We take $`X=(x,y)`$ and $`X^{}=(x,y^{})`$; i.e., we split the points in the $`y`$-direction only. Using the well-known short-distance expansions obtained in , $`\sigma `$ and the $`\mathrm{}_n^{\mathrm{}\xi ^4R}`$ are expanded in powers of $`(yy^{})`$, which is equivalent to expanding in powers of the curvature. These expansions are then substituted into $`K_{\text{div}}`$ and, assuming small curvatures, truncated at first order in the curvature for $`\mathrm{\Delta }\widehat{\phi }_{\mathrm{}}^2`$ and at second order for $`\mathrm{\Delta }𝒲_{\mathrm{}}`$. The mode-decomposition integrals in can then be evaluated with relative ease. In the present case, the equivalent procedure is to expand $`\sigma `$ and the $`\mathrm{}_n^{\mathrm{}\xi ^4R}`$ in powers of $`\lambda ^2`$, which is easily done; see Appendix B. We also take into account the periodicity of the two-spheres by converting our expansions in $`\lambda ^2`$ into expansions in $`(1\mathrm{cos}\lambda )`$. Defining $`z=\mathrm{cos}\lambda `$, we have $$\lambda ^2=2(1z)+\frac{1}{3}(1z)^2+\frac{4}{45}(1z)^3+\mathrm{}.$$ (4.19) We then substitute (4.19) for each $`\lambda ^2`$, truncating at the lowest order in $`(1z)`$ which will yield the correct flat-space limit. This replacement of $`\lambda ^2`$ by a finite series in $`(1z)`$ means that our expansions are only modified for large angular separations, where the renormalization terms are inherently ambiguous. Our choice simply corresponds to a natural extension of the flat-space heat kernel which respects the periodicity of the two-spheres for large angular separations. For more details, see Appendix B. Using this procedure, one finds that to first order in the curvature $`2\sigma `$ $`=`$ $`2r^2(1z)+{\displaystyle \frac{r^2}{3}}[1r^2(\varphi )^2](1z)^2,`$ (4.20) $`\mathrm{}_0^{\mathrm{}\xi ^4R}`$ $`=`$ $`1+{\displaystyle \frac{1}{6}}^4R_{\theta \theta }(1z),`$ (4.21) $`\mathrm{}_1^{\mathrm{}\xi ^4R}`$ $`=`$ $`\left({\displaystyle \frac{1}{6}}\xi \right)^4R.`$ (4.22) Inserting these expansions into (4.7) yields our approximation for the ‘divergent’ part of the four-dimensional heat kernel, $`K_{\text{div}}(X,X^{}\|s)`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi s)^2}}\mathrm{exp}\{m^2s{\displaystyle \frac{r^2}{2s}}(1z)\}[1+s({\displaystyle \frac{1}{6}}\xi )^4R+{\displaystyle \frac{1}{6}}^4R_{\theta \theta }(1z)`$ (4.23) $`{\displaystyle \frac{r^2}{12s}}[1r^2(\varphi )^2](1z)^2].`$ The mode decomposition (4.17) of this $`K_{\text{div}}`$ then boils down to evaluating the integrals $$J_\mathrm{}np_1^1𝑑zP_{\mathrm{}}(z)\mathrm{e}^{p(1z)}(1z)^n,$$ (4.24) where $`pr^2/2s`$ is a dimensionless parameter and $`n`$ is an integer. The integrals for $`n0`$ can be obtained from the $`n=0`$ result (3.9, 3.10) used in the flat-space case by differentiating with respect to $`p`$, yielding $$J_\mathrm{}n=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(\mathrm{}+k)!}{k!(\mathrm{}k)!}\frac{1}{(2p)^k}\left[\frac{(1)^k2^n}{(2p)^n}\frac{(k+n)!}{k!}(1)^{\mathrm{}}e^{2p}\underset{\alpha =0}{\overset{n}{}}\frac{2^n}{(2p)^\alpha }\frac{(k+\alpha )!}{k!}\frac{n!}{\alpha !(n\alpha )!}\right].$$ (4.25) The mode-decomposed heat kernel subtraction terms for a general four-dimensional spherically symmetric spacetime are then $$K_{\text{div}\|\mathrm{}}(x,x\|s)=\frac{\mathrm{e}^{m^2s}}{4\pi s}\left[J_\mathrm{}0+s\left(\frac{1}{6}\xi \right)^4RJ_\mathrm{}0+\frac{1}{6}^4R_{\theta \theta }J_\mathrm{}1\frac{r^2}{12s}[1r^2(\varphi )^2]J_\mathrm{}2\right].$$ (4.26) The first term in (4.26) is the mode decomposition for flat space, while the other terms carry the contributions due to the curvature. Meanwhile, the various parts of the $`J_\mathrm{}n`$ fulfill several roles. First, the $`k0`$ terms in (4.25) are associated with the centrifugal potential $`\mathrm{}(\mathrm{}+1)/r^2`$ induced by the the mode decomposition. This potential is ignored in the renormalization in two dimensions, since only the first (potential-independent) term in the Schwinger-DeWitt expansion of the heat kernel contributes divergences to the two-dimensional Green function. Second, then terms in (4.25) proportional to $`\mathrm{e}^{2p}=\mathrm{e}^{r^2/s}`$ enforce a Dirichlet boundary condition at $`r=0`$, which is required if the four-dimensional subtraction term is to be finite there \[see (4.16)\]. These results are to be compared with the subtraction term in two dimensions, which consists of the first term of (4.12): $$𝒦_{\mathrm{}\|\text{div}}(x,x\|s)=\frac{\mathrm{e}^{m^2s}}{4\pi s}.$$ (4.27) In contrast to $`K_{\text{div}\|\mathrm{}}`$, $`𝒦_{\mathrm{}\|\text{div}}`$ is independent of both the position, the two-metric $`h_{ab}`$, and the mode number $`\mathrm{}`$. It matches just the first term in the $`k=0`$ contribution to the flat space part of $`K_{\text{div}\|\mathrm{}}`$. As an example, let us consider a quantum field in Schwarzschild space, for which $`{}_{}{}^{4}R=0`$, $`{}_{}{}^{4}R_{\mu \nu }^{}=0`$, and $$[1r^2(\varphi )^2]=\frac{2M}{r},$$ (4.28) where $`M`$ is the black-hole mass. Figure 1 shows plots of $`K_{\text{div}\|\mathrm{}=0}(x,x\|s)`$ for fixed $`s`$ and various values of $`M/\sqrt{s}`$. Note that large values of $`M`$ cause the mode-decomposed subtraction terms to become negative. The anomaly in $`\widehat{\mathrm{\Phi }}^2`$ can now be found by integrating the difference of $`𝒦_{\mathrm{}\|\text{div}}`$ and $`K_{\text{div}\|\mathrm{}}`$ as in (4.18). We find $`\mathrm{\Delta }\widehat{\phi }_{\mathrm{}}^2`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}[I[m^2sJ_\mathrm{}0]{\displaystyle \frac{1}{m^2}}({\displaystyle \frac{1}{6}}\xi )^4RI[(m^2s)^2J_\mathrm{}0]{\displaystyle \frac{1}{6}}^4R_{\theta \theta }I[m^2sJ_\mathrm{}1]`$ (4.29) $`+{\displaystyle \frac{(mr)^2}{12}}[1r^2(\varphi )^2]I[J_\mathrm{}2]],`$ where $`I[(m^2s)^tJ_\mathrm{}n]`$ $``$ $`{\displaystyle \underset{k=2tn}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\mathrm{}+k)!}{k!(\mathrm{}k)!}}\mathrm{\hspace{0.17em}2}^n\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{(1)^k}{(mr)^{2n+2k}}}{\displaystyle \frac{(k+n)!}{k!}}(k+t+n2)!\right]`$ (4.30) $`(1)^{\mathrm{}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\mathrm{}+k)!}{k!(\mathrm{}k)!}}\mathrm{\hspace{0.17em}2}^n\left[{\displaystyle \underset{\alpha =0}{\overset{n}{}}}{\displaystyle \frac{(k+\alpha )!}{k!}}{\displaystyle \frac{n!}{\alpha !(n\alpha )!}}{\displaystyle \frac{K_{k+t+\alpha 1}(2mr)}{(mr)^{kt+\alpha +1}}}\right].`$ (The $`I[(m^2s)^tJ_\mathrm{}n]`$ result from integrating terms of the form $`(m^2s)^{t2}J_\mathrm{}n`$ over $`s`$.) For example, in Schwarzschild space the anomaly for the $`\mathrm{}=0`$ mode is $$\mathrm{\Delta }\widehat{\phi }_{\mathrm{}=0}^2=\frac{1}{2\pi }K_0(2mr)+\frac{M}{3\pi r}\left[\frac{1}{(mr)^2}\frac{2}{mr}K_1(2mr)2K_0(2mr)mrK_1(2mr)\right].$$ (4.31) Plots of $`\mathrm{\Delta }\widehat{\phi }_{\mathrm{}=0}^2`$ for various values of $`mM`$ are shown in Figure 2. Note that the anomaly in $`\widehat{\mathrm{\Phi }}^2`$ generally diverges at any point $`x^a`$ such that $`r(x^a)=0`$, while for asymptotically flat spaces it vanishes as $`r\mathrm{}`$. ## 5 The Dimensional-Reduction Anomaly in the Effective Action In the previous section we calculated the dimensional-reduction anomaly in $`\widehat{\mathrm{\Phi }}^2`$ for a general four-dimensional spherically symmetric space. We now use the same procedure to determine the anomaly in the effective action, denoted by $`\mathrm{\Delta }𝒲_{\mathrm{}}`$ in (2.16). Functional differentiation of $`\mathrm{\Delta }𝒲_{\mathrm{}}`$ with respect to the metric $`h_{ab}`$ would then give the corresponding anomaly $`\mathrm{\Delta }\widehat{𝒯}_{\mu \nu }_{\mathrm{}}`$ in the stress tensor in (2.17). For the four-dimensional effective action (4.6) the divergent part of the heat kernel consists of the first three terms of (4.7): $$K_{\text{div}}(X,X^{}\|s)=\frac{1}{(4\pi s)^2}\mathrm{exp}\left\{m^2s\frac{\sigma }{2s}\right\}\left[\mathrm{}_0^{\mathrm{}\xi ^4R}+s\mathrm{}_1^{\mathrm{}\xi ^4R}+s^2\mathrm{}_2^{\mathrm{}\xi ^4R}\right].$$ (5.1) As in the previous section, we split the points $`X`$, $`X^{}`$ in the angular direction only. Then one can write $`2\sigma `$ $`=`$ $`2r^2\left[(1z)+u(1z)^2+v(1z)^3+\mathrm{}\right],`$ (5.2) $`\mathrm{}_n^{\mathrm{}\xi ^4R}`$ $`=`$ $`\mathrm{}_{n(0)}^{\mathrm{}\xi ^4R}+\mathrm{}_{n(1)}^{\mathrm{}\xi ^4R}(1z)+\mathrm{}_{n(2)}^{\mathrm{}\xi ^4R}(1z)^2+\mathrm{},`$ (5.3) where $`z=\mathrm{cos}\lambda `$. From the calculations for the anomaly in $`\widehat{\mathrm{\Phi }}^2`$ we have seen that $`u=\frac{1}{6}[1r^2(\varphi )^2]`$, $`\mathrm{}_{0(0)}^{\mathrm{}\xi ^4R}=1`$, $`\mathrm{}_{0(1)}^{\mathrm{}\xi ^4R}=\frac{1}{6}^4R_{\theta \theta }`$, and $`\mathrm{}_{1(0)}^{\mathrm{}\xi ^4R}=\left(\frac{1}{6}\xi \right)^4R`$. The other $`\mathrm{}_{n(k)}^{\mathrm{}\xi ^4R}`$ and $`v`$ are found in Appendix B. Inserting these expansions into (5.1) and truncating at second order in the curvature, we find $`K_{\text{div}}(X,X^{}\|s)`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi s)^2}}\mathrm{exp}\{m^2s{\displaystyle \frac{r^2}{2s}}(1z)\}[\{1+s\mathrm{}_{1(0)}^{\mathrm{}\xi ^4R}+s^2\mathrm{}_{2(0)}^{\mathrm{}\xi ^4R}\}`$ (5.4) $`+\left\{\mathrm{}_{0(1)}^{\mathrm{}\xi ^4R}+s\mathrm{}_{1(1)}^{\mathrm{}\xi ^4R}\right\}(1z)+\left\{\mathrm{}_{0(2)}^{\mathrm{}\xi ^4R}{\displaystyle \frac{r^2u}{2s}}{\displaystyle \frac{r^2u}{2}}\mathrm{}_{1(0)}^{\mathrm{}\xi ^4R}\right\}(1z)^2`$ $`+\{{\displaystyle \frac{r^2u}{2s}}\mathrm{}_{0(1)}^{\mathrm{}\xi ^4R}{\displaystyle \frac{r^2v}{2s}}\}(1z)^3+{\displaystyle \frac{r^4u^2}{8s^2}}(1z)^4].`$ The decomposition of the heat kernel subtraction terms (5.4) is done in the same manner as in the previous section. Employing the definition (4.17) of the spherical decomposition and using the functions $`J_\mathrm{}n`$ of (4.24, 4.25), we obtain $`K_{\text{div}\|\mathrm{}}(x,x\|s)`$ $`=`$ $`\underset{x^{}x}{lim}2\pi r^2{\displaystyle _1^1}d(\mathrm{cos}\lambda )P_{\mathrm{}}(\mathrm{cos}\lambda )K_{\text{div}}(X,X^{}\|s)`$ (5.5) $`=`$ $`{\displaystyle \frac{\mathrm{e}^{m^2s}}{4\pi s}}[\{1+s\mathrm{}_{1(0)}^{\mathrm{}\xi ^4R}+s^2\mathrm{}_{2(0)}^{\mathrm{}\xi ^4R}\}J_\mathrm{}0+\{\mathrm{}_{0(1)}^{\mathrm{}\xi ^4R}+s\mathrm{}_{1(1)}^{\mathrm{}\xi ^4R}\}J_\mathrm{}1`$ $`+\left\{\mathrm{}_{0(2)}^{\mathrm{}\xi ^4R}{\displaystyle \frac{r^2u}{2s}}{\displaystyle \frac{r^2u}{2}}\mathrm{}_{1(0)}^{\mathrm{}\xi ^4R}\right\}J_\mathrm{}2`$ $`+\{{\displaystyle \frac{r^2v}{2s}}{\displaystyle \frac{r^2u}{2s}}\mathrm{}_{0(1)}^{\mathrm{}\xi ^4R}\}J_\mathrm{}3+{\displaystyle \frac{r^4u^2}{8s^2}}J_\mathrm{}4].`$ Meanwhile, the divergences in the effective action for the two-dimensional theory (2.6) arise from the first two terms of (4.12): $$𝒦_{\mathrm{}\|\text{div}}(x,x\|s)=\frac{\mathrm{e}^{m^2s}}{4\pi s}\left[1+s\left(\frac{1}{6}RV_{\mathrm{}}\right)\right].$$ (5.6) The anomaly in the effective action is then found by integrating the difference of (5.5, 5.6) as in (4.6): $`\mathrm{\Delta }𝒲_{\mathrm{}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^2x\sqrt{h}_0^{\mathrm{}}\frac{ds}{s}\left[𝒦_{\mathrm{}\|\text{div}}(x,x\|s)K_{\text{div}\|\mathrm{}}(x,x\|s)\right]}`$ (5.7) $`=`$ $`{\displaystyle \frac{m^2}{4\pi }}{\displaystyle }d^2x\sqrt{h}[I[J_\mathrm{}0]+{\displaystyle \frac{1}{m^2}}\mathrm{}_{1(0)}^{\mathrm{}\xi ^4R}I[m^2sJ_\mathrm{}0]+{\displaystyle \frac{1}{m^4}}\mathrm{}_{2(0)}^{\mathrm{}\xi ^4R}I[m^4s^2J_\mathrm{}0]+\mathrm{}_{0(1)}^{\mathrm{}\xi ^4R}I[J_\mathrm{}1]`$ $`+{\displaystyle \frac{1}{m^2}}\mathrm{}_{1(1)}^{\mathrm{}\xi ^4R}I[m^2sJ_\mathrm{}1]+(\mathrm{}_{0(2)}^{\mathrm{}\xi ^4R}{\displaystyle \frac{r^2u}{2}}\mathrm{}_{1(0)}^{\mathrm{}\xi ^4R})I[J_\mathrm{}2]{\displaystyle \frac{(mr)^2u}{2}}I[{\displaystyle \frac{1}{m^2s}}J_\mathrm{}2]`$ $`{\displaystyle \frac{(mr)^2}{2}}(u\mathrm{}_{0(1)}^{\mathrm{}\xi ^4R}+v)I[{\displaystyle \frac{1}{m^2s}}J_\mathrm{}3]+{\displaystyle \frac{(mr)^4u^2}{8}}I[{\displaystyle \frac{1}{m^4s^2}}J_\mathrm{}4]].`$ The $`I[(m^2s)^tJ_\mathrm{}n]`$ are given by (4.30). Using (5.7) and the values of $`u`$, $`v`$, and the $`\mathrm{}_{n(k)}^{\mathrm{}\xi ^4R}`$ given in Appendix B, one can compute the anomalous contribution to the stress tensor. ## 6 Conclusions In a $`D`$-dimensional spacetime which can be foliated by $`n`$-dimensional homogeneous subspaces, a field can be decomposed in terms of modes on the subspaces. This effectively converts the system from a single quantum field in $`D`$ dimensions to a collection of fields in $`(Dn)`$ dimensions. Quantities of interest for the original theory, such as the expectation value of the square of the field operator and the effective action, can then be written as sums of the corresponding objects from the dimensionally reduced theories. This relationship breaks down under renormalization, however, so that renormalized expectation values can be obtained by summing their lower-dimensional counterparts only if the contribution for each mode is modified by adding an anomalous contribution. This effect is the dimensional reduction anomaly. We have explicitly calculated the anomalous contributions to the expectation value of the square of the field operator and the effective action for the case of a massive scalar field propagating in a general four-dimensional spherically symmetric space. We have seen that the anomaly arises from several sources. One is the Dirichlet boundary condition imposed at $`r=0`$ due to the change in topology inherent in the spherical reduction of the spacetime. Other contributions are more local in nature, arising from the dimension-dependent contributions of the curvature and field potential to divergences. The resulting anomaly terms are constructed from the curvature, the dilaton field, and their covariant derivatives, and cannot be eliminated by further finite renormalization. The anomalies calculated in this paper may be of importance to recent attempts to calculate the stress tensor and Hawking radiation in black-hole spacetimes using quantum fields in two dimensions . These attempts are based on the dimensional reduction of a massless minimally-coupled quantum field in a Schwarzschild spacetime, followed by renormalization in two dimensions. We have seen, however, that the contributions of the dimensionally reduced fields should be modified by adding the corresponding anomaly term. We intend to return to a discussion of this interesting topic in a future publication. Acknowledgments: The author would like to thank Valeri Frolov and Andrei Zelnikov for useful discussions. This work was supported by the Natural Sciences and Engineering Research Council of Canada. ## Appendix A Spherical Decomposition of Curvatures Consider a line element of the form $$ds^2=g_{\mu \nu }dX^\mu dX^\nu =h_{ab}dx^adx^b+\rho ^2\mathrm{e}^{2\varphi }\omega _{ij}dy^idy^j,$$ (A.1) where $`h_{ab}=h_{ab}(x^c)`$ is an arbitrary two-dimensional metric and $`\omega _{ij}=\omega _{ij}(y^k)`$ is the metric of a two-sphere. The dilaton $`\varphi `$ is a function of the $`x^a`$ only, and $`\rho `$ is a constant with the dimensions of length. The radius of a two-sphere of fixed $`x^a`$ is $`r=\rho \mathrm{e}^\varphi `$. We wish to decompose our field theory in terms of modes on the two-sphere. This requires rewriting four-dimensional geometric quantities like the curvatures in terms of the corresponding curvatures for the metric $`h`$. Our notational conventions are as follows: four-dimensional covariant derivatives are denoted by $`()_{;a}`$, while $`\mathrm{}`$ is understood to represent the d’Alembertian with respect to $`g`$. Meanwhile, $``$, $`()_{\|a}`$ and $`\mathrm{\Delta }`$ are the two-dimensional covariant derivatives and d’Alembertian calculated using the metric $`h_{ab}`$. For the dilaton $`\varphi `$ we shall understand $`\varphi _a`$, $`\varphi _{ab}`$, etc. to denote multiple two-dimensional covariant derivatives of $`\varphi `$. For example, the four-dimensional d’Alembertian of an angle-independent scalar $`S`$ decomposes to $$\mathrm{}S=\mathrm{\Delta }S2\varphi S.$$ (A.2) In particular, $$\mathrm{}\varphi =\mathrm{\Delta }\varphi 2(\varphi )^2.$$ (A.3) For the given line element, the nonvanishing Christoffel symbols are $`{}_{}{}^{4}\mathrm{\Gamma }_{bc}^{a}[g]`$ $`=`$ $`{}_{}{}^{2}\mathrm{\Gamma }_{bc}^{a}[h],`$ (A.4) $`{}_{}{}^{4}\mathrm{\Gamma }_{ij}^{a}[g]`$ $`=`$ $`\varphi ^ag_{ij},`$ (A.5) $`{}_{}{}^{4}\mathrm{\Gamma }_{ja}^{i}[g]`$ $`=`$ $`\varphi _a\delta _j^i,`$ (A.6) $`{}_{}{}^{4}\mathrm{\Gamma }_{ij}^{k}[g]`$ $`=`$ $`{}_{}{}^{2}\mathrm{\Gamma }_{ij}^{k}[\omega ].`$ (A.7) Selecting coordinates $`(\theta ,\eta )`$ on the two-spheres, where $$d\omega _{ij}y^iy^j=d\theta ^2+\mathrm{sin}^2\theta d\eta ^2,$$ (A.8) one finds $${}_{}{}^{4}\mathrm{\Gamma }_{\eta \eta }^{\theta }[g]=\mathrm{sin}\theta \mathrm{cos}\theta ,^4\mathrm{\Gamma }_{\eta \theta }^\eta [g]=\frac{\mathrm{cos}\theta }{\mathrm{sin}\theta }.$$ (A.9) For convenience, we define the following commonly-occurring functions of the dilaton field: $`A`$ $`=`$ $`1r^2(\varphi )^2,`$ (A.10) $`B`$ $`=`$ $`\mathrm{\Delta }\varphi 2(\varphi )^2,`$ (A.11) $`T_{ab}`$ $`=`$ $`\varphi _{ab}\varphi _a\varphi _b,`$ (A.12) $`T`$ $`=`$ $`h^{ab}T_{ab}=\mathrm{\Delta }\varphi (\varphi )^2.`$ (A.13) Since the two-sphere metric has constant curvature $`{}_{}{}^{2}R[\omega ]=2`$, explicit reference to it may be dropped. Henceforth we shall assume all curvatures to be with respect to the two-dimensional metric $`h_{ab}`$ unless explicitly labelled otherwise. Using this notation, one can show that the only nonvanishing components of the four-dimensional curvatures are $`{}_{}{}^{4}R_{abcd}^{}[g]`$ $`=`$ $`{\displaystyle \frac{1}{2}}R(h_{ac}h_{bd}h_{ad}h_{bc}),`$ (A.14) $`{}_{}{}^{4}R_{aibj}^{}[g]`$ $`=`$ $`g_{ij}T_{ab},`$ (A.15) $`{}_{}{}^{4}R_{ijkm}^{}[g]`$ $`=`$ $`{\displaystyle \frac{A}{r^2}}(g_{ik}g_{jm}g_{im}g_{jk}),`$ (A.16) $`{}_{}{}^{4}R_{ab}^{}[g]`$ $`=`$ $`{\displaystyle \frac{1}{2}}Rh_{ab}+2T_{ab},`$ (A.17) $`{}_{}{}^{4}R_{ij}^{}[g]`$ $`=`$ $`g_{ij}\left[{\displaystyle \frac{1}{r^2}}+B\right],`$ (A.18) $`{}_{}{}^{4}R[g]`$ $`=`$ $`R+4\mathrm{\Delta }\varphi 6(\varphi )^2+{\displaystyle \frac{2}{r^2}},`$ (A.19) while the only nonvanishing $`{}_{}{}^{4}R_{\alpha \beta ;\gamma }^{}`$ are $`{}_{}{}^{4}R_{ab;c}^{}[g]`$ $`=`$ $`{\displaystyle \frac{1}{2}}h_{ab}R_{\|c}+2T_{ab\|c},`$ (A.20) $`{}_{}{}^{4}R_{am;n}^{}[g]`$ $`=`$ $`g_{mn}\left[\left({\displaystyle \frac{1}{2}}R+{\displaystyle \frac{1}{r^2}}+B\right)\varphi _a2T_{ab}\varphi ^b\right],`$ (A.21) $`{}_{}{}^{4}R_{mn;a}^{}[g]`$ $`=`$ $`g_{mn}\left({\displaystyle \frac{1}{r^2}}+B\right)_{\|a}.`$ (A.22) Also, $`{}_{}{}^{4}R_{mn;ab}^{}[g]`$ $`=`$ $`g_{mn}\left({\displaystyle \frac{1}{r^2}}+B\right)_{\|ab},`$ (A.23) $`{}_{}{}^{4}R_{mn;jk}^{}[g]`$ $`=`$ $`\left(g_{km}g_{nj}+g_{kn}g_{mj}\right)\left[\left({\displaystyle \frac{1}{2}}R+{\displaystyle \frac{1}{r^2}}+B\right)(\varphi )^22T_{ab}\varphi ^a\varphi ^b\right]`$ (A.24) $`g_{jk}g_{mn}\left({\displaystyle \frac{1}{r^2}}+B\right)_{\|a}\varphi ^a,`$ $`\mathrm{}^4R_{mn}[g]`$ $`=`$ $`g_{mn}\{[\mathrm{\Delta }2\varphi ]({\displaystyle \frac{1}{r^2}}+B)+R(\varphi )^22({\displaystyle \frac{1}{r^2}}+B)(\varphi )^2`$ (A.25) $`+4T_{ab}\varphi ^a\varphi ^b\},`$ $`{}_{}{}^{4}R_{;a}^{}`$ $`=`$ $`\left(R+4\mathrm{\Delta }\varphi 6(\varphi )^2+{\displaystyle \frac{2}{r^2}}\right)_{\|a},`$ (A.26) $`{}_{}{}^{4}R_{;ab}^{}`$ $`=`$ $`\left(R+4\mathrm{\Delta }\varphi 6(\varphi )^2+{\displaystyle \frac{2}{r^2}}\right)_{\|ab},`$ (A.27) $`{}_{}{}^{4}R_{;mn}^{}`$ $`=`$ $`g_{mn}\left(R+4\mathrm{\Delta }\varphi 6(\varphi )^2+{\displaystyle \frac{2}{r^2}}\right)_{\|a}\varphi ^a,`$ (A.28) $`\mathrm{}^4R`$ $`=`$ $`[\mathrm{\Delta }2\varphi ]\left(R+4\mathrm{\Delta }\varphi 6(\varphi )^2+{\displaystyle \frac{2}{r^2}}\right).`$ (A.29) ## Appendix B Point Splitting It will be necessary to write short-distance expansions for $`\sigma `$ and the $`D^{\frac{1}{2}}a_n`$ for $`X`$ and $`X^{}`$ separated along the two-spheres. We follow a method similar to that developed in . Without loss of generality we take the points to be split in the $`\theta `$ direction only, with angular separation $`\lambda =\theta \theta ^{}`$. Our procedure will be to calculate the desired quantities first as expansions in powers of $`\lambda ^2`$, and then to convert them to expansions in powers of $`(1\mathrm{cos}\lambda )`$ for use in the mode-decomposition calculations. We take as our ansatz for the geodetic interval $`\sigma `$ $$2\sigma (x,y;x^{},y^{})=(\stackrel{~}{r}\lambda )^2+U(\stackrel{~}{x})(\stackrel{~}{r}\lambda )^4+V(\stackrel{~}{x})(\stackrel{~}{r}\lambda )^6+\mathrm{},$$ (B.1) where $`\stackrel{~}{x}\frac{1}{2}(x+x^{})`$. Taking the derivative of $`\stackrel{~}{\sigma }`$ with respect to each of the coordinates and requiring $`\sigma =\frac{1}{2}g^{\alpha \beta }\sigma _\alpha \sigma _\beta `$ in the coincidence limit, one can show that $`U(x)`$ $`=`$ $`{\displaystyle \frac{1}{12}}(\varphi )^2,`$ (B.2) $`V(x)`$ $`=`$ $`{\displaystyle \frac{1}{90}}(\varphi )^4{\displaystyle \frac{1}{120}}\varphi ^a\varphi ^b\varphi _{ab},`$ (B.3) and $`(\sigma ^\theta )^2`$ $`=`$ $`\lambda ^2\left[1{\displaystyle \frac{1}{3}}r^2(\varphi )^2\lambda ^2+r^4\left({\displaystyle \frac{17}{180}}(\varphi )^4{\displaystyle \frac{1}{20}}\varphi ^a\varphi ^b\varphi _{ab}\right)\lambda ^4+\mathrm{}\right],`$ (B.4) $`\sigma ^\eta `$ $`=`$ $`0,`$ (B.5) $`\sigma ^a`$ $`=`$ $`{\displaystyle \frac{1}{2}}\varphi ^a(r\lambda )^2+\left[{\displaystyle \frac{1}{24}}\varphi ^{ab}\varphi _b+{\displaystyle \frac{1}{12}}(\varphi )^2\varphi ^a\right](r\lambda )^4+\mathrm{}.`$ (B.6) The expansion (B.1) for $`\sigma `$ can be converted into one in terms of $`(1\mathrm{cos}\lambda )`$ using $$\lambda ^2=2(1z)+\frac{1}{3}(1z)^2+\frac{4}{45}(1z)^3+\mathrm{},$$ (B.7) where $`z\mathrm{cos}\lambda `$. Defining the functions $`u(x)`$, $`v(x)`$ by $$2\sigma (x,y;x,y^{})=2r^2\left[(1z)+u(x)(1z)^2+v(x)(1z)^3+\mathrm{}\right],$$ (B.8) we obtain $`u(x)`$ $`=`$ $`{\displaystyle \frac{1}{6}}[1r^2(\varphi )^2],`$ (B.9) $`v(x)`$ $`=`$ $`{\displaystyle \frac{2}{45}}\left[1{\displaystyle \frac{5}{4}}r^2(\varphi )^2+r^4(\varphi )^4{\displaystyle \frac{3}{8}}r^4\varphi [(\varphi )^2]\right].`$ (B.10) Combining (B.4B.7) with the results of Appendix A and the short-distance expansions of , one can derive expansions for the $`D^{\frac{1}{2}}a_n`$ in powers of $`(1z)`$. Writing $$D^{\frac{1}{2}}a_n^{\mathrm{}\xi ^4R}=\mathrm{}_n^{\mathrm{}\xi ^4R}=\mathrm{}_{n(0)}^{\mathrm{}\xi ^4R}+\mathrm{}_{n(1)}^{\mathrm{}\xi ^4R}(1z)+\mathrm{}_{n(2)}^{\mathrm{}\xi ^4R}(1z)^2+\mathrm{},$$ (B.11) one can show that $`\mathrm{}_{0(0)}^{\mathrm{}\xi ^4R}`$ $`=`$ $`1,`$ (B.12) $`\mathrm{}_{0(1)}^{\mathrm{}\xi ^4R}`$ $`=`$ $`{\displaystyle \frac{1}{6}}(1+r^2B),`$ (B.13) $`\mathrm{}_{0(2)}^{\mathrm{}\xi ^4R}`$ $`=`$ $`{\displaystyle \frac{1}{90}}A^2+{\displaystyle \frac{1}{72}}(1+r^2B)^2+{\displaystyle \frac{1}{36}}(1+r^2B)(14r^2(\varphi )^2)`$ (B.14) $`+{\displaystyle \frac{r^4}{180}}[{\displaystyle \frac{3}{2}}R(\varphi )^2+6T_{ab}\varphi ^a\varphi ^b+2T_{ab}T^{ab}+12({\displaystyle \frac{1}{r^2}}+B)(\varphi )^2`$ $`+6({\displaystyle \frac{1}{r^2}}+B)_{\|a}\varphi ^a],`$ $`\mathrm{}_{1(0)}^{\mathrm{}\xi ^4R}`$ $`=`$ $`\left({\displaystyle \frac{1}{6}}\xi \right)\left(R+4\mathrm{\Delta }\varphi 6(\varphi )^2+{\displaystyle \frac{2}{r^2}}\right),`$ (B.15) $`\mathrm{}_{1(1)}^{\mathrm{}\xi ^4R}`$ $`=`$ $`{\displaystyle \frac{1}{6}}\left({\displaystyle \frac{1}{6}}\xi \right)\left[(1+r^2B)+r^2\varphi \right]\left(R+4\mathrm{\Delta }\varphi 6(\varphi )^2+{\displaystyle \frac{2}{r^2}}\right)`$ (B.16) $`+{\displaystyle \frac{r^2}{180}}[RT+3R(\varphi )^2+8T_{ab}T^{ab}+12T_{ab}\varphi ^a\varphi ^b+(R+4\mathrm{\Delta }\varphi 6(\varphi )^2+{\displaystyle \frac{2}{r^2}})_{\|a}\varphi ^a`$ $`+3[\mathrm{\Delta }2\varphi ]({\displaystyle \frac{1}{r^2}}+B)6({\displaystyle \frac{1}{r^2}}+B)(\varphi )^2]`$ $`+{\displaystyle \frac{1}{90r^2}}\left(2A^2+A(1+r^2B)2(1+r^2B)^2\right),`$ $`\mathrm{}_{2(0)}^{\mathrm{}\xi ^4R}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{6}}\xi \right)^2\left(R+4\mathrm{\Delta }\varphi 6(\varphi )^2+{\displaystyle \frac{2}{r^2}}\right)^2`$ (B.17) $`+{\displaystyle \frac{1}{6}}\left({\displaystyle \frac{1}{6}}\xi \right)\left[\mathrm{\Delta }2\varphi \right]\left(R+4\mathrm{\Delta }\varphi 6(\varphi )^2+{\displaystyle \frac{2}{r^2}}\right)`$ $`+{\displaystyle \frac{1}{180}}[[\mathrm{\Delta }2\varphi ](R+4\mathrm{\Delta }\varphi 6(\varphi )^2+{\displaystyle \frac{2}{r^2}})`$ $`+{\displaystyle \frac{1}{2}}R^22RT+4T_{ab}T^{ab}+{\displaystyle \frac{4}{r^4}}A^2{\displaystyle \frac{2}{r^4}}(1+r^2B)^2].`$ It is easily verified that for flat spacetime each of the $`\mathrm{}_{n(k)}^{\mathrm{}\xi ^4R}`$ vanishes, except for $`\mathrm{}_{0(0)}^{\mathrm{}\xi ^4R}`$.
warning/0003/cond-mat0003161.html	ar5iv	text	# Dynamic Approach to Weak First Order Phase Transitions ## Abstract A short-time dynamic approach to weak first order phase transitions is proposed. Taking the 2-dimensional Potts models as examples, from short-time behaviour of non-equilibrium relaxational processes starting from high temperature and zero temperature states,x pseudo critical points $`K^{}`$ and $`K^{}`$ are determined. A clear difference of the values for $`K^{}`$ and $`K^{}`$ distinguishes a weak first order transition from a second order one. At the pseudo critical points, pseudo critical exponents can be estimated. In recent years, much progress has been achieved in non-equilibrium critical dynamics. For example, in a dynamic process in which a system initially at a high temperature or a zero temperature state, is suddenly quenched to the critical temperature or nearby and then evolves dynamically, short-time universal scaling behaviour has been found jan89 ; hus89 . This phenomenon is rather fundamental. It exists not only in stochastic dynamics described by Langevin equations jan89 ; oer93 or Monte Carlo algorithms hus89 ; hum91 ; sta92 ; li94 ; gra95 ; sch95 , but also in deterministic dynamics described by fundamental microscopic equations of motion zhe99 . More interestingly, based on the short-time scaling form, it is possible to determine not only dynamic exponents but also static exponents as well as the critical temperature li95 ; luo98 . Since the measurements are carried out in the short-time regime, one does not suffer from critical slowing down. Compared with non-local cluster algorithms, the short-time dynamic approach does study properties of the original local dynamics and also applies to systems with quenched randomness. For a review, see Ref. zhe98 . Naturally, it is interesting and attractive to explore possible applications of short-time dynamics to first order phase transitions. Especially, due to large correlation lengths and small discontinuities, a weak first order transition presents quite similar behaviour as a second order one. It has long been challenging how to distinguish one from the other. Furthermore, slowing down in Monte Carlo simulations at first order transitions is even more severe than at second order ones. Non-local cluster algorithms also do not show much more efficiency. In numerical simulations at first order transitions in equilibrium, to locate the transition point one usually searches for the maximums of the specific heat, susceptibility, or a Binder cumulant constructed from energy bin92 . For a system with lattice size $`L`$, these maximums deviate from the real transition point by a power law $`1/L^d`$. To remove this power law deviation, special techniques have been introduced bor92 . With these techniques, first order transition points can be determined rather accurately from moderate lattice sizes, even for weak first order transitions. To distinguish a first order transition from a second order one, naively one may explore a signal for discontinuity of the order parameter by increasing the lattice sizes. Refined methods are typically based on the finite size scaling of the specific heat, susceptibility, order parameter, Binder cumulant of energy, or the transition point, e.g. see Refs. bin92 ; pri90 ; bin97 ; fer88 ; fer89 ; oli95 ; jan95 ; fer98 . However, when a first order transition is very weak, it becomes subtle. The lattice sizes one reaches in simulations hardly feel the difference between very large correlation lengths in weak first order transitions and divergent ones in second order transitions. The double peak structure of the energy distribution together with the finite size scaling shows its merit in this respect lee90 ; lee91 ; jan00 , but further efficient methods are still desired. In this letter, we propose a short-time dynamic approach to weak first order transitions. The idea is inspired by the existence of two pseudo critical points $`K^{}`$ and $`K^{}`$ near the weak first order transition point $`K_c`$ with $`K^{}<K_c<K^{}`$ gen75 ; fer92 . In equilibrium, numerical measurements of $`K^{}`$ and $`K^{}`$ are not easy since they are induced by metastable states. However, in short-time dynamics $`K^{}`$ and $`K^{}`$ can be determined rather accurately from two dynamic processes starting from high temperature and zero temperature states. In second order transitions, $`K^{}`$ and $`K^{}`$ overlap with the transition point $`K_c`$. Therefore, difference of $`K^{}`$ and $`K^{}`$ gives a criterion for a weak first order transition. As examples, we investigate the two-dimensional $`q`$-state Potts models. The transition point is exactly known at $`K_c=\mathrm{ln}(1+\sqrt{q})`$. The phase transition is second order for $`q4`$ and becomes first order for $`q5`$. For small $`q`$, the first order transitions are weak. Especially, for $`q=5`$ the transition is so weak that with standard methods one hardly sees a difference from a second order one. In second order transitions, it has been shown that at the critical point, short-time behaviour of physical observables is a power law in dynamic processes starting from both a random and an ordered state. Away from the critical point, the power law behaviour is modified by a scaling function zhe98 . We will demonstrate that it is different for first order transitions. An approximate power law behaviour will be observed only at the pseudo critical points $`K^{}`$ and $`K^{}`$ in proper dynamic processes. We begin our investigation by determining $`K^{}`$ for the 7-state Potts model. For this purpose, we consider a dynamic process in which the system initially in a random state, is suddenly quenched to $`K_c`$ or above, then evolves dynamically. We have performed simulations with the heat-bath algorithm. Lattice sizes are $`L=140`$ and $`280`$ and maximum updating times are $`t_{max}=2000`$ and $`6000`$ respectively. Total samples for averaging are $`4600`$ and errors are simply estimated by dividing the data into four subsamples. In Fig. 1(a), the second moment $`M^{(2)}(t)`$ with $`L=280`$ is displayed for $`K=1.293562`$ ($`K_c`$), $`1.294210`$ and $`1.294857`$ on a log-log scale. Apparently, at $`K_c`$ the curve bends downwards and does not show a power law behaviour due to the random initial state and the finite spatial correlation length in equilibrium. Actually, this already indicates that the transition is first order if we assume $`K_c`$ is known. What is interesting here is that at a slightly bigger $`K`$, which we denote by $`K^{}`$, one observes an approximate power law behaviour. The weaker the transition is, the cleaner the power law behaviour will be. When $`K`$ becomes bigger than $`K^{}`$, the curve bends upwards. Therefore, $`K^{}`$ looks like a critical point zhe98 . We can not prove that our $`K^{}`$ is the same as the pseudo critical point $`K^{}`$ defined in equilibrium, but we strongly believe so. In equilibrium, $`K^{}`$ is defined as a point at which the system presents approximate scaling behaviour similar to that at a critical point gen75 ; fer92 . In our short-time dynamic approach, practically we locate the pseudo critical point $`K^{}`$ by interpolating $`M^{(2)}(t)`$ among the three simulated $`K`$’s and searching for the best power law behaviour luo98 ; zhe98 . In short-time critical dynamics, it has been intensively discussed that universal behaviour emerges only after a time scale $`t_{mic}`$ which is large enough in microscopic sense. If a Monte Carlo time step (a sweep over all spins on the lattice) is considered to be a microscopic time unit, $`t_{mic}`$ is typically $`10`$ to some hundred time steps zhe98 . Similarly, in first order transitions, physical behaviour at macroscopic level is presented also only after $`t_{mic}`$. In the upper part of Fig. 2, $`K^{}`$ obtained with data in a time interval $`[t,t_{max}]`$ is shown. The results are stable and $`K^{}`$ is clearly above $`K_c`$. The final value for $`K^{}`$ is estimated to be $`K_{7s}^{}=1.293854(29)`$. This is consistent with the value $`K^{}=1.2945(9)`$ given in Ref. fer92 . However, the latter can hardly distinguish $`K^{}`$ from $`K_c`$ within the error. To determine $`K^{}`$, we study a dynamic process in which the system initially in an ordered state, is quenched to $`K_c`$ or below, and evolves dynamically. Here we have performed extra simulations for $`L=560`$, up to $`t_{max}=6000`$. Total samples for $`L=140`$, $`280`$ and $`560`$ are $`7000`$, $`1500`$ and $`135`$ respectively. In Fig. 1(b) the magnetisation with $`L=280`$ is plotted for $`K=1.2929`$, $`1.2930`$ and $`1.2931`$. The curve for $`K_c=1.293562`$ (not in the figure) is much above that for $`1.2931`$ and very far from power law behaviour. This is again a signal for a first order transition. The reason is clear. For first order transitions, with an ordered initial state the system will evolve to the ordered phase at $`K_c`$. However, at the pseudo critical point $`K^{}`$ we will observe approximate power law behaviour. Searching for a curve with the best power law behaviour from the three curves in Fig. 1(b), we determine the pseudo critical point $`K^{}`$. The results are presented in the lower part of Fig. 2. The values are clearly below $`K_c`$. Another interesting observable is the Binder cumulant $`U(t)M^{(2)}(t)/(M(t))^21`$. If a transition is second order, $`U(t)`$ obeys a power law at the transition point. Therefore it can also be used for the determination of $`K^{}`$. Results are included in Fig. 2. Summarising all these measurements leads to $`K_{7s}^{}=1.293008(7)`$. For the 5-state Potts model, the transition is extremely weak. One should carry out the simulations very carefully. To locate $`K^{}`$, we have first performed simulations with $`L=560`$ for $`K=1.174359`$ ($`K_c`$), $`1.174946`$, and $`1.175533`$, up to $`t_{max}=\mathrm{10\hspace{0.33em}000}`$ with $`1800`$ samples. The resulting $`K_{5s}^{}=1.17445(6)`$ is not accurate enough. Therefore another simulation has been carried out at $`K=1.174570`$, which is much closer to $`K^{}`$. In Fig. 3(a), the second moments for $`K=1.174359`$ ($`K_c`$) and $`1.174570`$ are displayed. With these data more accurate values for $`K^{}`$ are obtained and collected in the upper part of Fig. 4. We estimate the averaged $`K_{5s}^{}=1.174404(7)`$. Similar is the case for the determination of $`K^{}`$. We have first performed simulations with an ordered initial state with lattice sizes $`L=280`$ and $`560`$ for $`K=1.173890`$, $`1.174125`$, and $`1.174359`$ ($`K_c`$), up to $`t_{max}=\mathrm{10\hspace{0.33em}000}`$ with total samples $`725`$. From the data for the magnetisation we estimate a relatively rough $`K_{5s}^{}=1.17428(9)`$. Then we performed simulations at $`K=1.174280`$ and $`1.174359`$ ($`K_c`$) up to $`t_{max}=\mathrm{40\hspace{0.33em}000}`$. The results are not sensitive to whether we take $`t_{max}=\mathrm{10\hspace{0.33em}000}`$ or $`\mathrm{40\hspace{0.33em}000}`$. In Fig. 3(b), the magnetisation at $`K=1.174280`$ and $`1.174359`$ ($`K_c`$) are plotted. From the lower part of Fig. 4, we obtain a final value $`K_{5s}^{}=1.174322(2)`$. In Table 1, all results for $`K^{}`$ and $`K^{}`$ have been collected. For both the 7-state and the 5-state Potts model, $`K^{}`$ and $`K^{}`$ are clearly above and below the transition point $`K_c`$ respectively. Our short-time dynamic approach indeed provides a safe criterion for a weak first order transition. Since our dynamic measurements are carried out in the short-time regime when the spatial correlation length is still short, we can easily control the finite size effect. We also do not have the problem of generating independent configurations and therefore do not suffer from slowing down. After excluding the finite size effect, the measurements are sensitive enough to distinguish a finite but very large spatial correlation length in equilibrium from an infinite one. This is why our method is successful. With the pseudo critical points in hand, assuming similar scaling laws as in second order transitions zhe98 , one can estimate corresponding pseudo critical exponents. At $`K^{}`$, e.g, $$M^{(2)}(t)t^{c_2},c_2=(d2\beta /\nu )/z.$$ (1) At $`K^{}`$, for the magnetisation, $$M(t)t^{c_1},c_1=\beta /\nu z,$$ (2) while for the Binder cumulant $$U(t)t^{c_U},c_U=d/z.$$ (3) Here $`d`$ is the dimension of the lattice, $`\beta `$ and $`\nu `$ are the well known static exponents and $`z`$ is the dynamic exponent. However, the values of the exponents at $`K^{}`$ and $`K^{}`$ can be different. Plotting the observables vs. $`t`$ in log-log scale, one measures the corresponding exponents from the slopes. The results are given in Table 2. Here we should admit that accurate values for complete sets of exponents can not be obtained so easily. One still needs much more careful simulations. An important reason is that $`K^{}`$ and $`K^{}`$ are not real critical points. They are also rather close to each other. In conclusions, we have proposed a short-time dynamic approach to weak first order transitions. From non-equilibrium short-time behaviour of two dynamic processes starting from random and ordered initial states, pseudo critical points $`K^{}`$ and $`K^{}`$ are determined. Difference of $`K^{}`$ and $`K^{}`$ distinguishes a weak first order transition from a second order one. Since the measurements are carried out in short-time regimes, the method does not suffer from slowing down. Different from many techniques developed in simulations in equilibrium, our method is not based on the finite size scaling. A simple average of $`K^{}`$ and $`K^{}`$ gives a rather good estimate of the transition point $`K_c`$, especially for very weak transitions. For example, for the 5-state Potts model $`(K^{}+K^{})/2=1.174363`$ and the relative deviation from the exact $`K_c`$ is only the order of $`O(10^6)`$. It is interesting to investigate how to obtain an accurate $`K_c`$ for not too weak transitions. Furthermore, how other relevant observables like the specific heat and energy distribution evolve in non-equilibrium dynamics is also an important topic. It is challenging whether from short-time dynamics one can estimate the latent heat and the discontinuity of the order parameter in equilibrium. Acknowledgements: The authors thank Z.B. Li and Q. Wang for valuable comments. This work is supported in part by the Deutsche Forschungsgemeinschaft, Az. Schu 95/9-2 and TR 300/3-1.
warning/0003/astro-ph0003373.html	ar5iv	text	# Mode detection from line-profile variations ## 1. Historical development While the presence of radial velocity variations in $`\delta `$ Scuti itself was detected just at the beginning of the century (Wright 1900), the first hints about the possible presence in the same star of line profile variations (hereinafter LPV) were due to Struve about half a century later (Struve 1953). After a decade of quiescence the studies of the phenomena associated with LPV were resumed by several authors in the mid sixties and seventies. Several different phenomena were announced, such as variable emission in the cores of strong lines, large equivalent width variations and so on (for a review see the paper by Breger 1979). Unfortunately those findings were probably connected to the use of the photographic plate as a detector. As a matter of fact later studies of the same objects performed with photoelectric detectors were not able to confirm any of those findings. The first reliable observations of LPV in $`\delta `$ Scuti stars were performed at the beginning of the eighties with the advent of the Reticon detector. Several stars were observed by Campos & Smith (1980) and Smith (1982) but unluckily most of them for a few hours in one night only, i.e., while the presence of line profile variations was clearly detected, the data were not sufficient to analyze the variations and to detect periodicities. Generally the few data were phased with the periods available from photometry or radial velocity curves (usually at that time only one or two periods per star were known). Among the stars studied by these authors only 28 Aql had a sufficient dataset (4 consecutive nights) which allowed the analysis of line width and radial velocity variations. A period–finding program confirmed the presence in both parameters of the two known periods derived from photometry (Smith 1982). The successive papers showed, beside the changing shape of the line profiles, the presence of moving sub-features in the lines (e.g., Yang & Walker 1986, Walker et al. 1987). In these cases the data were also not adequate for independent period searches, although two attempts were made by Kennelly et al. (1991, 1992), who analyzed 2.5 hours and 5 hours of observations of $`k^2`$ Boo and $`\gamma `$ Boo, respectively. They did it by sampling at 10 km s<sup>-1</sup> intervals the intensity residuals in the line profiles derived from the subtraction of the average profile, and computing the Discrete Fourier Transform. Obviously, given the extremely short temporal baseline, the frequency resolution was very poor, allowing only a rough confirmation of the photometric periods. In the following years it was realized that more intensive campaigns on a few selected objects were more fruitful than sporadic observations of many targets. This was in part the consequence of the photometric studies that were showing that usually many close frequencies were present in the pulsation spectra and consequently longer baselines were needed in order to resolve these frequencies. Now before discussing in detail the results concerning these campaigns we shall briefly review the phenomenology associated with the LPV, the constraints that it imposes on the observational programmes and the techniques we can adopt to analyze the data. ## 2. Phenomenology Fig. 1 shows the LPV observed during one night in the spectrum of BV Cir; in order to clearly show the LPV the average spectrum has been subtracted from each spectrum. We see the presence of waves propagating across the profiles of the lines at 4501, 4508 and 4515 Å from short to long wavelengths, while in the continuum regions there are only noise fluctuations. In the figure the borders of the three lines have been marked by the solid vertical lines. As we can judge from the figure the typical time-scale for a feature to cross the line is of the order of a few hours, while the amplitudes of these features are of the order of a few thousandths of the continuum intensity. As is well known we can interpret such variations as the result of the presence of non-radial pulsations: in each zone of the visible stellar disk the pulsation velocity combines with the rotational one and therefore the flux from this zone, because of the Doppler effect, gives a contribution to the line profile at a distance from the line center corresponding to its radial velocity. In this sense, if the pulsation velocity is small with respect to the rotational one, the line profile broadened by rotation supplies a one-dimensional “Doppler imaging” of the stellar disk (Vogt & Penrod 1983). For this reason it is a common practice, when describing the variations across the line profiles, to transform wavelengths into velocities by means of the Doppler formula, assuming the zero velocity at the line center (see, e.g., Fig. 10). In the case of Fig. 1 the moving features indicate that the prevailing pulsation modes are crossing the visible stellar disk moving in the same direction of rotation (i.e., they are prograde as seen by the observer). If the set of spectra has a rather good phase coverage of the periods of the different pulsation modes we can get a first guess of some quantities that we need in the further data analysis. First, we can compute the average of the spectra and assume that this approximates the nonpulsation spectrum of the star (however it can be seen that this average spectrum tends to have broader lines that the true nonpulsation one, see, e.g., Hao 1998). Then we can estimate at each wavelength step (or for each pixel of the spectrum) the standard deviation of the normalized flux about the mean. As an example Fig. 2 shows in the upper panel the average spectrum of X Caeli while the lower panel shows the standard deviation of the individual spectra about the mean. This panel clearly shows that the variability is restricted inside the line profiles, furthermore it allows the definition for each line of the limits in which the profile variations are present. In the continuum zones the reciprocal of the standard deviation of the normalized fluxes about the continuum supplies an estimate of the $`S/N`$ of each spectrum, and the reciprocal of the average of the values of the standard deviations of all the spectra about the average one in the same zones gives an estimate of the average $`S/N`$ at the continuum level. For example in the X Cae case the lower panel of Fig. 2 suggests an average $`S/N`$ of about 230<sup>1</sup><sup>1</sup>1As $`S/N`$ in Astronomy is usually intended the ratio between amplitudes and not between powers as preferred by engineers.. Fig. 3 shows the standard deviation of the normalized flux in each pixel for six stars arranged, from the bottom to the top, in order of increasing projected rotational velocity. We can see that as the line profiles broaden the flux variation decreases, until for the fastest rotator ($`v\mathrm{sin}i=196`$ km s<sup>-1</sup>) this variation is barely discernible. It can be also appreciated that it becomes increasingly difficult, as $`v\mathrm{sin}i`$ increases, to find unblended lines and continuum regions which can be used for the spectrum normalization. While it is obvious that in order to detect high-degree modes it is necessary to study stars with high $`v\mathrm{sin}i`$ <sup>2</sup><sup>2</sup>2As a rule of thumb the highest detectable degree is of the order of $`v\mathrm{sin}i/W_i`$ (Kennelly et al. 1992), where $`W_i`$ is the intrinsic line width, i.e., the width that the line should have if rotational broadening was not present. , if we consider low-degree modes, it can be shown that in weak lines their $`S/N`$ ratio decreases approximately with $`(v\mathrm{sin}i)^{1/2}`$, and therefore these modes are more easily detectable in stars with low projected rotational velocity. Starting from the average line profile we can also get a first estimate of the projected rotational velocity and of the intrinsic line width. This can be accomplished by performing a non-linear least-squares fit of a rotationally broadened Gaussian profile to the observed one. Fig. 4 shows as an example the fit derived in this way on the FeII 4508 Å line of the star HD 2724 (Mantegazza & Poretti 1998). These two quantities can also be estimated using the zero-points of second and fourth moment (Balona 1986b)<sup>3</sup><sup>3</sup>3For the definition of the line moments see the paper by Aerts & Eyer, these proceedings.. As we have previously said, the average profile tends to be larger than the nonpulsation one. A way to get a better estimate of these two quantities is discussed at the end of this paper. ## 3. Observational constraints From what we have shown in the previous section we can estimate what the observational constraints should be in order to detect and to study the LPV in $`\delta `$ Scuti stars. The typical time-scales of the observed variations (a few hours) impose a limit on the integration times, which should be a fraction of these quantities. According to our experience this limitation is not very stringent for the mode detection, for instance in the case of FG Vir we have detected the 34.1 cd<sup>-1</sup> mode with an integration time of 15 min, which corresponds to about 35% of the pulsation period (see Table 1). For the mode identification the discourse is more delicate, because even if with a relatively long integration time we are still able to detect the fundamental frequency of a mode and to correct its amplitude for the damping due to the integration in time, its harmonics, which are of lesser amplitudes, can be lost in the noise, because of the larger damping due to their shorter periods, and therefore we would not be able to recover the correct shape of the signal. The amplitudes of the perturbation across the line profiles are for most of the modes of the order of a few thousandths of the continuum intensity, and therefore, to be reliably detected and measured, we need a $`S/N`$ of at least a few hundred. Finally, it remains to be evaluated what the minimum spectral resolution should be. Since, as it is explained in the paper by Aerts & Eyer (this proceedings), we can approximate the observed line profiles as the convolution between an intrinsic profile ($`W_i`$) and the rotational broadening one, the resolution should be $$R=\lambda /\delta \lambda c/W_i.$$ Since typically $`W_i10km/s`$, then $`R30000`$. The three requirements: short exposures, high $`S/N`$ and high resolution, impose that we collect a lot of photons in short times, and this implies that either we limit ourselves to the study of very bright targets, or we observe with large telescopes, or, by observing a wide spectral region, we find a way to add the information of the LPV contained in many of the observed lines without blurring it. Up to now 4 different approaches have been adopted to add the information of several lines: a) the straight average of the individual line profiles (e.g., Kennelly et al. 1996); b) the average of the line moments (e.g., Mantegazza & Poretti 1996), c) the computation of the cross-correlation function (e.g., Korzennik et al. 1995), d) the deconvolution of the observed profiles by the intrinsic one (Kennelly et al. 1998, see also the paper by Aerts & Eyer, these proceedings). None of these approaches are free of problems, and in particular all of them assume that the intrinsic profile should be the same for all the lines and therefore in all cases a careful selection of lines with similar characteristics should be done. ## 4. Analysis techniques Now we shall briefly review how we can study the variations in a line profile. Three different techniques have been adopted to search for the periodicities present in LPV: * the pixel-by-pixel analysis; * the frequency analysis of the moment time series; * Fourier Doppler imaging. ### 4.1. Pixel-by-pixel The pixel-by-pixel analysis is based on the fact that during an oscillation cycle the flux measured in the same pixel of the line profile (i.e., at the same wavelength) fluctuates with the same period. Fig. 5 shows this fact for the TiII 4501 line of the star X Caeli. The profiles have been phased according to the dominant pulsation mode (7.39 cd<sup>-1</sup>) and cover one complete cycle. The dot indicates for each profile the flux at a given fixed wavelength. Therefore we can extract from the set of all the profiles of the same line the time series constituted by the fluxes in the same pixel and then analyze them with the usual techniques developed to study one–dimensional time series (e.g., light-curves). We have seen in the previous section how the diagram showing the flux standard deviation in each pixel can be used to define the line borders and hence to detect which pixels can supply useful time series. Before extracting the individual pixel time series it is necessary to rebin the spectrograms in order to remove the observer’s velocity variations due to the Earth’s revolution and rotation. Even the latter cannot be neglected because its amplitude (about 0.5 km s<sup>-1</sup>) can be comparable to those of the pulsation modes. Two techniques have been the most frequently used to analyze the pixel-by-pixel time series: a) the least-squares power spectrum (Vaniĉek 1971), b) the CLEAN algorithm (Roberts et al. 1987) Fig. 6 shows the application of this approach to analyze the LPV in the FeII 4508 line of BV Cir (1998 campaign, Mantegazza et al. 1999). The upper left panel shows the average line profile, the central panel contains the pixel-by-pixel spectra computed with the least–squares technique without known constituents (i.e., the spectrum contains all signals present in the line variations, for details see below), and the lower panel shows the global least-squares spectrum (which is equivalent to a weighted average of the individual pixel-by-pixel spectra). In this case we are in the presence of rather complicated LPV which contain several modes with low, medium and high frequencies in the observer’s reference frame<sup>4</sup><sup>4</sup>4Because of stellar rotation a non-axisymmetric non-radial mode with an oscillation frequency $`\nu _0`$ in a reference frame co-rotating with the star is seen from the observer as having a frequency $`\nu =\nu _0m\mathrm{\Omega }`$, where $`m`$ is the mode azimuthal order and $`\mathrm{\Omega }`$ is the stellar rotational frequency.. Since the data were obtained in a single site campaign the peak of each mode is flanked by rather strong 1 cd<sup>-1</sup>aliases, and this, coupled with the intrinsic complexity of the pulsation spectrum, explains why the computed power spectrum is so entangled. In order to have a clear picture of the true pulsation spectrum is necessary to detect one by one all the periodicities or to analyze the data with the CLEAN algorithm. The result of the application of this technique is shown in Fig. 7. Other examples of the application of the CLEAN algorithm to the pixel-by-pixel analysis are in the papers by Bossi et al. (1998) and De Mey et al. (1998). The power spectrum of Fig. 7 is much more easily readable than that of Fig. 6 even if there are some aliases still present and probably in some cases aliases have been preferred to the true peaks. The reason is due to the fact that the CLEAN approach to select the peaks is completely automatic, and among other things it does not allow to select between competing aliases the use of the a priori knowledge, such as for example the fact that we know some correct periods from the light-curve analysis. In our opinion and according to our experience, the most reliable approach for looking for multi-periodicities is to proceed to detect one by one the periods with the least squares technique. This approach has been slightly modified to best analyze LPV (Mantegazza & Poretti 1999) in order to define a “global” least-squares spectrum which includes the information of the variability of the whole profile. #### The least-squares algorithm generalized to study LPV In general our data consist of a set of profiles of a spectral line $`P(\lambda _j,t_k)`$ ($`j`$ is the pixel number and $`t_k`$ is the time of the $`k`$-th spectrogram) whose global variance can be defined as $$\sigma _T=\underset{j,k}{}w_k^2\left[P(\lambda _j,t_k)P_0(\lambda _j)\right]^2$$ $`(1)`$ where $`P_0(\lambda _j)`$ is the time averaged profile and the $`w_k`$’s are the normalized weights derived from the $`S/N`$ of the individual spectrograms. The frequency analysis of the variations present in these profiles can be iteratively performed in the following way: if we have already detected $`m`$ periodic sinusoidal components (“known constituents”) and we are looking for the $`m+1`$, we explore the useful frequency range (0$`<\nu _i<\nu _{max}`$ cd<sup>-1</sup>) by fitting each pixel time series $`j`$ with the series $`p_{i,j}(t_k)=\overline{p}_i+{\displaystyle \underset{l=1,m}{}}A_{i,j,l}\mathrm{cos}(2\pi \nu _lt_k+\varphi _{i,j,l})`$ $`+A_{i,j,m+1}\mathrm{cos}(2\pi \nu _it_k+\varphi _{i,j,m+1})`$ (2) where $`\overline{p}_i,A_{i,j,l},\varphi _{i,j,l}`$ (with $`1lm+1`$) are the free parameters. Then we compute the global reduction of variance defined as $$RF_i=1\underset{j,k}{}w_k^2(p_{i,j}(t_k)P(\lambda _j,t_k))^2/\sigma _m$$ $`(3)`$ where $`\sigma _m`$ is the global residual variance after the fit of the line profile variations with the $`m`$ “known constituents”. The frequency $`\nu _i`$ giving the highest $`RF`$ (or one of its 1 cd<sup>-1</sup> aliases if there is a better agreement with the photometrically detected modes) is then selected as the $`m+1`$–th known constituent and the procedure is iterated again. The search ends when no dominant peaks are apparent in the last spectrum. On this point a more quantitative criterion should be developed such as the one adopted by Breger et al. (1995) to assess the physical reality of peaks detected from light-curve frequency analysis. At the end of this procedure, after having detected M known constituents, we can perform a final fit of $`P(\lambda _j,t_k)`$ and derive the functions: $`\overline{p}_M(\lambda _j)`$ (i.e., the best estimate of the average line profile), $`A_l(\lambda _j),\varphi _l(\lambda _j)`$ (with $`1<l<M`$) and also their formal errors. As an example we report in Fig. 8 the frequency analysis of the variations in the 4508Å line of the star HD2724. Each panel contains a “global least-squares spectrum” computed by introducing the frequencies of the terms detected in the previous spectra as “known constituents”. Fig. 9 summarizes in the upper panel the pulsation spectrum as derived from the frequency analysis of the previous figure. In the lower panel we show the same spectrum as obtained using the CLEAN algorithm. In this case the two approaches supply the same mode detection (the only discrepancy is on the selection of the lowest frequency peak where there is an uncertainty of 1 c/d, but as we can see in panel 4 of Fig. 8, the difference between the two peaks is very marginal). Fig. 10 shows the behaviors of amplitudes and phases across the 4508 Å line profile of the terms detected in the previous example (functions $`A_l(\lambda _j),\varphi _l(\lambda _j)`$ and their corresponding formal errors). The positions on the line profile are expressed in Doppler velocities as described in Section 2. The term at 0.07 cd<sup>-1</sup> is not reported because it was probably introduced by an instrumental effect (see Mantegazza & Poretti 1999). It can be observed that the phase diagram of the 5.31 cd<sup>-1</sup> term has the typical behavior of a low–degree retrograde mode (its phase increases in the same direction as the stellar rotation, i.e., with the wavelength). The opposite behavior is shown by the phases of the 8.58 cd<sup>-1</sup> term, which is typical of a moderately high–degree prograde mode. The diagrams that show the behavior of the phases of each mode across the line profile are very instructive to show the nature of the non-radial modes present in $`\delta `$ Scuti stars. In order to demonstrate this Fig. 11 shows a zoo of behaviors of such phases as we have derived them from our studies. Seven different modes are represented in the panels ranging from a high-degree retrograde mode (bottom panel) to a high-degree prograde mode and passing from an axisymmetric mode ($`m=0`$, third panel from the bottom). For each panel we give the name of the star to which the mode belongs, the frequency of the mode in the observer’s reference frame and its $`\mathrm{}`$ and $`m`$, as derived from the LPV fit technique described in the last section of this paper. The abscissae give the positions across the line profile with Doppler velocities (or equivalently wavelengths) increasing from left to right. +1 and –1 indicate a distance from the center corresponding to $`+v\mathrm{sin}i`$ and $`v\mathrm{sin}i`$ respectively (i.e., in the Doppler map they correspond approximately to the limbs of the visible stellar disk). ### 4.2. The moment time-series analysis The use of the line moments to study LPV was first introduced by Balona (1986a). Since the moments are also useful for mode identification, their definition is given in the paper by Aerts & Eyer (these proceedings). It suffices here to say that the zeroth order moment measures the equivalent width, while the first-order one the position of the line barycenter and hence it can be used as a measure of the stellar radial velocity. In his paper, Balona (1986a) gives some useful suggestions how to derive these quantities; to this end it is also useful to look at the paper by Balona et al. (1996). As an example Fig. 12 shows the behavior of the first three moments and of the zero-order one during one night of observations of the star X Cae (Mantegazza & Poretti 1996). Their variations are compared with the simultaneous variations in $`V`$ light and $`BV`$ color index. It can be observed that the first and third moment curves have a similar shape and are in anti-phase with respect to the light variations, while the second moment has a different behavior. As in the case of the pixel-by-pixel time series the moment time series can be analyzed with the same techniques used to study light-curves. Since the moments are quantities integrated on the whole line profile, they are more sensitive to low degree modes, because the variations due to the high degree oscillations tend to be averaged out by the integration. One interesting property of the moments is the fact that axisymmetric modes contribute with their fundamental frequency only to the variations of odd moments. Therefore the frequency analysis of the moment time series offers an easy way to detect these modes, and moreover, since often we face stars with a lot of excited modes, reducing the number of the modes affecting even moment time series makes it easier to analyze them. As an example of this fact we show in Table 1 the results of the frequency analysis of the first 5 moments of FG Vir (1995 campaign, Mantegazza et al. in preparation). The detected frequencies in each moment time series are listed for each mode beside the well known and reliable values derived from light-curve analysis (Breger et al. 1998). Since the spectroscopic baseline is rather short (5 consecutive nights) and the campaign was single site, some 1 cd<sup>-1</sup> aliases were picked up instead of the correct values. While up to 13 modes have been detected in the odd moments, the analysis of the even ones has allowed the detection of only two. According to what we said above it is evident that most of the modes should be axisymmetric. This is not surprising in view of the fact that the star is probably seen almost pole on, and therefore the detection of axisymmetric modes is favored by a selection effect. ### 4.3. Fourier Doppler Imaging (FDI) This technique was first introduced to study LPV in $`\delta `$ Scuti stars by Kennelly et al. (1992) and successively developed and improved by Kennelly et al. (1998). A theoretical study of its properties has been performed by Hao (1998). This technique allows at the same time the detection of the modes producing LPV and the estimation of their non-radial degree ($`\mathrm{}`$). Because it is based on the computation of a two-dimensional Fourier Transform, which considers at the same time temporal and spatial variations, this technique is mainly sensitive to high-degree modes (see Kennelly et al. 1998, where these authors show also that the technique is more sensitive to odd azimuthal order modes than to even ones) and requires targets with rather high $`v\mathrm{sin}i`$’s. For a description of it and an example of a two-dimensional Fourier amplitude spectrum see the paper by Aerts & Eyer (these proceedings, Fig. 8). One of the advantages of this method is that, by separating the detected modes both in temporal and spatial frequency, it allows, even if the temporal baseline is not adequate, the resolution of modes with close temporal frequencies but very different spatial ones. ## 5. Recent campaigns Tables 2 and 3 summarize the most recent spectroscopic campaigns on $`\delta `$ Scuti stars. In Table 2 we give for each star the epoch of the observations, the number of sites contributing to the campaign, whether or not there is simultaneous photometry, the adopted analysis technique and finally the reference. Table 3 gives more technical details such as the resolution, the spectral range, the average $`S/N`$ of the individual spectrograms at the continuum level, the number of observing nights, the total useful observing time, the number of gathered spectra, the average exposure time and finally the lines (or their number if they were many) whose LPV were analyzed. As we can see from Table 3 the observing parameters generally satisfy the requirements we derived in Section 3. We note moreover that moving from the oldest to the most recent campaigns their duration tends steadily to increase. This fact testifies that the observers became increasingly convinced that, given the complexity of the pulsation spectra, in order to get reliable mode detections baselines rivaling those of the photometric campaigns were necessary. Fig. 13 shows the position of these objects in the color-magnitude diagram. We can see that they are rather uniformly distributed across the instability strip in the giant region. Main Sequence objects are generally missing because their fainter apparent magnitudes make it difficult to get spectrograms satisfying the observational constraints discussed in Section 3. Table 4 summarizes the number of periods detected in each of these stars according to the technique adopted to analyze LPV. For comparison in the same table the number of periods derived from the light-curve analysis are reported. The uncertain detections are given within brackets. The RV column reports the detection from the radial velocity variations, where here for radial velocities we mean those derived with techniques different from the first moment computation (for instance with correlation techniques). The results regarding V480 Tau are extremely preliminary (Hao, private communication) and many periods are expected to be found. It is possible to see that in the cases where the temporal baseline is adequate the number of spectroscopic detections is comparable with those of photometric detections. Some of the best spectroscopically studied objects have only few photometrically detected modes, because careful photometric campaigns have not yet been performed on these stars. A better coordination in the selection of the spectroscopic and photometric targets should be desirable in the future in order to get the most complete picture possible of the pulsation spectra. Figures 14 and 15 show the best examples of pulsation spectra derived from spectroscopic analysis. For each star the projected rotational velocity is given below its name. The epoch of the observations is also given at the right of the star’s names because, as we will show in the next subsection, the pulsation spectra can change considerably with the time. The heavy vertical segment at the top of each panel marks the expected position of the fundamental radial mode. In the case of $`\theta `$ Tuc this line is dashed, because it belongs to a binary system, so the evaluation of its physical parameters is uncertain, and therefore the position of this line gives only the lower limit. We see that in the cases of FG Vir, X Cae, BV Cir and $`\tau `$ Peg there are a few modes with frequencies in the observer’s reference frame below this value. In the cases of X Cae, $`\tau `$ Peg and for one of the modes of BV Cir these are retrograde $`p`$ modes, which in the co-rotating reference frame have frequencies above the fundamental radial value. For the other two modes of BV Cir the question is open, because their identification is unclear, while the two modes of FG Vir could be $`g`$ modes (Breger et al. 1999). The amplitudes of the modes indicated with solid lines were derived from pixel-by-pixel analysis or from FDI and have been arbitrarily scaled for each star. The modes drawn with a dashed line for BB Phe, HD 101158, and $`\theta _2`$ Tau were detected from the moment or the radial velocity analysis, and therefore their amplitudes cannot be compared with those of the others. Several of the modes shown in these figures have not been photometrically detected. The number of modes with a spectroscopic detection only is indicated for each star in the last column of Table 4. We see therefore the complementarity between the two approaches and the need for both if we want to get the whole pulsation spectrum of the star. The higher the projected rotational velocity is, the higher the number of purely spectroscopically detected modes is. This is an obvious selection effect, because high–degree modes can be detected only in fast rotators. This is confirmed by the fact that, for instance, for a slow rotator such as FG Vir all the spectroscopically detected modes also have a photometric detection. In order to get an idea of how the pulsation spectrum of a fast rotating $`\delta `$ Scuti star should appear as seen in the co-rotating reference frame, the observed frequencies of the modes detected and identified by Kennelly et al. (1998) with the FDI in $`\tau `$ Peg have been corrected for the rotational effect assuming that all of them were sectoral. The result is shown in the lower panel of Fig. 16. For comparison the upper panel shows the pulsation spectrum in the observer’s reference frame and the dashed line shows the estimated position of the fundamental radial mode. We can see that in the co-rotating frame all the modes are clumped in the vicinity of the fundamental radial mode. The fact that a few of them have frequencies slightly lower than that of the fundamental radial one is probably due to the assumption that all the modes are sectoral, which probably for some modes is not true, and consequently their frequencies have been over-corrected toward low values. ### 5.1. Amplitude variations As we can see from Table 2, there are a few stars for which observations have been performed in two campaigns. The best cases are those of X Cae, BB Phe and BV Cir. The availability of two independent datasets is very useful because it allows: a) the check of the reality of the modes which have not independently been detected from photometry, b) the detection of variations in the mode amplitudes. For the above-quoted three stars most of the modes have been independently detected in the two data sets and hence we are quite confident about their reliability. In the cases of BB Phe and BV Cir there are large variations in the mode amplitudes. For both stars the strongest spectroscopic mode was different in the two seasons. Figure 15 shows the global least squares power spectra without known constituents of BV Cir for the two campaigns. We clearly see the different frequency content: in the 1996 data the strongest peak is at 13.85 cd<sup>-1</sup>, while in the 1998 data, even if this peak is still present, the highest one is at 17.28 cd<sup>-1</sup>. This last peak in 1996 was below the detectability level. Comparing the two spectra we can also easily perceive their different frequency resolution due to the different temporal baselines (6 and 12 days, respectively). Figure 18 compares the pulsation spectrum, as derived from the least-squares power spectrum analysis, of the two seasons. Seven terms are in common, but others are present only in one of the two datasets, in particular the three highest frequency terms present in 1998, and which are among the strongest, were not detectable in the 1995 data. Korzennik et al. (1995) on the basis of their FDI analysis of LPV in $`\upsilon `$ UMa claimed that there were strong variations in the mode amplitudes from night to night. In this case it is easy to demonstrate that the more plausible explanation is that this is due to beats between unresolved modes: analyzing the data on a night by night basis, where the observations in each night covered between 2.5 to 5 hours only, the frequency resolution was extremely poor (between 5 to 10 cd<sup>-1</sup>), hence the expected density of the pulsation spectrum amply allows the presence of these beats! ## 6. Mode identification by LPV fit in multiperiodic stars In the paper by Aerts & Eyer (these proceedings) it is stated that the best approach to the mode identification is given by the direct fit of LPV, but unfortunately, according to them, this is possible for mono-periodic pulsators only. Recently, we (Bossi et al. 1998, Mantegazza & Poretti 1999, Mantegazza et al. 1999) have developed a technique that allows the fit of LPV for multiperiodic pulsators and moreover, if simultaneous spectroscopic and photometric observations are available, the technique can proceed to the mode identification by simultaneously fitting both variations. In the presence of several simultaneously excited modes it is necessary to separate the contributions of each mode to the overall LPV and light variations. To do this we approximate the perturbations $`\mathrm{\Delta }p_j(\lambda ,t)`$ induced by mode $`j`$ on the line profile as: $$\mathrm{\Delta }p_j(\lambda ,t)=\underset{i}{}A_{ij}(\lambda )\mathrm{cos}\left(2\pi i\nu _{jt}+\varphi _{ij}(\lambda )\right)$$ $`(4)`$ where the sum is on the Fourier harmonics of the j–mode. We can also estimate the formal errors on $`\mathrm{\Delta }p_j(\lambda ,t)`$ ($`\delta \mathrm{\Delta }p_j(\lambda ,t)`$) from error propagation by $`\delta A_{ij}(\lambda )`$ and $`\delta \varphi _{ij}(\lambda )`$. These last quantities as well as $`A_{ij}(\lambda )`$ and $`\varphi _{ij}(\lambda )`$ have been derived in Section 5.1 by simultaneously fitting all the detected modes and their harmonics to the observed LPV. Usually with a $`S/N`$ of a few hundreds only the fundamental harmonic of each mode is easily detectable. Only for rather strong modes, such as for the dominant mode of X Cae (Mantegazza et al. 1999), the first harmonic is also detectable. Fig. 20 shows for instance the amplitude (right panels, dashed lines) and phase (left panels, dots) variations across the 4501 Å line profile of X Cae derived from this simultaneous least-squares fit of the fundamental frequency (upper panels) of the dominant mode and its first harmonic (lower panels). The reconstructed LPV (equation 4) due to this mode are shown as solid lines in Fig. 19. In this case the functions $`\mathrm{\Delta }p_j(\lambda ,t)`$ have been computed for ten equi-spaced phases of a complete cycle. We can try to fit the functions $`\mathrm{\Delta }p_j(\lambda ,t)`$ (eq. 4) with perturbations computed with a model of a non-radial pulsating star viewed at a certain inclination $`i`$. So for each plausible choice of $`\mathrm{},m,i`$ we can build a discriminant $$\sigma _p(\mathrm{},m,i)=\underset{\lambda }{}\underset{t}{}\frac{(\mathrm{\Delta }p_j(\lambda ,t)\mathrm{\Delta }p_c(\lambda ,t,\mathrm{},m,i))^2}{\delta \mathrm{\Delta }p_j(\lambda ,t)^2}$$ $`(5)`$ where $`\mathrm{\Delta }p_c(\lambda ,t,\mathrm{},m,i)`$ are the computed profile variations which best fit the observed ones. Moreover, if we have simultaneous photometric observations, we can obtain, by simultaneously fitting all the terms detected in the light curve, amplitude and phase with respective errors for the $`j`$ mode. Therefore we can calculate its light variations and relative errors ($`l_j(t)`$ and $`\delta ł_j(t)`$) and compare them with those predicted by the best fitting models and obtain the discriminant: $$\sigma _l(\mathrm{},m,i)=\underset{t}{}(l_j(t)l_c(t,\mathrm{},m,i))^2/\delta ł_j(t)^2.$$ $`(6)`$ A global discriminant is then defined as: $$\sigma _T(\mathrm{},m,i)=\sigma _p(\mathrm{},m,i)+\sigma _l(\mathrm{},m,i).$$ $`(7)`$ This is the function which is minimized with a non-linear least-squares fit for each detected $`j`$ mode and for any choice of $`\mathrm{},m,i`$. The model we used to compute the synthetic line profile variations $`\mathrm{\Delta }p_c(\lambda ,t,\mathrm{},m,i)`$ and light variations $`l_c(t,\mathrm{},m,i)`$ is the one described by Balona (1987). For each assigned $`\mathrm{},m,i`$ the computed profile variations can be modeled according to the amplitude and phase of vertical ($`v_r`$) and horizontal ($`v_r`$) velocities and flux variations. For $`\delta `$ Scuti stars usually $`v_hv_r`$ and therefore in order not to introduce into the model too many free parameters, we keep the usual theoretical relation (e.g., Heynderickx et al. 1994) $`v_h=74.4Q^2v_r`$ ($`Q`$ pulsation constant) and $`\psi _h=\psi _r`$. The observed light variations constrain strongly the computed flux variations, so it is very useful to have simultaneous spectroscopic and photometric data, otherwise in order to get meaningful physical results it is better to fit a simplified model which considers velocity variations only, neglecting flux variations. By applying this approach to fit the reconstructed profile and light variations due to the dominant mode of X Cae (solid lines of Fig. 17) we found that the least global discriminant is supplied by a non-radial mode with $`\mathrm{}=1`$, $`m=1`$, and $`i=70^{}`$. This result is in agreement with that derived from the independent dataset of the previous (1992) observing season (Mantegazza & Poretti 1998), and which was obtained from the moment variation fits with the technique developed by Balona (1987) (see also Balona et al. 1996), and which is different from the approach with the moments described in Aerts & Eyer (these proceedings). The LPV computed according to this model are represented as dashed lines in Fig. 19 and explains the 94% of the variance. The same model fits the observed $`B`$ variations due to this mode with a standard deviation of 0.4 mmag. Finally from the computed LPV we can also derive the behavior of amplitude and phase across the line profile both for the fundamental and first harmonic terms. These functions are represented as solid lines in the panels of Fig. 20. ### 6.1. Improved estimate of $`v\mathrm{sin}i`$ and $`W_i`$ As we have said in Section 2, the average line profile tends to be wider than the nonpulsation one and consequently there is the risk that both the projected rotational velocity and the intrinsic line width can be overestimated. A better way to estimate these two quantities is to give them as two more free parameters in the model which is non–linearly fit to LPV. For instance in the case of X Cae from the 4501Å line we get $`v\mathrm{sin}i`$=69.0 km/s and $`W_i`$= 11.9 km/s from the average line profile, while the non-linear LPV fit with the best-fitting mode ($`\mathrm{}=1`$, $`m=1`$) for the 7.39 cd<sup>-1</sup>mode supplies $`v\mathrm{sin}i`$=68.6 km/s and $`W_i`$=7.7 km/s. While the estimate of the projected rotational velocity has not appreciably changed, the intrinsic line width has been considerably decreased. ## 7. Conclusions In this paper we have presented the main observational characteristics of LPV in $`\delta `$ Scuti stars, and we have discussed how they can be observed and which techniques can be adopted to detect the pulsation modes. Moreover, a review of the main results obtained from the application of these techniques has been given. As we have seen, the number of stars carefully studied is at the moment scanty, so it is difficult to draw very general conclusions, although some facts, which can be useful as guidelines for future studies, can be pointed out: * The careful observation and analysis of LPV allows several pulsation modes to be detected. If data with adequate $`S/N`$ and temporal baselines are available the number of spectroscopically detected periods rivals the number which are photometrically detected. Many of these detections are purely spectroscopic, especially for stars with high $`v\mathrm{sin}i`$: photometric observations alone are not sufficient to get the complete pulsation spectrum. * The variability of the mode amplitudes is quite common. Sometimes some modes have amplitudes below the detectability threshold. Again, to get the complete pulsation spectrum, we need to observe the star at different epochs, and it may be desirable to get higher S/N data in order to better monitor the evolution of the amplitudes. * Period detections can be rather reliable: several of the spectroscopically detected modes have also been photometrically detected and moreover, for the stars with two spectroscopic campaigns, many of the periods have been independently detected in both datasets. Again, to do a good job, we need high $`S/N`$ with long baselines. * The three different analysis techniques presented in Section 4 tend to be complementary: the moment method is particularly suited for low-degree modes; FDI works preferably on high–degree ones and on stars with high $`v\mathrm{sin}i`$; the pixel-by-pixel analysis works both on low- and high-degree modes, and, for the detection of the low-degree ones, it is preferable to study stars with low $`v\mathrm{sin}i`$. For instance, with this technique we were able to detect several modes in FG Vir, which has a $`v\mathrm{sin}i`$ of only 20 km/s. * Given the complexity of the pulsation spectra, in order to get adequate frequency resolution and to avoid aliasing ambiguities, we need relatively long baselines (possibly comparable to the photometric ones) as well as multi-site campaigns. * We have already pointed out the need for data with good $`S/N`$. Given the very small amplitudes of most of the modes (a few thousandths of the continuum intensity), in order not to to merely detect them, but also to proceed to their identification, a $`S/N`$ better than 500 is desirable. * Since together with high $`S/N`$ data we also need high spectral and temporal resolutions and short exposure times (unless large telescopes are used for extended periods of time, which is almost a hopeless possibility), it is important to develop and to adopt techniques that add the information from many spectral lines. This would entail the use of medium size telescopes, which are more accessible. Much work is still to be done in this direction, especially if the added information is used not merely for mode detection but for identification, too. * As an immediate consequence of the previous items, it is evident that the best strategy to get a complete picture of the pulsation behavior of a $`\delta `$ Scuti star is to observe it simultaneously both photometrically and spectroscopically for several seasons and in the context of multi-site campaigns. In this respect the few researchers active in the field should coordinate their efforts in order to study a few carefully selected targets. Finally in the last part of this paper an approach to mode detection by LPV fitting in multiperiodic variables has been presented. This approach at the moment is not free of problems; for example, a better model of the LPV to fit to the data is needed, a model that at the same time is not too cumbersome, because the routine which uses it is called thousands of times during the minimization procedure. Another problem is the same one faced by the moment method described by Aerts & Eyer (this proceedings), i.e., the lack of confidence intervals for the minima of the discriminant, which makes it difficult to compare the competing modes. This second problem could be rather easily overcome once the first has been successfully addressed and the errors on the reconstruction of the LPV produced by each mode correctly estimated, because the discriminants (eqs. 5–7) have a $`\chi ^2`$ shape and therefore their statistical properties are well known. ##### Acknowledgments. I am grateful to E. Poretti for a critical reading of the manuscript and to M. 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warning/0003/hep-ph0003246.html	ar5iv	text	# 1 Introduction ## 1 Introduction It is difficult to overestimate the importance of the experimental discovery of the Higgs boson. It would not only help us to elucidate the dynamics responsible for electroweak symmetry breaking but it will most probably offer also an important clue as to the nature of the Physics beyond the Standard Model (SM). The paradigm for this new Physics, granted that a fundamental scalar drives electroweak symmetry breaking, is the Minimal Supersymmetric Standard Model (MSSM) : the most economical extension of the Standard Model that incorporates (softly-broken) Supersymmetry (SUSY). In spite of the uncertainties related to the origin of supersymmetry breaking (and therefore of the masses of the so far undetected supersymmetric particles), it is well known that the MSSM predicts the existence of a light Higgs particle with mass below about $`135`$ GeV (this bound depends sensitively on the top quark mass one uses; our present calculation intends to set a precise and firm bound). Unlike the case of the Standard Model (in which the mass of the Higgs boson is an unknown parameter), the mass of the light $`𝒞𝒫`$-even Higgs boson of the MSSM is calculable as a function of other masses of the model. A precise calculation of that mass is of prime importance for Higgs searches at LEP, Tevatron and the LHC, and is the topic of this paper. We recall at this point that the Higgs sector of the MSSM consists of two $`SU(2)`$ doublets, $`H_1`$ (which gives mass to down-type quarks and charged leptons) and $`H_2`$ (which gives mass to up-type quarks). The vacuum expectation values ($`v_{1,2}`$) of these doublets break the electroweak symmetry, after which, the Higgs spectrum contains two $`𝒞𝒫`$-even scalars ($`h^0`$ and $`H^0`$; with $`m_{h^0}m_{H^0}`$), one $`𝒞𝒫`$-odd pseudoscalar ($`A^0`$) and a pair of charged Higgses ($`H^\pm `$). At tree-level, the masses, couplings and mixing angles of these particles are determined by one unknown mass parameter (say $`m_{A^0}`$) and the parameter $`\beta `$, which measures the ratio $`v_2/v_1(\mathrm{tan}\beta )`$. In the limit $`m_{A^0}M_Z`$ all the Higgs particles except $`h^0`$ have masses $`m_{A^0}`$ and rearrange in a complete $`SU(2)`$ doublet almost decoupled from electroweak symmetry breaking, while $`h^0`$ remains light with $`m_{h^0}^2M_Z^2\mathrm{cos}^22\beta `$ and has SM properties. This bound (which applies for any value of $`m_{A^0}`$ and is saturated for $`m_{A^0}M_Z`$) is extremely important: it represents the limit that experimental bounds should reach to falsify the MSSM. In fact, the present experimental bound from LEP, including the latest data with up to $`\sqrt{s}=202`$ GeV, is $`m_{h^0}\stackrel{>}{_{}}107.7`$ GeV (for large $`m_{A^0}`$, case in which the SM limit is applicable; the limit falls to $`91`$ GeV for smaller $`m_{A^0}`$), which is well above this bound. This is not yet conclusive evidence against the MSSM because it does not take into account the radiatively corrected form of the mass bound. Radiative corrections to $`m_{h^0}^2`$ have been computed using three different techniques (or combinations of them): effective potential method , direct diagrammatic calculation and effective theory (or renormalization group) approach . The full one-loop radiative corrections to $`m_{h^0}`$ have been computed diagrammatically. The most important of these corrections come from top quark/squark loops and are given by $$\mathrm{\Delta }m_{h^0}^2=\frac{3m_t^4}{2\pi ^2v^2}\mathrm{ln}\frac{m_{\stackrel{~}{t}}^2}{m_t^2},$$ (1) where $`m_t`$ is the top quark mass, $`m_{\stackrel{~}{t}}`$ an average top-squark mass and $`v^2v_1^2+v_2^2=(246\mathrm{GeV})^2`$. This correction can be very large if $`m_{\stackrel{~}{t}}m_t`$, and in such case $`m_{h^0}`$ can evade easily the current experimental lower bound. This important $`𝒪(\alpha _t)`$ logarithmic correction to the dimensionless ratio $`m_{h^0}^2/m_t^2`$ \[here $`\alpha _th_t^2/(4\pi )`$, where $`h_t`$ is the top-quark Yukawa coupling\] can be most easily reproduced using renormalization group (RG) techniques. In addition, there is a finite (non-logarithmic) correction which may also be important, and which depends on the details of the top-squark spectrum. This correction is (assuming again for simplicity degenerate soft masses for the top-squarks) $$\mathrm{\Delta }m_{h^0}^2=\frac{3m_t^4}{2\pi ^2v^2}\left(\frac{X_t^2}{m_{\stackrel{~}{t}}^2}\frac{X_t^4}{12m_{\stackrel{~}{t}}^4}\right),$$ (2) where $`X_t=A_t+\mu \mathrm{cot}\beta `$ is the top-squark mixing parameter, $`A_t`$ the soft trilinear coupling associated to the top-Yukawa term in the superpotential and $`\mu `$ the supersymmetric Higgs mass parameter. Correction (2) is maximized for $`X_t^2=6m_{\stackrel{~}{t}}^2`$ (the so-called ‘maximal-mixing’ case). When using one-loop equations like (1) and (2) to compute the Higgs mass one has to decide whether to use on-shell (OS) or running values for the mass parameters that enter such formulae (and if the latter, at which scale to evaluate them). The difference between two such choices is of higher order, but can be non-negligible, especially because of the $`m_t^4`$-dependence of $`\mathrm{\Delta }m_{h^0}^2`$. Although RG techniques can be used to make an educated guess of the scale at which those mass parameters should be evaluated (see e.g. ), a precise answer to such questions could only be unambiguously given by a two-loop calculation like the one we perform in this paper. At two loops, radiative corrections to $`m_{h^0}^2`$ depend not only on the large top-Yukawa coupling but also on the QCD coupling $`g_3`$. It is reasonable to expect that the dominant two-loop corrections will be of order $`𝒪(\alpha _s\alpha _t[\mathrm{ln}(m_{\stackrel{~}{t}}^2/m_t^2)]^k)`$ and $`𝒪(\alpha _t^2[\mathrm{ln}(m_{\stackrel{~}{t}}^2/m_t^2)]^k),k=0,1,2`$. Terms with $`k=2`$ are the two-loop leading logarithmic contributions and can be obtained by RG techniques using one-loop RG equations; no true two-loop calculation is required and RG resummation will take into account such leading-logarithmic (LL) corrections to all loops. The $`k=1`$ terms are the two-loop next-to-leading-logarithmic (NTLL) corrections, which can be obtained (and resummed to all loops) with two-loop RG equations. Finally, the two-loop non-logarithmic terms ($`k=0`$) can be interpreted in the effective theory language as threshold corrections (at the supersymmetric scale set by the mass of the top-squarks) and require a genuine two-loop calculation; they simply cannot be obtained from RG arguments. The status of these higher-loop calculations of the radiatively corrected $`m_{h^0}^2`$ is the following. Higher-order logarithmic corrections were included in studies which used RG techniques almost since the dramatic impact of radiative corrections on $`m_{h^0}`$ was first recognized. Hempfling and Hoang were the first to perform a genuine two-loop calculation of $`m_{h^0}`$ which also included non-logarithmic terms. They computed the dominant two-loop radiative corrections \[to $`𝒪(\alpha _s\alpha _t)`$ and $`𝒪(\alpha _t^2)`$\] in the case $`\mathrm{tan}\beta 1`$ and zero top-squark mixing. Their computation also included the most important logarithmic corrections, which could be alternatively incorporated by RG resummation from one-loop results, as done e.g. in Ref. . In this last paper it was also pointed out that by a judicious choice of the renormalization scale at which to evaluate one-loop corrections, the higher order logarithmic corrections could be automatically taken into account. A similar idea was later implemented in to write down simple analytical approximations for the radiatively corrected $`m_{h^0}^2`$, obtained by iterative integration of RG equations. Besides being limited to a particularly simple value of $`\mathrm{tan}\beta `$, the calculation in Ref. missed the sizable impact of non-zero top-squark mixing in two-loop effects, that is, higher order corrections to the contribution written down in Eq. (2). Such corrections were first included to order $`𝒪(\alpha _s\alpha _t)`$ in the diagrammatic calculation , and by the effective potential method in Ref. . The effect of these corrections is to shift the values of $`X_t`$ that give maximal mixing, change the corresponding Higgs mass by up to $`10`$ GeV<sup>1</sup><sup>1</sup>1This correction is relative to the one-loop mass using on-shell parameters. The size of the correction would be much smaller if running parameters are used in the one-loop formula (1). and introduce an asymmetry in the dependence of $`m_{h^0}`$ with the sign of $`X_t`$. This two-loop top-squark mixing dependent correction was also explicitly isolated recently by the present authors in Ref. , which uses effective potential plus RG techniques. Besides confirming the independent diagrammatic results of Ref. we clarified the relation of these calculations to previous ones (in particular matching results expressed in different renormalization schemes; see also ). We also derived a compact formula for the Higgs mass (in the spirit of ) which took into account the most important radiative corrections, and used RG techniques to include in a compact way two-loop LL and NTLL corrections. With the $`𝒪(\alpha _s\alpha _t)`$ radiative corrections organized in this way, we find that the mixing-dependent genuine two-loop threshold corrections are generally small ($`\stackrel{<}{_{}}3`$ GeV). Nevertheless, the large computing effort just reviewed did not exhaust the potentially important radiative corrections: the two-loop $`𝒪(\alpha _t^2)`$ top-squark-mixing-dependent corrections to $`m_{h^0}`$ remained unknown to this day, while it is clear that they could compete in principle with the $`𝒪(\alpha _s\alpha _t)`$ contributions. The purpose of this paper is to complete the calculation performed in by using effective potential techniques (plus RG techniques) to compute such $`𝒪(\alpha _t^2)`$ contributions for general top-squark mixing parameters and any value of $`\mathrm{tan}\beta `$. The results in this paper can be considered the most complete and accurate approximation to $`m_{h^0}`$ presented in the literature. The structure of the paper is the following: the next Section describes the strategy of our calculation and presents some analytical formulae for $`m_{h^0}`$, obtained in the limit of $`m_{\stackrel{~}{t}}m_t`$. Section 3 goes one step ahead implementing the RG-improvement of such approximations and, in doing so, clarifies the organization of the higher order radiative corrections. This procedure is not only important to provide a clearer physical picture in connection with the effective field theory but also to classify those corrections calculated in Sec. 2 into a numerically dominant and compact part plus smaller finite threshold correction terms. In Section 4 we present our numerical results for the Higgs mass, illustrate the size of the new corrections and check the validity of our analytical approximation formulae. We draw some conclusions in Section 5. Several appendices are devoted to technical details of different aspects of the calculation. Appendix D is worth special mention as it contains the two-loop $`𝒪(\alpha _t^2)`$ MSSM effective potential used as starting point of our calculation and first computed in this paper. ## 2 $`𝒞𝒫`$-even Higgs boson masses to two-loop order The momentum-dependent mass-squared matrix for the $`𝒞𝒫`$-even Higgs bosons of the MSSM in the interaction eigenstate basis $`h_1,h_2`$ is $$_h^2(p^2)=\left[\begin{array}{cc}m_Z^2c_\beta ^2+m_{A^0}^2s_\beta ^2+\mathrm{\Delta }_{11}^2(p^2)& (m_Z^2+m_{A^0}^2)s_\beta c_\beta +\mathrm{\Delta }_{12}^2(p^2)\\ (m_Z^2+m_{A^0}^2)s_\beta c_\beta +\mathrm{\Delta }_{21}^2(p^2)& m_Z^2s_\beta ^2+m_{A^0}^2c_\beta ^2+\mathrm{\Delta }_{22}^2(p^2)\end{array}\right],$$ (3) where $`s_\beta \mathrm{sin}\beta `$ and $`c_\beta \mathrm{cos}\beta `$. The mass parameters $`m_Z`$ and $`m_{A^0}`$ are the (scale-dependent) running masses of the $`Z`$-boson and $`𝒞𝒫`$-odd Higgs boson $`A^0`$; they are related to the on-shell masses $`M_Z`$ and $`M_{A^0}`$ (we use capital letters to distinguish on-shell parameters from running ones) by $$m_Z^2=M_Z^2+\mathrm{Re}\mathrm{\Pi }_{ZZ}^T(M_Z^2),m_{A^0}^2=M_{A^0}^2+\mathrm{Re}\mathrm{\Pi }_{AA}(M_{A^0}^2)s_\beta ^2\frac{T_1}{v_1}c_\beta ^2\frac{T_2}{v_2},$$ (4) where $`\mathrm{\Pi }_{ZZ}^T`$ is the transverse part of the $`Z`$-boson self-energy and $`\mathrm{\Pi }_{AA}`$ the $`A^0`$-boson self-energy, $`T_1,T_2`$ are the tadpoles of the $`𝒞𝒫`$-even (real) fields $`h_1,h_2`$. Their explicit one-loop expressions can be found e.g. in Ref. . In (3), $`\mathrm{\Delta }^2`$ stands for the contributions from radiative corrections. They are $$\mathrm{\Delta }_{ij}^2(p^2)=\mathrm{\Pi }_{ij}(p^2)+\frac{T_i}{v_i}\delta _{ij},i,j=1,2,$$ (5) where $`\mathrm{\Pi }_{ij}`$ is the self-energy matrix of the Higgs fields $`h_1`$ and $`h_2`$. The masses, $`m_{h^0}`$, $`m_{H^0}`$, of the two $`𝒞𝒫`$-even Higgs bosons are then obtained from the real part of the poles of the propagator matrix, $$\mathrm{Det}\left[m_{h^0,H^0}^2\mathrm{𝟏}_h^2(m_{h^0,H^0}^2)\right]=0.$$ (6) The radiatively corrected mixing angle $`\alpha `$ is obtained as the angle of that rotation which diagonalizes $`_h^2`$ (for some choice of $`p^2`$, say $`p^2=m_{h^0}^2`$): $$\mathrm{tan}2\alpha =\frac{2(_h^2)_{12}}{(_h^2)_{11}(_h^2)_{22}}.$$ (7) Computing Higgs boson masses to a certain order of perturbation theory then requires calculating the self-energies and tadpoles in Eqs. (4) and (5) to that order. In the effective potential approach , self-energies and tadpoles can be calculated as derivatives of the Higgs potential $`V`$ according to: $$T_i=\left[\frac{V(h_1,h_2)}{h_i}\right]\|_{h_1=v_1,h_2=v_2},\mathrm{\Pi }_{ij}(0)=\left[\frac{^2V(h_1,h_2)}{h_ih_j}\right]\|_{h_1=v_1,h_2=v_2}.$$ (8) Note that the $`\mathrm{\Pi }`$’s obtained from derivatives of $`V`$ have zero external momentum. In the limit $`M_{A^0}M_Z`$, the lightest $`𝒞𝒫`$-even Higgs state lies along the direction of the breaking in field space , that is, $`\alpha \beta \pi /2+𝒪(m_Z^2/m_A^2)`$, and its radiatively corrected mass has a very simple expression $$M_{h^0}^2=\frac{4m_t^4}{v^2}\left(\frac{d}{dm_t^2}\right)^2V\mathrm{Re}\mathrm{\Pi }_{hh}(m_{h^0}^2)+\mathrm{Re}\mathrm{\Pi }_{hh}(0),$$ (9) which is exact up to corrections of order $`𝒪(m_Z^4/m_{A^0}^2)`$.<sup>2</sup><sup>2</sup>2 This formula can be proved as follows: If $`m_{A^0}^2m_Z^2`$, $`\alpha \beta \pi /2`$, and we can therefore use the approximation $`\mathrm{\Delta }m_{h^0}^2\mathrm{\Delta }_{11}^2c_\beta ^2+\mathrm{\Delta }_{22}^2s_\beta ^2+2\mathrm{\Delta }_{12}^2s_\beta c_\beta `$, up to higher order terms in $`m_Z^4/m_{A^0}^2`$. Observing that the potential $`V`$ depends on the fields $`h_1`$ and $`h_2`$ only through (field-dependent) top quark mass and the off-diagonal elements of the top-squark mass-squared matrix and using (8), we can easily express the partial derivatives of $`V`$ in terms of the total derivative in (9). A similar formula was already used in . In Eq. (9), $`V`$ is the projection of $`V(h_1,h_2)`$ along the light Higgs $`h=h_1c_\beta +h_2s_\beta `$: $`V(h)=V(h_1hc_\beta ,h_2hs_\beta )`$. Then $`V(h)`$ can be expressed as a function of $`m_t`$ using $`hm_t\sqrt{2}/(h_ts_\beta )`$. We decompose $`V`$ in its $`n^{th}`$-loop pieces $`V_n`$ (explicitly given in Appendix D) as $`V=V_0+V_1+V_2`$. The tree-level part $`V_0`$ is the only one in which we keep non-zero electroweak gauge couplings. We approximate the one-loop part $`V_1`$ by its $`𝒪(\alpha _t)`$ piece coming from top quark/squark loops. The two-loop part $`V_2`$ is approximated by $`V_{2s}+V_{2t}`$, where $`V_{2s}`$ is the $`𝒪(\alpha _s\alpha _t)`$ part and $`V_{2t}`$ the $`𝒪(\alpha _t^2)`$ one. Next, $`\mathrm{\Pi }_{hh}(p^2)`$ is the light Higgs self-energy at external momentum $`p`$, related to the self-energies of $`h_{1,2}`$ by $$\mathrm{\Pi }_{hh}(p^2)\mathrm{\Pi }_{11}(p^2)s_\alpha ^2+\mathrm{\Pi }_{22}(p^2)c_\alpha ^22\mathrm{\Pi }_{12}(p^2)s_\alpha c_\alpha .$$ (10) Notice that the self-energy difference in (9) involves non-zero external momentum and would require a diagrammatic two-loop calculation. However, throughout this paper we work in the approximation of neglecting in the radiative corrections all couplings other than $`h_t`$ or $`g_3`$. In that case, realizing that at tree level $`m_{h^0}`$ depends only on electroweak gauge couplings while its dependence on $`h_t`$ appears only at one-loop, we can write $$\mathrm{\Pi }_{hh}(m_{h^0}^2)\mathrm{\Pi }_{hh}(0)m_{h^0}^2\frac{d}{dp^2}\mathrm{\Pi }_{hh}(p^2)\|_{p^2=0},$$ (11) which gives rise to $`𝒪(\alpha _t^2)`$ contributions from the one-loop $`𝒪(\alpha _t)`$ self-energy $`\mathrm{\Pi }_{hh}`$. In Section 4 we present the numerical results of such procedure for the two-loop potential with $`𝒪(\alpha _s\alpha _t)`$ and $`𝒪(\alpha _t^2)`$ corrections included. The general expression for the $`𝒪(\alpha _s\alpha _t)`$ potential was first computed in while the complete $`𝒪(\alpha _t^2)`$ terms were still missing. Both contributions to $`V`$ are given in Appendix D. It is useful, both for a better understanding of the numerical results and for practical applications, to derive an analytical expression for the light Higgs mass in the case of a large hierarchy between the supersymmetric scale and the electroweak scale (say when the SUSY scale is of order 1 TeV). Such limit is interesting because it maximizes the radiative corrections to $`m_{h^0}^2`$ (so that it corresponds to the most pessimistic scenario for Higgs searches; the case one should be able to discard to rule out the MSSM), and at the same time simplifies the structure of the radiative corrections, avoiding the proliferation of a multitude of different supersymmetric thresholds. We consequently assume now that all supersymmetric particles have roughly the same mass $`M_SM_Z`$. In more detail, focusing on the particles relevant for the radiative corrections to $`m_{h^0}`$, we take equal soft masses $`M_{\stackrel{~}{Q}}=M_{\stackrel{~}{U}}=M_S`$ for the top-squarks (with diagonal masses $`m_{\stackrel{~}{t}}^2M_S^2+m_t^2`$). The two eigenvalues and mixing angle of the top-squark squared-mass matrix are then $$m_{\stackrel{~}{t}_1}^2=m_{\stackrel{~}{t}}^2+m_tX_t,m_{\stackrel{~}{t}_2}^2=m_{\stackrel{~}{t}}^2m_tX_t,s_t^2=c_t^2=\frac{1}{2},$$ (12) We also take the same mass $`M_S`$ for the gluino and the pseudoscalar Higgs \[this means in particular that we can use Eq. (9) for the light Higgs boson\]. In principle we admit the possibility that the $`\mu `$ parameter could be smaller than $`M_S`$, in which case we expect that one chargino and two neutralinos will have masses $`\|\mu \|`$ below the common supersymmetric threshold. In this situation, which broadly corresponds to the case of a common heavy SUSY scale, we find that, using the operator $$𝒟_m^2\frac{4m_t^4}{v^2}\left(\frac{d}{dm_t^2}\right)^2,$$ (13) the different parts entering (9) are $`𝒟_m^2V_0`$ $`=`$ $`m_Z^2\mathrm{cos}^22\beta ,`$ (14) $`𝒟_m^2V_1`$ $`=`$ $`{\displaystyle \frac{3m_t^4}{2\pi ^2v^2}}\left(\mathrm{ln}{\displaystyle \frac{M_S^2}{m_t^2}}+\widehat{X}_t^2{\displaystyle \frac{\widehat{X}_t^4}{12}}\right),`$ (15) $`𝒟_m^2V_{2s}`$ $`=`$ $`{\displaystyle \frac{\alpha _sm_t^4}{\pi ^3v^2}}\{\mathrm{ln}^2{\displaystyle \frac{M_S^2}{m_t^2}}2\mathrm{ln}^2{\displaystyle \frac{M_S^2}{Q^2}}+2\mathrm{ln}^2{\displaystyle \frac{m_t^2}{Q^2}}+\mathrm{ln}{\displaystyle \frac{m_t^2}{Q^2}}1+(1+2\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}+2\mathrm{ln}{\displaystyle \frac{m_t^2}{Q^2}})\widehat{X}_t`$ (16) $`+(12\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}})(\widehat{X}_t^2+{\displaystyle \frac{\widehat{X}_t^3}{3}}){\displaystyle \frac{\widehat{X}_t^4}{12}}\},`$ $`𝒟_m^2V_{2t}`$ $`=`$ $`{\displaystyle \frac{3\alpha _tm_t^4}{16\pi ^3v^2}}\{9\mathrm{ln}^2{\displaystyle \frac{M_S^2}{Q^2}}6\mathrm{ln}{\displaystyle \frac{m_t^2}{Q^2}}\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}3\mathrm{ln}^2{\displaystyle \frac{m_t^2}{Q^2}}+2[3f_2(\widehat{\mu })3f_1(\widehat{\mu })8]\mathrm{ln}{\displaystyle \frac{M_S^2}{m_t^2}}`$ (17) $`+`$ $`6\widehat{\mu }^2\left(1\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}\right)2(4+\widehat{\mu }^2)f_1(\widehat{\mu })+4f_3(\widehat{\mu }){\displaystyle \frac{\pi ^2}{3}}`$ $`+`$ $`\left[(33+6\widehat{\mu }^2)\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}106\widehat{\mu }^24f_2(\widehat{\mu })+(46\widehat{\mu }^2)f_1(\widehat{\mu })\right]\widehat{X}_t^2`$ $`+`$ $`\left[4(7+\widehat{\mu }^2)\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}+23+4\widehat{\mu }^2+2f_2(\widehat{\mu })2(12\widehat{\mu }^2)f_1(\widehat{\mu })\right]{\displaystyle \frac{\widehat{X}_t^4}{4}}`$ $`+`$ $`{\displaystyle \frac{1}{2}}s_\beta ^2\widehat{X}_t^6(\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}1)+c_\beta ^2[3\mathrm{ln}^2{\displaystyle \frac{M_S^2}{m_t^2}}+7\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}4\mathrm{ln}{\displaystyle \frac{m_t^2}{Q^2}}3+60K+{\displaystyle \frac{4\pi ^2}{3}}`$ $`+`$ $`\left(1224K18\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}\right)\widehat{X}_t^2\left(3+16K3\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}\right)(4\widehat{X}_t\widehat{Y}_t+\widehat{Y}_t^2)`$ $`+`$ $`\left(6+{\displaystyle \frac{11}{2}}\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}\right)\widehat{X}_t^4+\left(4+16K2\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}\right)\widehat{X}_t^3\widehat{Y}_t`$ $`+`$ $`({\displaystyle \frac{14}{3}}+24K3\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}})\widehat{X}_t^2\widehat{Y}_t^2({\displaystyle \frac{19}{12}}+8K{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}})\widehat{X}_t^4\widehat{Y}_t^2]\},`$ The notations used are $`\widehat{X}_t=\widehat{A}_t+\widehat{\mu }\mathrm{cot}\beta `$, $`\widehat{Y}_t=\widehat{A}_t\widehat{\mu }\mathrm{tan}\beta `$, with reduced parameters $`\widehat{z}z/M_S`$, and (see Appendix A) $`K0.1953256`$. We also use the following non-singular functions of $`\widehat{\mu }`$ $`f_1(\widehat{\mu })`$ $`=`$ $`{\displaystyle \frac{\widehat{\mu }^2}{1\widehat{\mu }^2}}\mathrm{ln}\widehat{\mu }^2,`$ $`f_2(\widehat{\mu })`$ $`=`$ $`{\displaystyle \frac{1}{1\widehat{\mu }^2}}\left[1+{\displaystyle \frac{\widehat{\mu }^2}{1\widehat{\mu }^2}}\mathrm{ln}\widehat{\mu }^2\right],`$ $`f_3(\widehat{\mu })`$ $`=`$ $`{\displaystyle \frac{(1+2\widehat{\mu }^2+2\widehat{\mu }^4)}{(1\widehat{\mu }^2)^2}}\left[\mathrm{ln}\widehat{\mu }^2\mathrm{ln}(1\widehat{\mu }^2)+Li_2(\widehat{\mu }^2){\displaystyle \frac{\pi ^2}{6}}\widehat{\mu }^2\mathrm{ln}\widehat{\mu }^2\right],`$ (18) with $`f_1(0)=0`$, $`f_2(0)=1`$, $`f_3(0)=\pi ^2/6`$ and $`f_1(1)=1`$, $`f_2(1)=1/2`$, $`f_3(1)=9/4`$. Finally, the correction for non-zero external momentum in Eq. (9) is given by (see Appendix C) $$\mathrm{Re}\left[\mathrm{\Pi }_{hh}(m_{h^0}^2)+\mathrm{\Pi }_{hh}(0)\right]=\frac{h_t^2}{16\pi ^2}m_{h^0}^2s_\beta ^2\left(3\mathrm{ln}\frac{m_t^2}{Q^2}+2\frac{\widehat{X}_t^2}{2}\right).$$ (19) The parameters that appear in these expressions are running parameters, evaluated in the $`\overline{\mathrm{DR}}`$ -scheme and satisfy MSSM RG equations. In fact, it can be checked that the (physical) Higgs mass, given by Eq. (9), is renormalization-scale independent (up to two-loop order), as it should. This scale independence is at the root of the RG-resummation procedure discussed in the next section. It is evident that for different values of the renormalization scale, the magnitude of the two-loop corrections will change, so that it should be possible to choose the scale in such a way that the bulk of the corrections is transferred to the one-loop terms (which depend on the scale implicitly). Therefore, the magnitude and relevance of the two-loop corrections depends on the definition of the mass parameters that enter the one-loop corrections. It is in this respect convenient to write down the two-loop expressions just obtained in the particular case in which all mass parameters in the one-loop correction are the OS ones. This is also useful to compare with explicit diagrammatic calculations. The relationships between running and OS parameters are listed in Appendix C. Using them, we obtain for the two-loop correction to $`m_{h^0}^2`$: $`\mathrm{\Delta }m_{h^0}^2`$ $`=`$ $`{\displaystyle \frac{\alpha _sm_t^4}{\pi ^3v^2}}\left\{3\mathrm{ln}^2{\displaystyle \frac{M_S^2}{m_t^2}}6\mathrm{ln}{\displaystyle \frac{M_S^2}{m_t^2}}+6\widehat{X}_t3\mathrm{ln}{\displaystyle \frac{M_S^2}{m_t^2}}\widehat{X}_t^2{\displaystyle \frac{3}{4}}\widehat{X}_t^4\right\}`$ (20) $`+`$ $`{\displaystyle \frac{3\alpha _tm_t^4}{16\pi ^3v^2}}\{(3\mathrm{ln}^2{\displaystyle \frac{M_S^2}{m_t^2}}+13\mathrm{ln}{\displaystyle \frac{M_S^2}{m_t^2}})s_\beta ^2{\displaystyle \frac{7}{2}}{\displaystyle \frac{\pi ^2}{3}}3\widehat{\mu }^2(11\widehat{\mu }^2+3\widehat{\mu }^4)f_1(\widehat{\mu })`$ $`3(1\widehat{\mu }^2)^2\mathrm{ln}(1\widehat{\mu }^2)+3f_2(\widehat{\mu })+4f_3(\widehat{\mu })+c_\beta ^2\left(60K+{\displaystyle \frac{13}{2}}+{\displaystyle \frac{4\pi ^2}{3}}\right)`$ $`+\left[3s_\beta ^2\mathrm{ln}{\displaystyle \frac{M_S^2}{m_t^2}}+{\displaystyle \frac{73}{2}}+9\widehat{\mu }^2+f_1(\widehat{\mu })7f_2(\widehat{\mu })c_\beta ^2\left({\displaystyle \frac{69}{2}}+24K\right)\right]\widehat{X}_t^2`$ $`+{\displaystyle \frac{1}{6}}\left[269\widehat{\mu }^2+3f_1(\widehat{\mu })+3f_2(\widehat{\mu })+{\displaystyle \frac{61}{2}}c_\beta ^2\right]\widehat{X}_t^4+{\displaystyle \frac{s_\beta ^2}{2}}\widehat{X}_t^6`$ $`+3(1\widehat{\mu }^2)\left[(23\widehat{\mu }^2)f_1(\widehat{\mu })(1+3\widehat{\mu }^2)\mathrm{ln}(1\widehat{\mu }^2)\right]\left(\widehat{X}_t^2{\displaystyle \frac{\widehat{X}_t^4}{6}}\right)`$ $`+c_\beta ^2[(316K\pi \sqrt{3})(4\widehat{X}_t\widehat{Y}_t+\widehat{Y}_t^2)+(16K+{\displaystyle \frac{2\pi }{\sqrt{3}}})\widehat{X}_t^3\widehat{Y}_t`$ $`+({\displaystyle \frac{4}{3}}+24K+\pi \sqrt{3})\widehat{X}_t^2\widehat{Y}_t^2({\displaystyle \frac{7}{12}}+8K+{\displaystyle \frac{\pi }{2\sqrt{3}}})\widehat{X}_t^4\widehat{Y}_t^2]`$ $`+(2\widehat{X}_t{\displaystyle \frac{\widehat{X}_t^3}{3}})[(3+{\displaystyle \frac{2\pi }{\sqrt{3}}})c_\beta ^2\widehat{X}_t\widehat{Y}_t^2(3\mathrm{ln}{\displaystyle \frac{m_tX_t}{M_S^2}}+\mathrm{ln}4)s_\beta ^2\widehat{X}_t^3]\}`$ We emphasize that this expression gives the two-loop corrections when the one-loop contribution (15) is expressed in terms of OS parameters, that is, $$\left[\mathrm{\Delta }m_{h^0}^2\right]_{1loop}^{\mathrm{OS}}=\frac{3g^2M_t^4}{8\pi ^2M_W^2}\left[\mathrm{ln}\frac{M_{\stackrel{~}{t}}^2}{M_t^2}+\left(\frac{X_t^{\mathrm{OS}}}{M_{\stackrel{~}{t}}}\right)^2\frac{1}{12}\left(\frac{X_t^{\mathrm{OS}}}{M_{\stackrel{~}{t}}}\right)^4\right].$$ (21) Several features of Eq. (20) are worth commenting. First, if we restrict Eq. (20) to $`\mathrm{tan}\beta 1`$ and zero $`A_t`$, to compare with the result of Ref. , we find the same logarithmic terms. However, the $`𝒪(\alpha _t^2)`$ finite term is different. In particular, that term is sensitive to the value of the parameter $`\mu `$, contrary to what is stated in Ref. . Nevertheless, the result quoted for that finite term in Ref. is inside the range we would find by varying $`\widehat{\mu }^2`$ from 0 to 1, and the impact of this $`\mu `$-dependence on the final Higgs mass is quite small. Second, we see that radiative corrections no longer depend on $`A_t`$ and $`\mu `$ in the combination $`X_t`$ that appears through the off-diagonal entry of the top-squark mass matrix: besides the explicit dependence on the parameter $`\mu `$ already noticed, the quantity $`Y_t`$ also introduces a different combination of $`A_t`$ and $`\mu `$. This dependence on $`Y_t`$ originates from the $`H\stackrel{~}{t}\stackrel{~}{b}`$ and $`H\stackrel{~}{t}\stackrel{~}{t}`$ diagrams of Fig. 8. Third, although roughly speaking the top-Yukawa correction has a small pre-factor $`3/16`$ in comparison with the QCD correction, this does not guarantee that the new contributions will be negligible compared to the QCD one. In fact, we will see that for two-loop top-squark-mixing-dependent corrections of (20), the top Yukawa contributions have opposite signs as that of the QCD corrections and could be as much as $`60\%`$ of the latter (see Fig. 6). In the next Section, we will follow RG methods and reorganize these corrections in the effective theory language, with the most important corrections of Eq. (20) reshuffled in a RG-motivated one-loop formula. ## 3 Renormalization group resummation Before illustrating in Section 4 the impact of the newly computed corrections on the Higgs mass, we show in the following how the use of renormalization group techniques allows us to write the previous complicated corrections \[see Eq. (20)\] in a simpler and more transparent way, while at the same time it clarifies the connection to the RG programme, which can be used to improve the precision of the mass formula by resummation of higher order corrections. We already applied this idea in Ref. to the $`𝒪(\alpha _s\alpha _t)`$ two-loop corrections. By a convenient (and physically well motivated) choice of the scale at which to evaluate running parameters in the one-loop mass correction one can absorb large logarithms in Eq. (20). The RG evolution of the parameters is given by the corresponding one-loop RG functions listed in Appendix B. We use the following equations to relate supersymmetric running parameters at different scales \[cf. Eqs. (B.23) and (B.24)\]: $$m_{\stackrel{~}{t}}^2(Q)=m_{\stackrel{~}{t}}^2(Q^{})\left\{1+\frac{1}{16\pi ^2}\left[\frac{16}{3}g_3^2\frac{3}{2}h_t^2\left(\widehat{X}_t^2s_\beta ^2+\widehat{Y}_t^2c_\beta ^2+c_\beta ^2+22\widehat{\mu }^2\right)\right]\mathrm{ln}\frac{Q^{}^2}{Q^2}\right\},$$ (22) $$X_t(Q)=X_t(Q^{})\frac{1}{16\pi ^2}\left[\frac{16}{3}g_3^2M_S+3h_t^2(X_t+X_ts_\beta ^2+Y_tc_\beta ^2)\right]\mathrm{ln}\frac{Q^{}^2}{Q^2},$$ (23) where we have used $`A_t=X_ts_\beta ^2+Y_tc_\beta ^2`$, $`A_t^2+\mu ^2=X_t^2s_\beta ^2+Y_t^2c_\beta ^2`$ and $`m_{H_2}^2+\mu ^2=m_{A^0}^2c_\beta ^2`$. Notice that, to the order we work, it is sufficient to use these one-loop LL approximations to the full RG evolution because we are concerned with parameters that appear in a one-loop order term. The Standard Model $`\overline{\mathrm{MS}}`$ top quark mass $`\overline{m}_t`$ and the Higgs VEV $`\overline{v}`$ are related to the on-shell mass $`M_t`$ and MSSM VEV $`v`$ by \[cf. Eqs. (Appendix C: One-loop self-energies) and (C.10), from which relevant terms can be easily identified\] $$\overline{m}_t^2(Q)=M_t^2\left[1\frac{g_3^2}{6\pi ^2}\left(43\mathrm{ln}\frac{m_t^2}{Q^2}\right)+\frac{h_t^2s_\beta ^2}{32\pi ^2}\left(83\mathrm{ln}\frac{m_t^2}{Q^2}\right)\right],$$ (24) $$\overline{v}^2(Q)=v^2(Q)\left[1+\frac{h_t^2s_\beta ^2}{32\pi ^2}\widehat{X}_t^2\right].$$ (25) We also use one-loop LL solutions of the SM RG equations to relate these parameters at different scales: $$\overline{m}_t^2(Q)=\overline{m}_t^2(Q^{})\left[1+\frac{1}{16\pi ^2}\left(8g_3^2\frac{3}{2}h_t^2s_\beta ^2\right)\mathrm{ln}\frac{Q^{}^2}{Q^2}\right],$$ (26) $$\overline{v}^2(Q)=\overline{v}^2(Q^{})\left[1+\frac{3h_t^2s_\beta ^2}{16\pi ^2}\mathrm{ln}\frac{Q^{}^2}{Q^2}\right].$$ (27) Using the above equations, we find the following compact expression for the Higgs boson mass, which is one of the main results of this paper $$M_{h^0}^2=M_Z^2\mathrm{cos}^22\beta +\frac{3\overline{m}_t^4(Q_t)}{2\pi ^2\overline{v}^2(Q_1^{})}\mathrm{ln}\frac{m_{\stackrel{~}{t}}^2(Q_{\stackrel{~}{t}})}{\overline{m}_t^2(Q_t^{})}+\mathrm{\Delta }_{\mathrm{th}}^{(1)}m_{h^0}^2+\mathrm{\Delta }_{\mathrm{th}}^{(2)}m_{h^0}^2.$$ (28) The one-loop threshold correction is $$\mathrm{\Delta }_{\mathrm{th}}^{(1)}m_{h^0}^2=\frac{3\overline{m}_t^4(Q_{\mathrm{th}})}{2\pi ^2\overline{v}^2(Q_2^{})}\left[\widehat{X}_t^2(Q_{\mathrm{th}})\frac{\widehat{X}_t^4(Q_{\mathrm{th}})}{12}\right],$$ (29) and the two-loop threshold correction reads $`\mathrm{\Delta }_{\mathrm{th}}^{(2)}m_{h^0}^2`$ $`=`$ $`{\displaystyle \frac{\alpha _sm_t^4}{\pi ^3v^2}}\left[2\widehat{X}_t\widehat{X}_t^2+{\displaystyle \frac{7}{3}}\widehat{X}_t^3+{\displaystyle \frac{1}{12}}\widehat{X}_t^4{\displaystyle \frac{1}{6}}\widehat{X}_t^5\right]`$ (30) $`+`$ $`{\displaystyle \frac{3\alpha _tm_t^4}{16\pi ^3v^2}}\{R_0(\widehat{\mu })+R_2(\widehat{\mu })\widehat{X}_t^2+R_4(\widehat{\mu })\widehat{X}_t^4{\displaystyle \frac{1}{2}}s_\beta ^2\widehat{X}_t^6`$ $`+`$ $`c_\beta ^2[60K{\displaystyle \frac{9}{2}}+{\displaystyle \frac{4\pi ^2}{3}}(3+16K)(4\widehat{X}_t\widehat{Y}_t+\widehat{Y}_t^2)+(1524K)\widehat{X}_t^2`$ $``$ $`{\displaystyle \frac{25}{4}}\widehat{X}_t^4+(4+16K)\widehat{X}_t^3\widehat{Y}_t+({\displaystyle \frac{14}{3}}+24K)\widehat{X}_t^2\widehat{Y}_t^2({\displaystyle \frac{19}{12}}+8K)\widehat{X}_t^4\widehat{Y}_t^2]\}.`$ We have used the short-hand notation $`R_0(\widehat{\mu })`$ $`=`$ $`{\displaystyle \frac{9}{2}}{\displaystyle \frac{\pi ^2}{3}}+6\widehat{\mu }^2(11+2\widehat{\mu }^2)f_1(\widehat{\mu })+3f_2(\widehat{\mu })+4f_3(\widehat{\mu }),`$ $`R_2(\widehat{\mu })`$ $`=`$ $`11\widehat{\mu }^2[6+6f_1(\widehat{\mu })+10f_2(\widehat{\mu })],`$ $`R_4(\widehat{\mu })`$ $`=`$ $`6+\widehat{\mu }^2[1+f_1(\widehat{\mu })+f_2(\widehat{\mu })].`$ (31) The scales required in (28,29) are $$Q_t=\sqrt{m_tm_{\stackrel{~}{t}}},Q_t^{}=(m_tm_{\stackrel{~}{t}}^2)^{1/3},Q_{\stackrel{~}{t}}=Q_{\mathrm{th}}=m_{\stackrel{~}{t}},$$ $$Q_1^{}=e^{1/3}m_t0.7m_t,Q_2^{}=e^{1/3}m_t1.4m_t.$$ (32) It is a non-trivial check of our calculation that the values of the scales (32) required to re-absorb the large $`\mathrm{ln}(M_S^2/m_t^2)`$ logarithms in the two-loop corrections are consistent with the ones obtained in for the QCD corrections alone. We see, in particular, that the uncertainty found there in the definition of the scales $`Q_t^{}`$ and $`Q_{\stackrel{~}{t}}`$ is here resolved by the need of absorbing the new radiative corrections. We still find a somewhat complicated expression for the threshold correction $`\mathrm{\Delta }_{\mathrm{th}}^{(2)}m_{h^0}^2`$, due to the fact that we have kept free the $`\mu `$ parameter. Expressions for the two limiting cases of heavy $`\mu `$ ($`\mu M_S`$) and light $`\mu `$ ($`\mu M_s`$) can be readily derived. In both cases, the resulting threshold correction is much simpler than the general case (30) and contains no more logarithms. Explicitly, for $`\mu M_s`$ we find $$R_0(0)=\frac{\pi ^2}{3}\frac{3}{2},R_2(0)=11,R_4(0)=6,$$ (33) and for $`\mu M_S`$: $$R_0(1)=7\frac{\pi ^2}{3},R_2(1)=16,R_4(1)=\frac{13}{2}.$$ (34) It is perhaps convenient to make more explicit the connection between our results and those obtained in the RG approach (see, e.g., Ref. ). To be concrete, let us assume $`\|\mu \|=M_S`$ so that all supersymmetric particles (including charginos and neutralinos) have masses of order $`M_S`$; below that scale, the effective theory is the SM. The light Higgs quartic coupling $`\lambda `$ at $`M_S`$ consists of a tree-level part plus higher-order threshold corrections which arise from the heavy decoupling supersymmetric particles, it can be evolved down to the electroweak scale, say $`Q=m_t`$, using the SM RGEs; at that scale $`\lambda `$ is related to the physical Higgs mass. This procedure should reproduce all the logarithmic corrections we have found. More explicitly, defining $`\beta _\lambda d\lambda /d\mathrm{ln}Q^2`$, we can write $$\lambda (𝒬_t)=\lambda (𝒬_{\stackrel{~}{t}})_{Q=𝒬_t}^{𝒬_{\stackrel{~}{t}}}\beta _\lambda d\mathrm{ln}Q^2.$$ (35) We use a special notation for the high and low scales between which we run $`\lambda `$ to distinguish them from other definitions of $`m_t`$ and $`m_{\stackrel{~}{t}}`$ that appear in the paper. These quantities are defined by: $$𝒬_t\overline{m}_t(𝒬_t),𝒬_{\stackrel{~}{t}}m_{\stackrel{~}{t}}(𝒬_{\stackrel{~}{t}}),$$ (36) i.e., they are the running masses evaluated at a scale equal to the corresponding mass. This is the natural definition in the RG approach. Making a loop expansion of $`\beta _\lambda `$ in (35) and a further expansion around a particular value of $`Q`$ (say the low energy limit of the running interval, $`𝒬_t`$), we obtain to the two-loop order $$\lambda (𝒬_t)\lambda (𝒬_{\stackrel{~}{t}})\left[\beta _\lambda ^{(1)}(𝒬_t)+\beta _\lambda ^{(2)}(𝒬_t)\right]\mathrm{ln}\frac{𝒬_{\stackrel{~}{t}}^2}{𝒬_t^2}\frac{1}{2}\frac{d\beta _\lambda ^{(1)}}{d\mathrm{ln}Q^2}(𝒬_t)\mathrm{ln}^2\frac{𝒬_{\stackrel{~}{t}}^2}{𝒬_t^2}+\mathrm{}$$ (37) where the one- and two-loop contributions to $`\beta _\lambda `$ are approximated by \[neglected all couplings other than the strong gauge coupling $`g_3`$ and the SM top Yukawa coupling $`g_t`$ ($`h_ts_\beta `$)\] $`\beta _\lambda ^{(1)}`$ $`=`$ $`{\displaystyle \frac{3g_t^2}{8\pi ^2}}(g_t^2+\lambda ),`$ $`\beta _\lambda ^{(2)}`$ $`=`$ $`{\displaystyle \frac{2g_t^4}{(16\pi ^2)^2}}(15g_t^216g_3^2).`$ (38) We note that for a correct two-loop computation it is necessary to retain also the $`\lambda g_t^2`$ term in $`\beta _\lambda ^{(1)}`$ because $`\lambda `$ gets one-loop contributions proportional to $`g_t^4`$. $`d\beta _\lambda ^{(1)}/d\mathrm{ln}Q^2`$ can be calculated from the one-loop RG evolution of $`g_t`$ $$\frac{dg_t^2}{d\mathrm{ln}Q^2}=\frac{g_t^2}{32\pi ^2}(9g_t^216g_3^2).$$ (39) Once $`\lambda (𝒬_t)`$ is obtained from (37), we extract the physical Higgs mass using the SM relation : $$\lambda (𝒬_t)\overline{v}^2(𝒬_t)=M_{h^0}^2\left(1\frac{g_t^2}{8\pi ^2}\right).$$ (40) This correction arises from wave-function renormalization and takes into account the fact that the physical mass is defined on-shell, and not at zero external momentum. Its physical content is therefore similar to the correction (11) in our effective potential approach. According to (37), the large LL and NTLL corrections to $`M_{h^0}^2`$ arise solely from $`\lambda (𝒬_t)`$. Additional radiative contributions in (40), coming from $`\overline{v}^2(𝒬_t)`$ and the wave-function correction factor, affect the large logarithmic terms only through multiplication of $`\lambda (𝒬_t)`$. It is therefore clear that it is enough for our purposes to know these correction factors at one-loop order. Based on this observation, we can combine both factors together using (25) to write the simpler formula $$M_{h^0}^2=\lambda (𝒬_t)\overline{v}^2(𝒬_1^{}),$$ (41) with $`𝒬_1^{}=e^{1/3}𝒬_t`$, in accordance with (32). It is now straightforward to show perfect agreement of $`M_{h^0}`$ as obtained from the above expression with our results (28-30). All logarithmic corrections up to two-loops are exactly reproduced while the finite part agrees if one uses as boundary condition at the SUSY scale $$\lambda (𝒬_{\stackrel{~}{t}})=\frac{1}{4}(g_1^2+g_2^2)\mathrm{cos}^22\beta +\delta _{\mathrm{th}}^{(1)}\lambda +\delta _{\mathrm{th}}^{(2)}\lambda ,$$ (42) with $`\delta _{\mathrm{th}}^{(1)}\lambda `$ $`=`$ $`{\displaystyle \frac{3g_t^4(Q_{\mathrm{th}})}{8\pi ^2}}\left[\widehat{X}_t^2(Q_{\mathrm{th}}){\displaystyle \frac{\widehat{X}_t^4(Q_{\mathrm{th}})}{12}}\right],`$ $`\delta _{\mathrm{th}}^{(2)}\lambda `$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_{\mathrm{th}}^{(2)}m_{h^0}^2}{v^2}}.`$ (43) To summarize, we find full agreement between our approximate formula (28) for the Higgs boson mass and the RG-improved mass calculated in the RG (or effective theory) approach, to two-loop order. The connection to the effective theory language clarifies the origin of the different terms in (28), and rewrites them in a very convenient way, absorbing the large (logarithmic) two-loop effects in the one-loop correction and leaving behind two-loop threshold corrections which are numerically small, as we will see in the next Section. Note that this applies in particular to the sizable top-squark-mixing-dependent corrections of Eq. (20), the bulk of which is transferred to the RG-reshuffled one-loop threshold correction of Eq. (29). Knowing the boundary condition, $`\lambda (𝒬_{\stackrel{~}{t}})`$, one can integrate (35) numerically by solving a coupled set of differential equations (describing the two-loop evolution of $`\lambda `$, $`g_3`$, $`g_t`$), find $`\lambda (𝒬_t)`$ and use (41) to get the Higgs mass. The final result will be the full RG-improved value of $`M_{h^0}`$ and will resum LL and NTLL corrections to all loops \[numerical integration includes all the terms from the expansion around $`𝒬_t`$ which were neglected in (37)\]. In this respect, note that our compact formula, Eq. (28), which has been found by requiring that logarithmic contributions are correctly reproduced up to two-loops only, contains in fact logarithmic corrections of higher order. It can be shown that these higher order logarithmic corrections do not match exactly the correct ones (obtained by the RG method) if we use simple one-loop approximations \[like those given in Eqs. (24,25)\] to evaluate the parameters in (28) at their corresponding scales. However, evaluation of those parameters by means of a full numerical integration \[similar to that in Eq. (35) for $`\lambda `$\] would correctly take into account the LL (but not the NTLL) terms to all loops. Nevertheless, as we will see in the next Section, the error made in neglecting logarithmic corrections of higher order is very small for SUSY scales of interest \[below $`M_S𝒪(1)`$ TeV\]. If $`M_S`$ turns out to be significantly larger than that (starting to be in conflict with naturalness criteria), then one should revert to the numerical RG integration of $`\lambda `$ to get a reliable estimate of the Higgs mass. Our results for the boundary condition $`\lambda (𝒬_{\stackrel{~}{t}})`$ will still be useful in such a case. ## 4 Numerical results In this section we present numerical results from our two-loop study. For the one-loop analysis we closely follow Ref. , which has included complete radiative corrections from the dominant top quark/squark sector and the sub-dominant gauge/Higgs boson and neutralino/chargino sectors. In what follows, we shall concentrate on two-loop radiative corrections. We start by sketching the procedure for this analysis, which is the following: we first take as inputs the on-shell mass parameters<sup>3</sup><sup>3</sup>3For the top-squark sector, we can alternatively take as inputs the on-shell top-squark masses and mixing angle. $`M_{A^0}`$, $`M_t`$, $`M_{\stackrel{~}{Q}}^{\mathrm{OS}}`$, $`M_{\stackrel{~}{U}}^{\mathrm{OS}}`$ and $`A_t^{\mathrm{OS}}`$. From them we can determine the values of the corresponding running parameters at any renormalization scale $`Q`$. To do this, we have to calculate the one-loop self-energy diagrams for Higgses and top-squarks (the latter are collected in Appendix C). We also input $`\mathrm{tan}\beta `$ and $`\mu `$ parameters, and convert $`\alpha _s(M_Z)=0.118`$ to the MSSM $`\overline{\mathrm{DR}}`$ running value. Next we calculate in the MSSM the two-loop corrections to the $`𝒞𝒫`$-even Higgs mass matrix, $`\mathrm{\Delta }_{ij}^2`$, from the two-loop potential (D.5) and (D.6) using Eq. (8). Numerically, the partial derivatives in these equations are replaced by finite differences in $`h_1,h_2`$, i.e. we vary the values of these fields by a finite amount and recalculate the field-dependent top-quark mass $`m_t^2=\frac{1}{2}h_t^2h_2^2`$ and top-squark masses $`m_{\stackrel{~}{t}_1}`$, $`m_{\stackrel{~}{t}_2}`$, mixing angle $`\theta _{\stackrel{~}{t}}`$ from Eq. (B.5). With these new parameters, the two-loop potential is reevaluated and their variations from the reference values \[calculated at $`h_1^2+h_2^2=(246\mathrm{GeV})^2`$\] are found. Finally, equipped with the corrections $`\mathrm{\Delta }_{ij}^2`$, we compute the lightest $`𝒞𝒫`$-even Higgs boson mass by solving Eq. (6). Several approximations have been made to quantities in Eqs. (3-5), in particular we neglect all dimensionless couplings other than the top-Yukawa coupling $`h_t`$ and the QCD gauge coupling $`g_3`$. In this way we pick up the dominant radiative effects only, what we term throughout leading corrections. We notice that the two-loop self-energy of the $`Z`$-boson and the non-zero external momentum corrections to two-loop Higgs boson self-energies can be neglected in our calculation since all these corrections are higher order effects in the leading approximation. However, we need to calculate $`\mathrm{\Pi }_{AA}`$ to the two-loop level since it has $`𝒪(\alpha _s\alpha _t)`$ and $`𝒪(\alpha _t^2)`$ corrections and in principle could contribute to (4) at the same order as $`\mathrm{\Delta }_{ij}^2`$. It is not possible to obtain these self-energies in our current approach, and explicit two-loop calculation of the corresponding two-point functions are needed. Fortunately, the correction to $`m_{h^0}`$ from $`\mathrm{\Pi }_{AA}`$ is always numerically negligible for large $`m_{A^0}`$ as can be easily seen from the structure of the Higgs mass matrix (3). That is, (9) is correct for large $`m_{A^0}`$ and we can safely neglect the $`\mathrm{\Pi }_{AA}`$ corrections. (For $`m_{A^0}m_Z`$, a complete two-loop calculation of $`m_{h^0}`$ would need $`\mathrm{\Pi }_{AA}`$.) Fig. 1 is used as calibration: we compare in it our numerical results for $`M_{h^0}`$ including only up to two-loop $`𝒪(\alpha _s\alpha _t)`$ corrections with the mass obtained by the program FeynHiggs which uses the explicit two-loop diagrammatic results of Ref. . We choose two sets of parameters<sup>4</sup><sup>4</sup>4We assume $`M_3`$ is positive hereafter. For a negative $`M_3`$ our formulae still apply simply by simultaneous sign changes in $`X_t`$ and $`Y_t`$.: (a) $`M_{A^0}=M_3=M_{\stackrel{~}{Q}}^{\mathrm{OS}}=M_{\stackrel{~}{U}}^{\mathrm{OS}}=M_S=500`$ GeV, $`\mu =200`$ GeV and (b) $`M_{A^0}=M_3=M_{\stackrel{~}{Q}}^{\mathrm{OS}}=M_{\stackrel{~}{U}}^{\mathrm{OS}}=M_S=1`$ TeV, $`\mu =500`$ GeV. For each case, results for two values of $`\mathrm{tan}\beta `$ ($`1.6`$ and $`20`$) are plotted. We find good agreement (given the fact that they are two independent programs) between both one-loop (shown in dotted and dashed lines) and two-loop QCD corrected (shown in dot-dashed and solid lines) masses; this shows numerically that the two approaches are equivalent to that order. This equivalence is easily understood since the effective potential, as a generating functional , encompasses all tadpole and self-energy diagrams (as well as all other multi-point functions) which are calculated in . The effective potential approach is more efficient for the purpose of calculating $`M_{h^0}`$ and much simpler to implement in a Fortran program since it requires evaluating only one set of two-loop functions. In Fig. 2 we show the Higgs boson mass $`M_{h^0}`$ vs. the (on-shell) SUSY scale $`M_S`$, for two values of the top-squark mixing parameters $`\widehat{X}_t^{\mathrm{OS}}`$ ($`0`$ and $`2`$). All the physical masses $`M_{A^0}`$, $`M_{\stackrel{~}{Q}}^{\mathrm{OS}}`$, $`M_{\stackrel{~}{U}}^{\mathrm{OS}}`$ and $`M_3`$ have been set to $`M_S`$ (and the same will be done for all the following plots). The dashed, dot-dashed and solid lines in this figure correspond to masses $`M_{h^0}`$ corrected to one-loop, two-loop $`𝒪(\alpha _s\alpha _t)`$ and two-loop $`𝒪(\alpha _s\alpha _t+\alpha _t^2)`$ order<sup>5</sup><sup>5</sup>5In general, we try to follow the rule that denser lines correspond to more precise approximations.. Fig. 2a ($`\widehat{X}_t^{\mathrm{OS}}=0`$) corresponds to the case of minimal left-right top-squark mixing, and the two-loop $`𝒪(\alpha _t^2)`$ corrections are generally small, $`\stackrel{<}{}2`$ GeV. For Fig. 2b ($`\widehat{X}_t^{\mathrm{OS}}=2`$), which roughly corresponds to the maximal left-right top-squark mixing case, we find that the two-loop $`𝒪(\alpha _t^2)`$ corrections are sizable ($`5`$ GeV). In Fig. 3 we examine the upper limit on the Higgs boson mass $`M_{h^0}`$ by including the dominant two-loop corrections. We show corrected massed to the two-loop $`𝒪(\alpha _s\alpha _t)`$ and $`𝒪(\alpha _s\alpha _t+\alpha _t^2)`$ orders in dot-dashed and solid lines for $`\widehat{X}_t^{\mathrm{OS}}=0,2`$ and the top quark pole mass $`M_t=175`$ GeV. We see that maximal values for $`M_{h^0}`$ of $`129`$ GeV can be reached for large $`\mathrm{tan}\beta `$ and left-right top-squark mixing parameter $`\widehat{X}_t^{\mathrm{OS}}2`$. Without two-loop $`𝒪(\alpha _t^2)`$ corrections, the upper bound of $`M_{h^0}`$ would be at $`124`$ GeV. We also show the Higgs boson masses for $`M_t=180`$ and 170 GeV (including all two-loop dominant corrections) in dashed and dot-dashed lines; the masses are increased or decreased by $``$ 5 GeV respectively. We remark that this upper bound on $`M_{h^0}`$ is asymmetric with respect to $`\widehat{X}_t^{\mathrm{OS}}`$. For $`\widehat{X}_t^{\mathrm{OS}}=2`$ and $`M_t=175`$ GeV, we find the bound is about $`5`$ GeV lower. As is well know, this asymmetry arises from the two-loop $`𝒪(\alpha _s\alpha _t)`$ corrections . In Fig. 4 we compare results from our analytical approximation formula for $`M_{h^0}`$ in Sec. 2 with those obtained by full numerical evaluations. They are shown in dot-dashed and solid lines respectively. The analytical approximation formula works remarkably well: it is good to a precision of $`\stackrel{<}{}0.5`$ GeV for almost all the parameter space. The analytical approximation has a complicated dependence on the $`\mu `$-parameter. Numerically this dependence is quite weak: varying $`\mu `$ from 100 GeV to 1 TeV for a fixed $`\widehat{X}_t^{\mathrm{OS}}`$ changes the Higgs boson mass by less than 1 GeV. We emphasize that the analytical formula is useful for several reasons: (1) the logarithmic and finite corrections can be easily separated, and one can weight the relative importance of these terms; (2) all terms can be traced back to the potential, so one can easily locate the particles giving the biggest contributions; (3) the formula can significantly simplify the numerical evaluations of $`M_{h^0}`$ to a good precision. In Fig. 5 we further compare the results for our RG-corrected Higgs boson masses, Eqs. (28-30), with those of the full numerical evaluation. For comparison, we have also shown two-loop $`𝒪(\alpha _s\alpha _t)`$ corrections and their RG-corrected results; they have been studied previously in . As mentioned in Sect. 3, the good agreement between these curves is an indication of the smallness of the logarithmic corrections beyond two-loops and illustrates the accuracy of our results. Finally in Figs. 6 and 7 we detail the size of the two-loop top-squark-mixing-dependent corrections in the OS-scheme and their corresponding finite threshold corrections in the RG approach. Fig. 6 shows in dotted lines two-loop masses without including the top-squark-mixing-dependent corrections of Eq. (20). Refs. have already calculated the QCD corrections, and they are depicted in dashed lines. The difference of the solid and dashed lines is the two-loop $`𝒪(\alpha _t^2)`$ terms which are calculated in this paper. We see clearly that these terms are sizable: for large mixing parameters, they increase $`M_{h^0}`$ by about 4 GeV and $`23`$ GeV for small and large $`\mathrm{tan}\beta `$ respectively. Fig. 7 shows the effect of two-loop threshold corrections $`\mathrm{\Delta }_{\mathrm{th}}^{(2)}m_{h^0}^2`$ evaluated following the RG-inspired analysis of Sect. 3. The dotted lines show the Higgs boson mass neglecting these corrections; this would have been obtained by integrating two-loop RG equations (with the two-loop boundary threshold correction being set to zero), as we have shown in the second part of Sec. 3. Two-loop masses without the $`𝒪(\alpha _t^2)`$ threshold correction and the complete two-loop results are shown in dashed and solid lines respectively. The RG reshuffling of radiative corrections has absorbed the main part of the two-loop top-squark-mixing-dependent terms of Eq. (20) into the RG-corrected one-loop term $`\mathrm{\Delta }_{\mathrm{th}}^{(1)}m_{h^0}^2`$; the remaining genuine two-loop threshold corrections (in the sense of the effective field theory) are generally small, $`\stackrel{<}{}3`$ GeV. ## 5 Conclusions In this paper we calculate radiative corrections to the lightest MSSM $`𝒞𝒫`$-even Higgs boson mass to the two-loop $`𝒪(\alpha _t^2)`$ order. Our analysis extends existing two-loop diagrammatic results using a simpler effective potential approach and provides the most complete and accurate calculation presented in the literature. We also derive useful analytical approximation formulae, applicable when the supersymmetric particles are heavy, which accurately reproduce results from the full numerical study. Our calculation includes effects which can have an impact on the final Higgs mass but were neglected by previous studies. In particular, the two-loop $`𝒪(\alpha _t^2)`$ top-squark-mixing-dependent corrections to $`M_{h^0}^2`$ \[see Eq. (20)\] are calculated for the first term in this paper and are numerically important. We further simplify our analytical formula by reshuffling higher order logarithmic corrections (using RG techniques) in a compact one-loop expression \[Eq. (28)\]. In that expression all mass parameters are evaluated at appropriate renormalization scales chosen to reproduce the numerically most important leading and next-to-leading logarithmic corrections. The remaining two-loop finite terms can be interpreted as threshold corrections, and are numerically less important. This RG rewriting clarifies the structure of the two-loop corrections to $`M_{h^0}^2`$, identifies the most important contributions and links our work to the effective theory or RG approach, as we have shown in detail in Sec. 3. To summarize our numerical results, we have shown that two-loop top Yukawa corrections to $`M_{h^0}`$ are sizable for the maximal top-squark mixing case. They can increase the Higgs boson mass $`M_{h^0}`$ by as much as $`5`$ GeV (among which the top-squark-mixing-dependent corrections account for about 4 GeV) for small $`\mathrm{tan}\beta `$ where $`h_t`$ is large. The upper bound on $`M_{h^0}`$ is $`129\pm 5`$ GeV for $`M_t=175\pm 5`$ GeV. Our final approximation formulae (20-21) and (28-30) have been shown to excellently agree with the full numerical results and can be easily implemented in precision numerical studies. Although we have focussed in this paper on the Higgs mass, it is worth mentioning that we have also presented in Appendix D the MSSM two-loop effective potential including top-quark Yukawa contributions (for general top-squark mixing parameters and any $`\mathrm{tan}\beta `$). This knowledge may well prove useful for other studies. ## Acknowledgments We thank André Hoang for correspondence. R.-J.Z. was supported in part by a DOE grant No. DE-FG02-95ER40896 and in part by the Wisconsin Alumni Research Foundation. ## Appendix A: One- and two-loop scalar functions ### A.1 One-loop scalar functions In this subsection we define the scalar functions $`A_0`$, $`B_0`$, $`B_1`$, $`B_{22}`$ and $`G`$, which appear in one-loop self-energy calculations. The $`A_0`$ function is defined by the following momentum integral in $`d=42ϵ`$ dimensions $$A_0(m^2)=16\pi ^2\mu ^{4d}\frac{d^dp}{i(2\pi )^d}\frac{1}{p^2m^2+i\epsilon }=m^2\left(\frac{1}{ϵ}+1\mathrm{ln}\frac{m^2}{Q^2}\right),$$ (A.1) where $`Q^2=4\pi \mu ^2e^{\gamma _E}`$ is the renormalization scale, with $`\gamma _E`$ the Euler constant. The $`B_0`$ function is $`B_0(p^2,m_1^2,m_2^2)`$ $`=`$ $`16\pi ^2\mu ^{4d}{\displaystyle \frac{d^dq}{i(2\pi )^d}\frac{1}{[q^2m_1^2+i\epsilon ][(qp)^2m_2^2+i\epsilon ]}}`$ (A.2) $`=`$ $`{\displaystyle \frac{1}{ϵ}}{\displaystyle _0^1}𝑑x\mathrm{ln}{\displaystyle \frac{(1x)m_1^2+xm_2^2x(1x)p^2i\epsilon }{Q^2}}.`$ The remaining functions can be related to $`A_0`$ and $`B_0`$ as follows $`B_1(p^2,m_1^2,m_2^2)`$ $`=`$ $`{\displaystyle \frac{1}{2p^2}}\left[A_0(m_2^2)A_0(m_1^2)+(p^2+m_1^2m_2^2)B_0(p^2,m_1^2,m_2^2)\right],`$ (A.3) $`B_{22}(p^2,m_1^2,m_2^2)`$ $`=`$ $`{\displaystyle \frac{1}{6}}[A_0(m_2^2)+2m_1^2B_0(p^2,m_1^2,m_2^2)(p^2+m_1^2m_2^2)B_1(p^2,m_1^2,m_2^2)`$ (A.4) $`+m_1^2+m_2^2{\displaystyle \frac{p^2}{3}}],`$ $`G(p^2,m_1^2,m_2^2)`$ $`=`$ $`(p^2m_1^2m_2^2)B_0(p^2,m_1^2,m_2^2)A_0(m_1^2)A_0(m_2^2).`$ (A.5) In all one-loop expressions of radiative corrections, we adopt a (modified) minimal subtraction procedure to remove poles in $`ϵ`$ and keep only finite (real) parts of the above functions. Some useful expressions for these functions in limiting cases are (after minimal subtraction) $`B_0(0,m_1^2,m_2^2)=1\mathrm{ln}{\displaystyle \frac{m_1^2}{Q^2}}+{\displaystyle \frac{m_2^2}{m_1^2m_2^2}}\mathrm{ln}{\displaystyle \frac{m_2^2}{m_1^2}},`$ (A.6) $`B_0(m_1^2,m_2^2,0)=2\mathrm{ln}{\displaystyle \frac{m_1^2}{Q^2}}\left(1{\displaystyle \frac{m_2^2}{m_1^2}}\right)\mathrm{ln}\left(1{\displaystyle \frac{m_2^2}{m_1^2}}\right){\displaystyle \frac{m_2^2}{m_1^2}}\mathrm{ln}{\displaystyle \frac{m_2^2}{m_1^2}},`$ (A.7) $`{\displaystyle \frac{d}{dp^2}}B_0(p^2,m^2,m^2)\|_{p^2=0}={\displaystyle \frac{1}{6m^2}},`$ (A.8) $`B_0(m^2,m^2,m^2)=\mathrm{ln}{\displaystyle \frac{m^2}{Q^2}}+2{\displaystyle \frac{\pi }{\sqrt{3}}}.`$ (A.9) ### A.2 Two-loop scalar functions In this subsection we collect some useful formulae of zero-point two-loop scalar functions. They have been studied extensively by several groups using two different methods: a differential equation method and an integral Mellin-Barnes transformation method ; their results all agree. Here we mainly follow Ref. . The momentum integrals appearing in a two-loop effective potential calculation can be reduced to the following two types of scalar functions \[corresponding to the topologies of two distinct zero-point two-loop irreducible Feynman diagrams (the figure-8 and sunset diagrams)\]: $$J(m_1^2,m_2^2)=(16\pi ^2\mu ^{4d})^2\frac{d^dpd^dq}{(2\pi )^{2d}}\frac{1}{[p^2m_1^2+i\epsilon ][q^2m_2^2+i\epsilon ]},$$ (A.10) and $$I(m_1^2,m_2^2,m_3^2)=(16\pi ^2\mu ^{4d})^2\frac{d^dpd^dq}{(2\pi )^{2d}}\frac{1}{[p^2m_1^2+i\epsilon ][q^2m_2^2+i\epsilon ][(p+q)^2m_3^2+i\epsilon ]}.$$ (A.11) The function $`J`$ is symmetric in $`m_1,m_2`$ and $`I`$ symmetric in $`m_1,m_2`$ and $`m_3`$. The function $`J`$ can be reduced to the product of one-loop scalar functions as $$J(m_1^2,m_2^2)=A_0(m_1^2)A_0(m_2^2).$$ (A.12) The function $`I`$ satisfies the following first-order partial differential equation $`R^2{\displaystyle \frac{}{m_3^2}}I(m_1^2,m_2^2,m_3^2)=(d3)(m_3^2m_1^2m_2^2)I(m_1^2,m_2^2,m_3^2)`$ (A.13) $`+`$ $`(d2)\left[{\displaystyle \frac{m_3^2m_1^2+m_2^2}{2m_3^2}}J(m_1^2,m_3^2)+{\displaystyle \frac{m_3^2+m_1^2m_2^2}{2m_3^2}}J(m_2^2,m_3^2)J(m_1^2,m_2^2)\right],`$ where $$R^2=m_1^4+m_2^4+m_3^42m_1^2m_2^22m_1^2m_3^22m_2^2m_3^2.$$ (A.14) This differential equation can be used to solve for the $`I`$ function. The initial value of this function can be evaluated from (A.13) which reduces to a simple algebraic equation when $`m_3=m_1+m_2`$, i.e. $`R=0`$. In our calculation, any Feynman diagram in the two-loop effective potential is subtracted by all its possible one-loop sub-diagrams; this is done by replacing the $`I`$ and $`J`$ functions as follows : $`I(m_1^2,m_2^2,m_3^2)`$ $``$ $`\widehat{I}(m_1^2,m_2^2,m_3^2)=I(m_1^2,m_2^2,m_3^2){\displaystyle \frac{1}{ϵ}}\left[A_0(m_1^2)+A_0(m_2^2)+A_0(m_3^2)\right],`$ $`J(m_1^2,m_2^2)`$ $``$ $`\widehat{J}(m_1^2,m_2^2)=J(m_1^2,m_2^2)+{\displaystyle \frac{1}{ϵ}}\left[m_1^2A_0(m_2^2)+m_2^2A_0(m_1^2)\right].`$ (A.15) It is then straightforward to show $$\widehat{J}(m_1^2,m_2^2)=\frac{m_1^2m_2^2}{ϵ^2}+m_1^2m_2^2\left(1\mathrm{ln}\frac{m_1^2}{Q^2}\right)\left(1\mathrm{ln}\frac{m_2^2}{Q^2}\right),$$ (A.16) and with some effort $`\widehat{I}(m_1^2,m_2^2,m_3^2)`$ $`=`$ $`{\displaystyle \frac{1}{2ϵ^2}}(m_1^2+m_2^2+m_3^2){\displaystyle \frac{1}{2ϵ}}(m_1^2+m_2^2+m_3^2)`$ (A.17) $``$ $`{\displaystyle \frac{1}{2}}[(m_1^2+m_2^2+m_3^2)\mathrm{ln}{\displaystyle \frac{m_2^2}{Q^2}}\mathrm{ln}{\displaystyle \frac{m_3^2}{Q^2}}+(m_1^2m_2^2+m_3^2)\mathrm{ln}{\displaystyle \frac{m_1^2}{Q^2}}\mathrm{ln}{\displaystyle \frac{m_3^2}{Q^2}}`$ $`+`$ $`(m_1^2+m_2^2m_3^2)\mathrm{ln}{\displaystyle \frac{m_1^2}{Q^2}}\mathrm{ln}{\displaystyle \frac{m_2^2}{Q^2}}4\left(m_1^2\mathrm{ln}{\displaystyle \frac{m_1^2}{Q^2}}+m_2^2\mathrm{ln}{\displaystyle \frac{m_2^2}{Q^2}}+m_3^2\mathrm{ln}{\displaystyle \frac{m_3^2}{Q^2}}\right)`$ $`+`$ $`\xi (m_1^2,m_2^2,m_3^2)+5(m_1^2+m_2^2+m_3^2)],`$ where (for $`R^2>0`$) $`\xi `$ is given by $`\xi (m_1^2,m_2^2,m_3^2)`$ $`=`$ $`R[2\mathrm{ln}\left({\displaystyle \frac{m_3^2+m_1^2m_2^2R}{2m_3^2}}\right)\mathrm{ln}\left({\displaystyle \frac{m_3^2m_1^2+m_2^2R}{2m_3^2}}\right)\mathrm{ln}{\displaystyle \frac{m_1^2}{m_3^2}}\mathrm{ln}{\displaystyle \frac{m_2^2}{m_3^2}}`$ (A.18) $``$ $`2Li_2\left({\displaystyle \frac{m_3^2+m_1^2m_2^2R}{2m_3^2}}\right)2Li_2\left({\displaystyle \frac{m_3^2m_1^2+m_2^2R}{2m_3^2}}\right)+{\displaystyle \frac{\pi ^2}{3}}],`$ where $`Li_2(x)`$ is the dilogarithm function $$Li_2(x)=_0^1𝑑y\frac{\mathrm{ln}(1xy)}{y}.$$ (A.19) In the region where $`R^2<0`$, (A.18) should be replaced by its analytical continuation. Equivalent expressions for $`\xi `$ also appear in and ; we find that (A.18) is most convenient for series expansions. We also define a function $`L`$ for future use $$L(m_1^2,m_2^2,m_3^2)=J(m_2^2,m_3^2)J(m_1^2,m_2^2)J(m_1^2,m_3^2)(m_1^2m_2^2m_3^2)I(m_1^2,m_2^2,m_3^2).$$ (A.20) Performing a (modified) minimal subtraction (by removing the single and double poles in $`ϵ`$), it is the finite (real) parts of (A.16) and (A.17) that we use in our two-loop effective potential expressions. We will also omit the carets of $`\widehat{I}`$ and $`\widehat{J}`$ to simplify the notation. When computing the two-loop potential, some argument of the $`I`$ function, e.g. the bottom-quark mass $`m_b`$, tree-level Higgs boson mass $`m_{h^0}`$, can be taken to be zero. The function $`I`$ is well-behaved in these limiting cases: $`I(m_1^2,m_2^2,0)`$ $`=`$ $`m_1^2\mathrm{ln}{\displaystyle \frac{m_1^2}{Q^2}}\mathrm{ln}{\displaystyle \frac{m_2^2}{Q^2}}(m_1^2m_2^2)\mathrm{ln}{\displaystyle \frac{m_1^2m_2^2}{Q^2}}\mathrm{ln}{\displaystyle \frac{m_1^2}{m_2^2}}+{\displaystyle \frac{1}{2}}(m_1^2m_2^2)\mathrm{ln}^2{\displaystyle \frac{m_1^2}{Q^2}}`$ $`+`$ $`2\left(m_1^2\mathrm{ln}{\displaystyle \frac{m_1^2}{Q^2}}+m_2^2\mathrm{ln}{\displaystyle \frac{m_2^2}{Q^2}}\right){\displaystyle \frac{5}{2}}(m_1^2+m_2^2)+(m_1^2m_2^2)\left[{\displaystyle \frac{\pi ^2}{6}}+Li_2\left({\displaystyle \frac{m_2^2}{m_1^2}}\right)\right],`$ $`I(m^2,0,0)`$ $`=`$ $`m^2\left({\displaystyle \frac{1}{2}}\mathrm{ln}^2{\displaystyle \frac{m^2}{Q^2}}2\mathrm{ln}{\displaystyle \frac{m^2}{Q^2}}+{\displaystyle \frac{5}{2}}+{\displaystyle \frac{\pi ^2}{6}}\right),`$ (A.22) where we have kept only the finite terms as explained before. In (A.2 Two-loop scalar functions) we have implicitly assumed $`m_1m_2`$. The symmetry of the above expression for $`I(m_1^2,m_2^2,0)`$ in $`m_1`$ and $`m_2`$ \[which obviously follows from the definition (A.15) of $`I`$\] can be explicitly checked by using the identity $$Li_2(x)=Li_2(x^1)\frac{1}{2}\mathrm{ln}^2(x)\frac{\pi ^2}{6}.$$ (A.23) Finally, we collect expansion formulae for the function $`\xi `$ which we use in the derivation of an analytical approximation formula for the two-loop Higgs boson mass corrections. The $`\xi `$ functions we find can be reduced to one of the different types we list below using the relation $$\xi (m_1^2,m_2^2,m_3^2)=m_1^2\xi (1,m_2^2/m_1^2,m_3^2/m_1^2).$$ (A.24) (1) For $`0r1`$ and $`0ϵ1`$: $`\xi (1,r,ϵ)`$ $`=`$ $`(1r)\left\{{\displaystyle \frac{\pi ^2}{3}}+\left[\mathrm{ln}ϵ2\mathrm{ln}(1r)\right]\mathrm{ln}r2Li_2(r)\right\}`$ (A.25) $``$ $`ϵ\left\{22\mathrm{ln}ϵ+\mathrm{ln}r+{\displaystyle \frac{1+r}{1r}}\left[{\displaystyle \frac{\pi ^2}{3}}+\left(\mathrm{ln}ϵ2\mathrm{ln}(1r)1\right)\mathrm{ln}r2Li_2(r)\right]\right\}`$ $`+`$ $`{\displaystyle \frac{ϵ^2}{(1r)^3}}\{({\displaystyle \frac{3}{2}}\mathrm{ln}ϵ)(1r^2){\displaystyle \frac{2\pi ^2}{3}}r[2\mathrm{ln}ϵ+r4\mathrm{ln}(1r)]r\mathrm{ln}r`$ $`+4rLi_2(r){\displaystyle \frac{}{}}\}+𝒪(ϵ^3).`$ If $`r>1`$, one uses $`\xi (1,r,ϵ)=r\xi (1,1/r,ϵ/r)`$ and the above expression. Two particular cases of the previous expansion are: (1a) For $`0ϵ_1,ϵ_21`$: $`\xi (1,ϵ_1,ϵ_2)`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{3}}+\mathrm{ln}ϵ_1\mathrm{ln}ϵ_22\left(1+{\displaystyle \frac{\pi ^2}{3}}+\mathrm{ln}ϵ_1\mathrm{ln}ϵ_2\right)ϵ_1ϵ_2`$ $`+`$ $`[(2{\displaystyle \frac{\pi ^2}{3}}+2\mathrm{ln}ϵ_1\mathrm{ln}ϵ_1\mathrm{ln}ϵ_2)ϵ_1+({\displaystyle \frac{3}{2}}\mathrm{ln}ϵ_1)ϵ_1^2+(ϵ_1ϵ_2)]+𝒪(ϵ_1^mϵ_2^n),`$ with $`m+n=3`$, and (1b) For $`0\|ϵ_1\|,ϵ_21`$: $`\xi (1,1+ϵ_1,ϵ_2)`$ $`=`$ $`2(4+ϵ_12\mathrm{ln}ϵ_2)ϵ_2+\left({\displaystyle \frac{8}{9}}{\displaystyle \frac{1}{3}}\mathrm{ln}ϵ_2\right)ϵ_2^2`$ (A.27) $`+`$ $`\left[2\mathrm{ln}ϵ_2+\left({\displaystyle \frac{7}{18}}+{\displaystyle \frac{1}{6}}\mathrm{ln}ϵ_2\right)ϵ_2\right]ϵ_1^2`$ $`+`$ $`\left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\mathrm{ln}ϵ_2\right)ϵ_1^3+\left({\displaystyle \frac{2}{9}}{\displaystyle \frac{1}{3}}\mathrm{ln}ϵ_2\right)ϵ_1^4+𝒪(ϵ_1^mϵ_2^n),`$ with $`m+2n5`$. Finally we also give (2) For $`\|ϵ_1\|,\|ϵ_2\|1`$: $`\xi (1,1+ϵ_1,1+ϵ_2)`$ $`=`$ $`36K+(8K1)ϵ_1ϵ_2+\left({\displaystyle \frac{5}{36}}{\displaystyle \frac{8}{3}}K\right)ϵ_1^2ϵ_2^2`$ (A.28) $`+`$ $`\{12Kϵ_1+(18K)ϵ_1^2+({\displaystyle \frac{8}{3}}K{\displaystyle \frac{2}{9}})ϵ_1^3+({\displaystyle \frac{1}{108}}{\displaystyle \frac{16}{9}}K)ϵ_1^4`$ $`+`$ $`[{\displaystyle \frac{ϵ_1^2}{6}}+({\displaystyle \frac{11}{54}}+{\displaystyle \frac{8}{9}}K)ϵ_1^3]ϵ_2+(ϵ_1ϵ_2)\}+𝒪(ϵ_1^mϵ_2^n),`$ with $`m+n=5`$. In this expansion the constant number $`K`$ is given by $$K=\frac{1}{\sqrt{3}}_0^{\pi /6}𝑑x\mathrm{ln}(2\mathrm{cos}x)0.1953256.$$ (A.29) ## Appendix B: MSSM in the leading approximation The general structure of the MSSM is quite complicated, with many different fields and field mixings. This makes the computation of the complete potential prohibitive at two-loops. However, it is a good aproximation to keep only those terms of the MSSM Lagrangian which depend on the $`SU(3)`$ gauge coupling $`g_3`$ and the top Yukawa $`h_t`$ (and neglect the electroweak gauge couplings $`g_1,g_2`$ and the rest of the Yukawa couplings). We call this the leading approximation and it greatly simplifies our two-loop effective potential calculation. In this Appendix, we summarize the necessary Feynman rules for computing the two-loop potential in this leading approximation and also some MSSM renormalization group equations, useful to check the scale invariance of the potential. ### B.1 Masses and Feynman rules The Higgs sector scalar potential in the leading approximation is $$V_{\mathrm{Higgs}}=(m_{H_1}^2+\mu ^2)\|H_1\|^2+(m_{H_2}^2+\mu ^2)\|H_2\|^2+B_\mu (H_1H_2+\mathrm{H}.\mathrm{c}.),$$ (B.1) where $`m_{H_1},m_{H_2}`$ and $`B_\mu `$ \[with dimensions of (mass)<sup>2</sup>\] are the soft-breaking Higgs mass parameters, and $`\mu `$ the supersymmetric Higgs-boson mass parameter. Although we do not write the quartic Higgs couplings, which depend on the electroweak gauge coupling constants, they are responsible for the tree-level mass of the lightest Higgs boson, which we of course include in our calculations. The $`SU(2)`$ doublet Higgs fields $`H_1`$ and $`H_2`$ can be written as follows: $$H_1=\left[\begin{array}{c}(h_1+ia_1)/\sqrt{2}\\ h_1^{}\end{array}\right],H_2=\left[\begin{array}{c}h_2^+\\ (h_2+ia_2)/\sqrt{2}\end{array}\right].$$ (B.2) In our approximation, the mass-squared matrices for $`𝒞𝒫`$-even and odd Higgs fields are $$_\pm ^2=\left(\begin{array}{cc}m_{H_1}^2+\mu ^2& \pm B_\mu \\ \pm B_\mu & m_{H_2}^2+\mu ^2\end{array}\right),$$ (B.3) where the positive (negative) sign applies to the $`𝒞𝒫`$-even (odd) fields respectively. The charged Higgs fields have the same mass-squared matrix $`_{}^2`$ as the $`𝒞𝒫`$-odd Higgses. The $`𝒞𝒫`$-even interaction eigenstates $`h_1,h_2`$ are rotated by the angle $`\alpha `$ into the mass eigenstates $`H^0`$ and $`h^0`$. Similarly, the $`𝒞𝒫`$-odd states $`a_1,a_2`$ (charged states $`h_1^+,h_2^+`$) are rotated into mass eigenstates $`G^0`$ and $`A^0`$ ($`G^+`$ and $`H^+`$) by the angle $`\beta `$. This angle $`\beta `$ is conventionally defined in terms of the $`𝒞𝒫`$-even Higgs field VEVs, $`h_{1,2}=v_{1,2}`$, by $`\mathrm{tan}\beta =v_2/v_1`$. The fact that $`\beta `$ diagonalizes $`_{}^2`$ is obvious when the minimization conditions of the potential (B.1), $`m_{H_1}^2+\mu ^2=B_\mu \mathrm{tan}\beta `$ and $`m_{H_2}^2+\mu ^2=B_\mu \mathrm{cot}\beta `$, are imposed and the soft parameters in the matrix are replaced by $`\mathrm{tan}\beta `$ and $`m_{A^0}^2=B_\mu (\mathrm{tan}\beta +\mathrm{cot}\beta )`$. Since we have neglected all $`g_1,g_2`$ related terms in (B.1), in our approximation (we use shorthand notations $`c_\beta =\mathrm{cos}\beta `$, $`s_\beta =\mathrm{sin}\beta `$, etc.) $$c_\alpha =s_\beta ,\mathrm{and}s_\alpha =c_\beta .$$ (B.4) This approximation is excellent when $`M_{A^0}M_Z`$ but would fail for $`M_{A^0}M_Z`$. The effect is numerically relevant for the tree level masses and we take it into account, but it may be consistently neglected in the two-loop corrections. The (field-dependent) top and bottom squark mass-squared matrices (neglecting the $`D`$-terms) are<sup>6</sup><sup>6</sup>6In this revised version, we have also included bottom Yukawa terms in the Feynman rules, they will be used in the expanded two-loop effective potential expression (D.6). $$_{\stackrel{~}{t}}^2=\left[\begin{array}{cc}M_{\stackrel{~}{Q}}^2+\frac{1}{2}h_t^2h_2^2& \frac{1}{\sqrt{2}}h_t(A_th_2+\mu h_1)\\ \frac{1}{\sqrt{2}}h_t(A_th_2+\mu h_1)& M_{\stackrel{~}{U}}^2+\frac{1}{2}h_t^2h_2^2\end{array}\right],$$ (B.5) $$_{\stackrel{~}{b}}^2=\left[\begin{array}{cc}M_{\stackrel{~}{Q}}^2+\frac{1}{2}h_b^2h_1^2& \frac{1}{\sqrt{2}}h_b(A_bh_1+\mu h_2)\\ \frac{1}{\sqrt{2}}h_b(A_bh_1+\mu h_2)& M_{\stackrel{~}{D}}^2+\frac{1}{2}h_b^2h_1^2\end{array}\right],$$ (B.6) where $`M_{\stackrel{~}{Q}}`$, $`M_{\stackrel{~}{U}}`$ ($`M_{\stackrel{~}{D}}`$) are soft-breaking mass parameters of the left- and right-handed top(bottom)-squarks $`\stackrel{~}{Q}`$ and $`\stackrel{~}{U}`$ ($`\stackrel{~}{D}`$); $`A_t`$ and $`A_b`$ are the usual trilinear soft-breaking parameters. We denote the mass eigenvalues of the matrix (B.5) by $`m_{\stackrel{~}{t}_1}`$, $`m_{\stackrel{~}{t}_2}`$ and the mixing angle by $`\theta _{\stackrel{~}{t}}`$, and the corresponding quantities for the matrix (B.6) by $`m_{\stackrel{~}{b}_1}`$, $`m_{\stackrel{~}{b}_2}`$ and $`\theta _{\stackrel{~}{b}}`$. The Feynman rules for Higgs/Goldstone-boson-squark trilinear coupling are simply $`i\lambda `$, with $`\lambda `$ as listed below: $`\lambda _{H^+\stackrel{~}{t}_1\stackrel{~}{b}_1}`$ $`=`$ $`h_tc_\beta [(c_tm_t+s_tY_t)c_b+m_bs_bs_t]h_bs_\beta [(c_bm_b+s_bY_b)c_t+m_ts_ts_b],`$ $`\lambda _{H^+\stackrel{~}{t}_1\stackrel{~}{b}_2}`$ $`=`$ $`h_tc_\beta [(c_tm_t+s_tY_t)s_bm_bc_bs_t]+h_bs_\beta [(s_bm_bc_bY_b)c_tm_ts_tc_b],`$ $`\lambda _{H^+\stackrel{~}{t}_2\stackrel{~}{b}_1}`$ $`=`$ $`h_tc_\beta [(s_tm_tc_tY_t)c_bm_bs_bc_t]+h_bs_\beta [(c_bm_b+s_bY_b)s_tm_tc_ts_b],`$ $`\lambda _{H^+\stackrel{~}{t}_2\stackrel{~}{b}_2}`$ $`=`$ $`h_tc_\beta [(s_tm_tc_tY_t)s_b+m_bc_bc_t]h_bs_\beta [(s_bm_bc_bY_b)s_t+m_tc_tc_b],`$ (B.7) $`\lambda _{G^+\stackrel{~}{t}_1\stackrel{~}{b}_1}`$ $`=`$ $`h_ts_\beta (c_tm_t+s_tX_t)c_b+h_bc_\beta (c_bm_b+s_bX_b)c_t,`$ $`\lambda _{G^+\stackrel{~}{t}_1\stackrel{~}{b}_2}`$ $`=`$ $`h_ts_\beta (c_tm_t+s_tX_t)s_b+h_bc_\beta (s_bm_b+c_bX_b)c_t,`$ $`\lambda _{G^+\stackrel{~}{t}_2\stackrel{~}{b}_1}`$ $`=`$ $`h_ts_\beta (s_tm_tc_tX_t)c_bh_bc_\beta (c_bm_b+s_bX_b)s_t,`$ $`\lambda _{G^+\stackrel{~}{t}_2\stackrel{~}{b}_2}`$ $`=`$ $`h_ts_\beta (s_tm_t+c_tX_t)s_b+h_bc_\beta (s_bm_bc_bX_b)s_t,`$ (B.8) and $`\lambda _{H^0\stackrel{~}{t}_1\stackrel{~}{t}_1}=\sqrt{2}h_t(m_t+s_tc_tY_t^\alpha )s_\alpha ,`$ $`\lambda _{H^0\stackrel{~}{t}_2\stackrel{~}{t}_2}=\sqrt{2}h_t(m_ts_tc_tY_t^\alpha )s_\alpha ,`$ $`\lambda _{h^0\stackrel{~}{t}_1\stackrel{~}{t}_1}=\sqrt{2}h_t(m_t+s_tc_tX_t^\alpha )c_\alpha ,`$ $`\lambda _{h^0\stackrel{~}{t}_2\stackrel{~}{t}_2}=\sqrt{2}h_t(m_ts_tc_tX_t^\alpha )c_\alpha ,`$ $`\lambda _{H^0\stackrel{~}{t}_1\stackrel{~}{t}_2}={\displaystyle \frac{1}{\sqrt{2}}}h_tc_{2t}Y_t^\alpha s_\alpha ,`$ $`\lambda _{h^0\stackrel{~}{t}_1\stackrel{~}{t}_2}={\displaystyle \frac{1}{\sqrt{2}}}h_tc_{2t}X_t^\alpha c_\alpha ,`$ $`\lambda _{A^0\stackrel{~}{t}_1\stackrel{~}{t}_2}=\lambda _{A^0\stackrel{~}{t}_2\stackrel{~}{t}_1}={\displaystyle \frac{1}{\sqrt{2}}}h_tY_tc_\beta ,`$ $`\lambda _{G^0\stackrel{~}{t}_1\stackrel{~}{t}_2}=\lambda _{G^0\stackrel{~}{t}_2\stackrel{~}{t}_1}={\displaystyle \frac{1}{\sqrt{2}}}h_tX_ts_\beta ,`$ (B.9) where $`c_t=\mathrm{cos}\theta _{\stackrel{~}{t}}`$, $`s_t=\mathrm{sin}\theta _{\stackrel{~}{t}}`$, $`c_{2t}=\mathrm{cos}2\theta _{\stackrel{~}{t}}`$ (with similar expressions for $`\theta _{\stackrel{~}{b}}`$ functions) and $$X_t=A_t+\mu \mathrm{cot}\beta ,Y_t=A_t\mu \mathrm{tan}\beta ,$$ (B.10) $$X_b=A_b+\mu \mathrm{tan}\beta ,Y_b=A_b\mu \mathrm{cot}\beta .$$ (B.11) In addition, we find convenient to define the $`\alpha `$-dependent quantities $$X_t^\alpha =A_t\mu \mathrm{tan}\alpha ,Y_t^\alpha =A_t+\mu \mathrm{cot}\alpha ,$$ (B.12) $$X_b^\alpha =A_b\mu \mathrm{cot}\alpha ,Y_b^\alpha =A_b+\mu \mathrm{tan}\alpha ,$$ (B.13) which tend to the corresponding quantities without the $`\alpha `$ label ($`X_t^\alpha X_t`$, etc) in the limit $`m_AM_Z`$. Couplings similar to the above ones but for bottom squarks can be obtained directly from (B.7), (B.8) and (B.9): simply make everywhere the replacements $`\{h_th_b,m_tm_b,\theta _{\stackrel{~}{t}}\theta _{\stackrel{~}{b}},X_t^{(\alpha )}X_b^{(\alpha )},Y_t^{(\alpha )}Y_b^{(\alpha )}\}`$ and $`\{c_\alpha s_\alpha ,c_\beta s_\beta \}`$ for the couplings to $`\{H^+,A^0,H^0\}`$ or $`\{c_\alpha s_\alpha ,c_\beta s_\beta \}`$ for the couplings to $`\{G^+,G^0,h^0\}`$. The couplings of squarks to neutralinos and charginos are very simple in the leading approximation, since the gaugino-Higgsino mixing can be neglected and the only interactions are Higgsino-squark interactions. The Feynman rules for the $`\stackrel{~}{h}_i^0t\stackrel{~}{t}_j`$ couplings can be written as $`i(a𝒫_L+b𝒫_R)`$ and that of $`\stackrel{~}{h}^+t\stackrel{~}{b}_L`$ as $`i𝒞^1(a𝒫_L+b𝒫_R)`$ ($`𝒫_{L,R}`$ are chiral projectors and $`𝒞`$ the charge-conjugation matrix), with $`a_{\stackrel{~}{h}_1^0t\stackrel{~}{t}_1}=ia_{\stackrel{~}{h}_2^0t\stackrel{~}{t}_1}=b_{\stackrel{~}{h}_1^0t\stackrel{~}{t}_2}=ib_{\stackrel{~}{h}_2^0t\stackrel{~}{t}_2}={\displaystyle \frac{h_t}{\sqrt{2}}}c_t,`$ $`a_{\stackrel{~}{h}_1^0t\stackrel{~}{t}_2}=ia_{\stackrel{~}{h}_2^0t\stackrel{~}{t}_2}=b_{\stackrel{~}{h}_1^0t\stackrel{~}{t}_1}=ib_{\stackrel{~}{h}_2^0t\stackrel{~}{t}_1}={\displaystyle \frac{h_t}{\sqrt{2}}}s_t,`$ $`a_{\stackrel{~}{h}^+t\stackrel{~}{b}_L}=h_t,a_{\stackrel{~}{h}^+b\stackrel{~}{t}_1}=h_ts_t,a_{\stackrel{~}{h}^+b\stackrel{~}{t}_2}=h_tc_t,`$ (B.14) when $`\mu >0`$; for $`\mu <0`$, we only need to interchange $`a_{\stackrel{~}{h}_1^0t\stackrel{~}{t}_i}`$ and $`a_{\stackrel{~}{h}_2^0t\stackrel{~}{t}_i}`$, as well as $`b_{\stackrel{~}{h}_1^0t\stackrel{~}{t}_i}`$ and $`b_{\stackrel{~}{h}_2^0t\stackrel{~}{t}_i}`$. Other Feynman rules of $`𝒪(g_3)`$ and $`𝒪(h_t)`$ vertices are exactly the same as in the general MSSM and we do not present them explicitly. ### B.2 Renormalization group equations The MSSM RGEs that we will use to check the invariance of the potential to two-loop order under renormalization scale transformations are the following. First, we need the two-loop RGEs for those parameters entering in the tree-level potential (B.1) $`{\displaystyle \frac{m_{H_2}^2}{\mathrm{ln}Q^2}}`$ $`=`$ $`{\displaystyle \frac{3h_t^2}{16\pi ^2}}_t^2+{\displaystyle \frac{16g_3^2h_t^2}{(16\pi ^2)^2}}(_t^2+2M_3^22M_3A_t){\displaystyle \frac{18h_t^4}{(16\pi ^2)^2}}(_t^2+A_t^2),`$ (B.15) $`{\displaystyle \frac{\mathrm{ln}\mu }{\mathrm{ln}Q^2}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{ln}h_2}{\mathrm{ln}Q^2}}={\displaystyle \frac{3h_t^2}{32\pi ^2}}+{\displaystyle \frac{8g_3^2h_t^2}{(16\pi ^2)^2}}{\displaystyle \frac{9}{2}}{\displaystyle \frac{h_t^4}{(16\pi ^2)^2}},`$ (B.16) $`{\displaystyle \frac{B_\mu }{\mathrm{ln}Q^2}}`$ $`=`$ $`{\displaystyle \frac{3h_t^2}{16\pi ^2}}\left({\displaystyle \frac{B_\mu }{2}}+A_t\mu \right)+{\displaystyle \frac{16g_3^2h_t^2}{(16\pi ^2)^2}}\left({\displaystyle \frac{B_\mu }{2}}+A_t\mu M_3\mu \right)`$ (B.17) $``$ $`{\displaystyle \frac{9h_t^4}{(16\pi ^2)^2}}\left({\displaystyle \frac{B_\mu }{2}}+2A_t\mu \right),`$ where $`_t^2=m_{H_2}^2+M_{\stackrel{~}{Q}}^2+M_{\stackrel{~}{U}}^2+A_t^2`$. Then we need one-loop RGEs for those masses entering in the one-loop potential $`16\pi ^2{\displaystyle \frac{m_t^2}{\mathrm{ln}Q^2}}`$ $`=`$ $`\left({\displaystyle \frac{16g_3^2}{3}}+3h_t^2\right)m_t^2,`$ (B.18) $`16\pi ^2{\displaystyle \frac{m_{\stackrel{~}{t}_1}^2}{\mathrm{ln}Q^2}}`$ $`=`$ $`{\displaystyle \frac{16g_3^2}{3}}\left[m_t^2+M_3^2s_{2t}m_t\left(M_3{\displaystyle \frac{X_t}{2}}\right)\right]`$ (B.19) $`+`$ $`h_t^2\left[3m_t^2+(1+s_t^2)_t^2+3s_{2t}m_t\left(A_t+{\displaystyle \frac{3X_t}{2}}\right)\right],`$ $`16\pi ^2{\displaystyle \frac{m_{\stackrel{~}{t}_2}^2}{\mathrm{ln}Q^2}}`$ $`=`$ $`{\displaystyle \frac{16g_3^2}{3}}\left[m_t^2+M_3^2+s_{2t}m_t\left(M_3{\displaystyle \frac{X_t}{2}}\right)\right]`$ (B.20) $`+`$ $`h_t^2\left[3m_t^2+(1+c_t^2)_t^23s_{2t}m_t\left(A_t+{\displaystyle \frac{3X_t}{2}}\right)\right],`$ $`16\pi ^2{\displaystyle \frac{m_{H_n^0}^2}{\mathrm{ln}Q^2}}`$ $`=`$ $`3h_t^2\left[\mu ^2+D_n_t^2+E_n\left({\displaystyle \frac{B_\mu }{2}}+A_t\mu \right)\right],`$ (B.21) $`16\pi ^2{\displaystyle \frac{m_{H_n^+}^2}{\mathrm{ln}Q^2}}`$ $`=`$ $`3h_t^2\left[\mu ^2+D_{n+2}_t^2+E_{n+2}\left({\displaystyle \frac{B_\mu }{2}}+A_t\mu \right)\right],`$ (B.22) where $`D_n=s_\alpha ^2,c_\alpha ^2,s_\beta ^2,c_\beta ^2`$ and $`E_n=s_{2\alpha },s_{2\alpha },s_{2\beta },s_{2\beta }`$ for $`n=1,2,3,4`$. \[Here we use the $`\alpha `$ angle to keep track the $`H^0`$ and $`h^0`$ contributions; it can be replaced by the $`\beta `$ angle as in (B.4) in the leading approximation.\] The ordering of the Higgs/Goldstone bosons are $`H_n^0=H^0,h^0,G^0`$ and $`A^0`$ for $`n=1,2,3,4`$ and $`H_n^+=G^+,H^+`$ for $`n=1,2`$. Eqs. (B.19-B.22) seem unfamiliar, but they follow directly from (B.3), (B.5) and the one-loop MSSM RGEs of soft parameters entering those equations. Using (B.19) and (B.20), we find one-loop RGEs for $`X_t`$ and $`m_{\stackrel{~}{t}}^2`$, the (arithmetic) average of the (squared) top squark masses. They are $`16\pi ^2{\displaystyle \frac{X_t}{\mathrm{ln}Q^2}}`$ $`=`$ $`{\displaystyle \frac{16}{3}}g_3^2M_3+3h_t^2(A_t+X_t).`$ (B.23) $`16\pi ^2{\displaystyle \frac{m_{\stackrel{~}{t}}^2}{\mathrm{ln}Q^2}}`$ $`=`$ $`{\displaystyle \frac{16}{3}}g_3^2(m_t^2+M_3^2)+h_t^2\left(3m_t^2+{\displaystyle \frac{3}{2}}_t^2\right),`$ (B.24) these two equations are used in Sec. 3 for the RG discussion of the formula for the Higgs boson mass $`M_{h^0}`$. Eq. (B.23) can also be derived from (B.10) and one-loop RGEs of $`A_t`$, $`\mu `$ and $`\mathrm{tan}\beta `$. ## Appendix C: One-loop self-energies In this appendix, we collect formulae for those MSSM one-loop self-energies which are necessary for our analysis. We present these self-energies in the leading approximation of keeping only $`h_t`$ and $`g_3`$-dependent terms, as explained in Appendix B; their full form can be found in Ref. , which we follow for notation. (See also for top quark/squark self-energies.) –Top quark: $`16\pi ^2\mathrm{\Sigma }_t(p^2)`$ $`=`$ $`{\displaystyle \frac{4g_3^2}{3}}\{m_t[B_1(p^2,m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{t}_1}^2)+B_1(p^2,m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{t}_2}^2)]m_t(53\mathrm{ln}{\displaystyle \frac{m_t^2}{Q^2}})`$ (C.1) $``$ $`s_{2t}m_{\stackrel{~}{g}}[B_0(p^2,m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{t}_1}^2)B_0(p^2,m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{t}_2}^2)]\}`$ $`+`$ $`{\displaystyle \frac{h_t^2}{2}}m_t\{c_\beta ^2[2B_1(p^2,m_t^2,m_{A^0}^2)+B_1(p^2,m_b^2,m_{A^0}^2)]`$ $`+`$ $`s_\beta ^2\left[2B_1(p^2,m_t^2,m_Z^2)+B_1(p^2,m_b^2,m_Z^2)\right]`$ $`+`$ $`B_1(p^2,\mu ^2,m_{\stackrel{~}{t}_1}^2)+B_1(p^2,\mu ^2,m_{\stackrel{~}{t}_2}^2)+B_1(p^2,\mu ^2,m_{\stackrel{~}{b}_L}^2)\},`$ where we have assumed all heavy Higgs bosons have mass $`m_{A^0}`$ much larger than the masses of the light Higgs and $`W`$-boson, taken to be $`m_Z`$. From (C.1) we find the running top-quark mass at the scale $`Q`$ (under the simplified assumptions of a common heavy SUSY scale $`M_S`$ while the $`\mu `$ parameter is left free, see Sec. 2) $`m_t^2(Q)`$ $`=`$ $`M_t^2\{1{\displaystyle \frac{g_3^2}{6\pi ^2}}[53\mathrm{ln}{\displaystyle \frac{m_t^2}{Q^2}}+\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}\widehat{X}_t]`$ $`+`$ $`{\displaystyle \frac{3h_t^2}{32\pi ^2}}[(1+c_\beta ^2)({\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}})+s_\beta ^2({\displaystyle \frac{8}{3}}\mathrm{ln}{\displaystyle \frac{m_t^2}{Q^2}}){\displaystyle \frac{\widehat{\mu }^2}{1\widehat{\mu }^2}}(1+{\displaystyle \frac{\widehat{\mu }^2}{1\widehat{\mu }^2}}\mathrm{ln}\widehat{\mu }^2\left)\right]\}.`$ In this equation we have neglected the external momentum and used (A.3) and (A.6). We have used the reduced parameters $`\widehat{X}_tX_t/M_S`$, $`\widehat{\mu }\mu /M_S`$ and $`M_t`$ is the top quark pole mass (we use capital letters to denote on-shell mass parameters). –Top squarks: $`16\pi ^2\mathrm{\Pi }_{\stackrel{~}{t}_L\stackrel{~}{t}_L}(p^2)`$ $`=`$ $`{\displaystyle \frac{8g_3^2}{3}}\{G(p^2,m_{\stackrel{~}{g}}^2,m_t^2)+c_t^2[A_0(m_{\stackrel{~}{t}_1}^2)(p^2+m_{\stackrel{~}{t}_1}^2)B_0(p^2,m_{\stackrel{~}{t}_1}^2,0)]`$ (C.3) $`+`$ $`s_t^2[A_0(m_{\stackrel{~}{t}_2}^2)(p^2+m_{\stackrel{~}{t}_2}^2)B_0(p^2,m_{\stackrel{~}{t}_2}^2,0)]\}`$ $`+`$ $`h_t^2\left[s_t^2A_0(m_{\stackrel{~}{t}_1}^2)+c_t^2A_0(m_{\stackrel{~}{t}_2}^2)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{4}{}}}D_nA_0(m_{H_n^0}^2)+G(p^2,\mu ^2,m_t^2)\right]`$ $`+`$ $`{\displaystyle \underset{n=1}{\overset{4}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\lambda _{H_n^0\stackrel{~}{t}_L\stackrel{~}{t}_i}^2B_0(p^2,m_{H_n^0}^2,m_{\stackrel{~}{t}_i}^2)+{\displaystyle \underset{n=1}{\overset{2}{}}}\lambda _{H_n^+\stackrel{~}{t}_L\stackrel{~}{b}_L}^2B_0(p^2,m_{H_n^+}^2,m_{\stackrel{~}{b}_L}^2),`$ $`16\pi ^2\mathrm{\Pi }_{\stackrel{~}{t}_R\stackrel{~}{t}_R}(p^2)`$ $`=`$ $`{\displaystyle \frac{8g_3^2}{3}}\{G(p^2,m_{\stackrel{~}{g}}^2,m_t^2)+s_t^2[A_0(m_{\stackrel{~}{t}_1}^2)(p^2+m_{\stackrel{~}{t}_1}^2)B_0(p^2,m_{\stackrel{~}{t}_1}^2,0)]`$ (C.4) $`+`$ $`c_t^2[A_0(m_{\stackrel{~}{t}_2}^2)(p^2+m_{\stackrel{~}{t}_2}^2)B_0(p^2,m_{\stackrel{~}{t}_2}^2,0)]\}`$ $`+`$ $`h_t^2[c_t^2A_0(m_{\stackrel{~}{t}_1}^2)+s_t^2A_0(m_{\stackrel{~}{t}_2}^2)+A_0(m_{\stackrel{~}{b}_L}^2)`$ $`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{4}{}}}D_nA_0(m_{H_n^0}^2)+{\displaystyle \underset{n=1}{\overset{2}{}}}D_{n+2}A_0(m_{H_n^+}^2)+G(p^2,\mu ^2,m_t^2)+G(p^2,\mu ^2,m_b^2)]`$ $`+`$ $`{\displaystyle \underset{n=1}{\overset{4}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\lambda _{H_n^0\stackrel{~}{t}_R\stackrel{~}{t}_i}^2B_0(p^2,m_{H_n^0}^2,m_{\stackrel{~}{t}_i}^2)+{\displaystyle \underset{n=1}{\overset{2}{}}}\lambda _{H_n^+\stackrel{~}{t}_R\stackrel{~}{b}_L}^2B_0(p^2,m_{H_n^+}^2,m_{\stackrel{~}{b}_L}^2),`$ $`16\pi ^2\mathrm{\Pi }_{\stackrel{~}{t}_L\stackrel{~}{t}_R}(p^2)`$ $`=`$ $`{\displaystyle \frac{4g_3^2}{3}}[s_{2t}(p^2+m_{\stackrel{~}{t}_1}^2)B_0(p^2,m_{\stackrel{~}{t}_1}^2,0)+s_{2t}(p^2+m_{\stackrel{~}{t}_2}^2)B_0(p^2,m_{\stackrel{~}{t}_2}^2,0)`$ (C.5) $`+`$ $`4m_{\stackrel{~}{g}}m_tB_0(p^2,m_{\stackrel{~}{g}}^2,m_t^2)]+{\displaystyle \frac{3}{2}}h^2_ts_{2t}[A_0(m_{\stackrel{~}{t}_1}^2)A_0(m_{\stackrel{~}{t}_2}^2)]`$ $`+`$ $`{\displaystyle \underset{n=1}{\overset{4}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\lambda _{H_n^0\stackrel{~}{t}_L\stackrel{~}{t}_i}\lambda _{H_n^0\stackrel{~}{t}_R\stackrel{~}{t}_i}B_0(p^2,m_{H_n^0}^2,m_{\stackrel{~}{t}_i}^2)`$ $`+`$ $`{\displaystyle \underset{n=1}{\overset{2}{}}}\lambda _{H_n^+\stackrel{~}{t}_L\stackrel{~}{b}_L}\lambda _{H_n^+\stackrel{~}{t}_R\stackrel{~}{b}_L}B_0(p^2,m_{H_n^+}^2,m_{\stackrel{~}{b}_L}^2),`$ where $`\lambda _{H^0\stackrel{~}{t}_L\stackrel{~}{t}_1}=c_t\lambda _{H^0\stackrel{~}{t}_1\stackrel{~}{t}_1}s_t\lambda _{H^0\stackrel{~}{t}_2\stackrel{~}{t}_1}`$, $`\lambda _{H^0\stackrel{~}{t}_R\stackrel{~}{t}_1}=s_t\lambda _{H^0\stackrel{~}{t}_1\stackrel{~}{t}_1}+c_t\lambda _{H^0\stackrel{~}{t}_2\stackrel{~}{t}_1}`$, etc.. The symbols $`D_n`$ are defined after (B.10). From (C.3-C.5) we derive relations between running and on-shell top-squark masses and mixing parameters using the following one-loop relationships (for $`c_t^2=s_t^2=1/2`$): $`M_{\stackrel{~}{t}_1}^2`$ $`=`$ $`M_{\stackrel{~}{Q}}^2+m_t^2+m_tX_t{\displaystyle \frac{1}{2}}\mathrm{Re}\left[\mathrm{\Pi }_{\stackrel{~}{t}_L\stackrel{~}{t}_L}(M_{\stackrel{~}{t}_1}^2)+\mathrm{\Pi }_{\stackrel{~}{t}_R\stackrel{~}{t}_R}(M_{\stackrel{~}{t}_1}^2)\right]\mathrm{Re}\mathrm{\Pi }_{\stackrel{~}{t}_L\stackrel{~}{t}_R}(M_{\stackrel{~}{t}_1}^2),`$ $`M_{\stackrel{~}{t}_2}^2`$ $`=`$ $`M_{\stackrel{~}{Q}}^2+m_t^2m_tX_t{\displaystyle \frac{1}{2}}\mathrm{Re}\left[\mathrm{\Pi }_{\stackrel{~}{t}_L\stackrel{~}{t}_L}(M_{\stackrel{~}{t}_2}^2)+\mathrm{\Pi }_{\stackrel{~}{t}_R\stackrel{~}{t}_R}(M_{\stackrel{~}{t}_2}^2)\right]+\mathrm{Re}\mathrm{\Pi }_{\stackrel{~}{t}_L\stackrel{~}{t}_R}(M_{\stackrel{~}{t}_2}^2),`$ (C.6) we obtain (assuming again a common heavy SUSY scale $`M_S`$ and leaving free the $`\mu `$-parameter): $`m_{\stackrel{~}{t}}^2(Q)`$ $`=`$ $`M_{\stackrel{~}{t}}^2\{1{\displaystyle \frac{g_3^2}{3\pi ^2}}(2\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}})+{\displaystyle \frac{3h_t^2}{32\pi ^2}}[(\widehat{X}_t^2s_\beta ^2+\widehat{Y}_t^2c_\beta ^2)(2\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}})`$ (C.7) $`+`$ $`c_\beta ^2\left(1{\displaystyle \frac{\pi }{\sqrt{3}}}\widehat{Y}_t^2\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}\right)`$ $`+`$ $`\widehat{\mu }^4\mathrm{ln}\widehat{\mu }^2+(1\widehat{\mu }^2)(32\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}})(1\widehat{\mu }^2)^2\mathrm{ln}(1\widehat{\mu }^2)]\},`$ $`m_tX_t(Q)`$ $`=`$ $`M_tX_t^{\mathrm{OS}}+{\displaystyle \frac{g_3^2}{12\pi ^2}}m_tM_S\left[4(2\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}})+2\widehat{X}_t\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}\right]`$ (C.8) $`+`$ $`{\displaystyle \frac{3h_t^2}{16\pi ^2}}m_t\{(X_ts_\beta ^2+Y_tc_\beta ^2)(2\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}}){\displaystyle \frac{\pi }{\sqrt{3}}}Y_tc_\beta ^2+X_t(1{\displaystyle \frac{3}{2}}\mathrm{ln}{\displaystyle \frac{M_S^2}{Q^2}})`$ $``$ $`{\displaystyle \frac{1}{2}}\left[1\widehat{\mu }^2+\widehat{\mu }^4\mathrm{ln}\widehat{\mu }^2+(1\widehat{\mu }^4)\mathrm{ln}(1\widehat{\mu }^2)\right]X_t`$ $`+`$ $`({\displaystyle \frac{1}{2}}+{\displaystyle \frac{\pi }{3\sqrt{3}}})c_\beta ^2\widehat{Y}_t^2X_t{\displaystyle \frac{1}{2}}s_\beta ^2\widehat{X}_t^2X_t\mathrm{ln}\left({\displaystyle \frac{m_tX_t}{M_S^2}}\right){\displaystyle \frac{1}{3}}s_\beta ^2\widehat{X}_t^2X_t\mathrm{ln}2\},`$ where we have used (A.7), (A.9) and the definition $`\widehat{Y}_t(A_t\mu \mathrm{tan}\beta )/M_S`$. $`W`$ boson: $`16\pi ^2\mathrm{\Pi }_{WW}^T(p^2)`$ $`=`$ $`3g^2\{2B_{22}(p^2,m_t^2,m_b^2)+{\displaystyle \frac{1}{2}}G(p^2,m_t^2,m_b^2)2c_t^2[B_{22}(p^2,m_{\stackrel{~}{t}_1}^2,m_{\stackrel{~}{b}_L}^2){\displaystyle \frac{1}{4}}A_0(m_{\stackrel{~}{t}_1}^2)]`$ (C.9) $``$ $`2s_t^2[B_{22}(p^2,m_{\stackrel{~}{t}_2}^2,m_{\stackrel{~}{b}_L}^2){\displaystyle \frac{1}{4}}A_0(m_{\stackrel{~}{t}_2}^2)]+{\displaystyle \frac{1}{2}}A_0(m_{\stackrel{~}{b}_L}^2)\}.`$ This gives (under the assumption of the simplified SUSY spectrum of Sec. 2, described already for previous self-energies) $$v^2(Q)=\frac{4}{g^2}[M_W^2+\mathrm{Re}\mathrm{\Pi }_{WW}^T(M_W^2)]=\frac{4M_W^2}{g^2}\left[1\frac{h_t^2s_\beta ^2}{32\pi ^2}\left(6\mathrm{ln}\frac{m_t^2}{Q^2}+3+\widehat{X}_t^2\right)\right],$$ (C.10) where we have neglected the external momentum in (C.9) and used (A.3-A.6). –Higgs boson: We need only the difference $`16\pi ^2[\mathrm{\Pi }_{hh}(m_{h^0}^2)\mathrm{\Pi }_{hh}(0)]`$ $`=`$ $`3h_t^2m_{h^0}^2s_\beta ^2\left[B_0(0,m_t^2,m_t^2)4m_t^2{\displaystyle \frac{d}{dp^2}}B_0(p^2,m_t^2,m_t^2)\|_{p^2=0}\right]`$ (C.11) $`+`$ $`3m_{h^0}^2{\displaystyle \underset{i,j}{}}\lambda _{h^0\stackrel{~}{t}_i\stackrel{~}{t}_j}^2{\displaystyle \frac{d}{dp^2}}B_0(p^2,m_{\stackrel{~}{t}_i}^2,m_{\stackrel{~}{t}_j}^2)\|_{p^2=0},`$ where $`\lambda _{h\stackrel{~}{t}_i\stackrel{~}{t}_j}`$ are defined in (B.9). Using (A.6) and (A.7) we get (19). ## Appendix D: MSSM effective potential to the two-loop order In this Appendix we present the MSSM effective potential for the (real) neutral components of the Higgs fields up to the two-loop level in the leading approximation (which neglects all dimensionless couplings except $`h_t`$ and $`g_3`$). We first write the potential as $$V(h_1,h_2)=V_{\mathrm{vac}}+V_0(h_1,h_2)+V_1(h_1,h_2)+V_2(h_1,h_2),$$ (D.1) where $`V_{\mathrm{vac}}`$ is a field-independent vacuum energy term<sup>7</sup><sup>7</sup>7This term is a function of the soft-breaking parameters; it is needed for the invariance of the potential under a RG transformation.. The tree-level potential $`V_0`$ is $$V_0(h_1,h_2)=\frac{1}{2}(m_{H_1}^2+\mu ^2)h_1^2+\frac{1}{2}(m_{H_2}^2+\mu ^2)h_2^2+B_\mu h_1h_2,$$ (D.2) which simply follows from substituting Eq. (B.2) into the MSSM Higgs sector scalar potential (B.1). The one-loop potential is well known and the $`𝒪(\alpha _s\alpha _t)`$ part of the two-loop potential was computed in ; we list them here for completeness and for future reference. The complete one-loop potential in Laudau gauge is<sup>8</sup><sup>8</sup>8We adopt the (modified) $`\overline{\mathrm{DR}}`$ -scheme of Ref. . $`16\pi ^2V_1(h_1,h_2)`$ $`=`$ $`{\displaystyle \underset{f}{}}N_c^f\left[{\displaystyle \underset{i=1,2}{}}H(m_{\stackrel{~}{f}_i}^2)2H(m_f^2)\right]+3H(m_W^2)+{\displaystyle \frac{3}{2}}H(m_Z^2)`$ (D.3) $`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{4}{}}}H(m_{H_n^0}^2)+{\displaystyle \underset{n=1}{\overset{2}{}}}H(m_{H_n^+}^2)2{\displaystyle \underset{i=1}{\overset{2}{}}}H(m_{\stackrel{~}{\chi }_i^+}^2){\displaystyle \underset{i=1}{\overset{4}{}}}H(m_{\stackrel{~}{\chi }_i^0}^2),`$ where $`f`$ sums over all the (s)quarks and (s)leptons, $`N_c^f`$ is the color factor, 3 for (s)quarks and 1 for (s)leptons. Following the leading approximation, we keep only the numerically important parts, i.e., those from top (s)quarks. In Eq. (D.3), $`\stackrel{~}{\chi }_i^+(i=1,2)`$ and $`\stackrel{~}{\chi }_i^0(i=1,2,3,4)`$ represent charginos and neutralinos, and the function $`H`$ is $$H(m^2)=\frac{m^4}{2}\left(\mathrm{ln}\frac{m^2}{Q^2}\frac{3}{2}\right).$$ (D.4) The QCD contribution to the two-loop effective potential in the MSSM is $`(16\pi ^2)^2V_{2s}(h_1,h_2)=8g_3^2\{J(m_t^2,m_t^2)2m_t^2I(m_t^2,m_t^2,0)`$ $`+{\displaystyle \frac{1}{2}}(c_t^4+s_t^4){\displaystyle \underset{i=1}{\overset{2}{}}}J(m_{\stackrel{~}{t}_i}^2,m_{\stackrel{~}{t}_i}^2)+2s_t^2c_t^2J(m_{\stackrel{~}{t}_1}^2,m_{\stackrel{~}{t}_2}^2)+{\displaystyle \underset{i=1}{\overset{2}{}}}m_{\stackrel{~}{t}_i}^2I(m_{\stackrel{~}{t}_i}^2,m_{\stackrel{~}{t}_i}^2,0)`$ $`+{\displaystyle \underset{i=1}{\overset{2}{}}}L(m_{\stackrel{~}{t}_i}^2,m_{\stackrel{~}{g}}^2,m_t^2)4m_{\stackrel{~}{g}}m_ts_tc_t\left[I(m_{\stackrel{~}{t}_1}^2,m_{\stackrel{~}{g}}^2,m_t^2)I(m_{\stackrel{~}{t}_2}^2,m_{\stackrel{~}{g}}^2,m_t^2)\right]`$ $`+[m_tm_b,m_{\stackrel{~}{t}_i}m_{\stackrel{~}{b}_i},\theta _{\stackrel{~}{t}}\theta _{\stackrel{~}{b}}]\},`$ (D.5) where $`\stackrel{~}{g}`$ is the gluino, with tree-level mass given by the $`SU(3)`$ gaugino soft mass $`M_3`$. The last term, obtained by interchanging variables, gives the contribution from sbottoms. Note that there is no mixed contribution involving stops and sbottoms, even if such mixed couplings exists \[from $`SU(3)`$ $`D`$-terms\]. The two-loop scalar functions $`I`$, $`J`$ and $`L`$ in Eq. (D.5) are given in Appendix A, Eqs. (A.17), (A.16) and (A.20).<sup>9</sup><sup>9</sup>9The procedure we have followed of subtracting all possible one-loop subdivergences to define these functions is an alternative to the direct way used in Ref. . The direct way requires the computation of some one-loop quantities to order $`𝒪(ϵ)`$; perhaps we find the subtraction method simpler. We have explicitly checked that, in the particular limit studied in , we exactly reproduce their unexpanded mass formula, Eq. (11) of , which shows the equivalence of both methods. The top Yukawa contribution to the two-loop potential is a new result of this paper. The relevant Feynman diagrams are shown in Fig. 8. To simplify the final result, we neglect left-right mixings in the bottom-squark sector and the gaugino-Higgsino mixings in the neutralino-chargino sector (under this assumption, the Higgsino masses are simply $`\|\mu \|`$); these simplifications are valid in the leading approximation. Using the Feynman rules given in Appendix B, we find (the last diagram of Fig. 8 is of order $`h_t^2`$ but does not contribute to $`m_{h^0}`$) the top and bottom Yukawa contribution to the two-loop potential<sup>10</sup><sup>10</sup>10 In this revised version, we have also included the bottom Yukawa contributions for completeness. All analyses in the main text use only the top Yukawa contributions as in the previous version.: $`(16\pi ^2)^2(V_{2t}(h_1,h_2)+V_{2b}(h_1,h_2))=`$ $`[3h_t^2\{{\displaystyle \underset{n=1}{\overset{4}{}}}{\displaystyle \frac{D_n}{2}}[L(m_{H_n^0}^2,m_t^2,m_t^2)\pm 2m_t^2I(m_{H_n^0}^2,m_t^2,m_t^2)+{\displaystyle \underset{i=1}{\overset{2}{}}}J(m_{\stackrel{~}{t}_i}^2,m_{H_n^0}^2)]`$ $`+{\displaystyle \underset{n=1}{\overset{2}{}}}D_{n+2}\left[s_t^2J(m_{\stackrel{~}{t}_1}^2,m_{H_n^+}^2)+c_t^2J(m_{\stackrel{~}{t}_2}^2,m_{H_n^+}^2)+c_b^2J(m_{\stackrel{~}{b}_1}^2,m_{H_n^+}^2)+s_b^2J(m_{\stackrel{~}{b}_2}^2,m_{H_n^+}^2)\right]`$ $`+s_t^2\left[c_b^2J(m_{\stackrel{~}{t}_1}^2,m_{\stackrel{~}{b}_1}^2)+s_b^2J(m_{\stackrel{~}{t}_1}^2,m_{\stackrel{~}{b}_2}^2)\right]+c_t^2\left[c_b^2J(m_{\stackrel{~}{t}_2}^2,m_{\stackrel{~}{b}_1}^2)+s_b^2J(m_{\stackrel{~}{t}_2}^2,m_{\stackrel{~}{b}_2}^2)\right]`$ $`+s_{2t}^2{\displaystyle \underset{i=1}{\overset{2}{}}}J(m_{\stackrel{~}{t}_i}^2,m_{\stackrel{~}{t}_i}^2)+c_{4t}J(m_{\stackrel{~}{t}_1}^2,m_{\stackrel{~}{t}_2}^2)+L(m_{\stackrel{~}{t}_1}^2,\mu ^2,m_t^2)+L(m_{\stackrel{~}{t}_2}^2,\mu ^2,m_t^2)\}`$ $`+3\{(h_t^2s_t^2+h_b^2c_t^2)L(m_{\stackrel{~}{t}_1}^2,\mu ^2,m_b^2)+(h_t^2c_t^2+h_b^2s_t^2)L(m_{\stackrel{~}{t}_2}^2,\mu ^2,m_b^2)`$ $`2\mu m_bh_th_bs_{2t}[I(m_{\stackrel{~}{t}_1}^2,\mu ^2,m_b^2)I(m_{\stackrel{~}{t}_2}^2,\mu ^2,m_b^2)]\}`$ $`+\{h_th_b,m_tm_b,m_{\stackrel{~}{t}_i}m_{\stackrel{~}{b}_i},\theta _{\stackrel{~}{t}}\theta _{\stackrel{~}{b}},D_kD_k^{}\}]`$ $`+3\{(h_t^2s_\beta ^2+h_b^2c_\beta ^2)L(m_{G^+}^2,m_t^2,m_b^2)+(h_t^2c_\beta ^2+h_b^2s_\beta ^2)L(m_{H^+}^2,m_t^2,m_b^2)`$ $`+2m_tm_bh_th_bs_{2\beta }[I(m_{G^+}^2,m_t^2,m_b^2)I(m_{H^+}^2,m_t^2,m_b^2)]\}`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \underset{i,j=1}{\overset{2}{}}}{\displaystyle \underset{n=1}{\overset{4}{}}}\lambda _{H_n^0\stackrel{~}{f}_i\stackrel{~}{f}_j}^2I(m_{H_n^0}^2,m_{\stackrel{~}{f}_i}^2,m_{\stackrel{~}{f}_j}^2)3{\displaystyle \underset{i,j=1}{\overset{2}{}}}{\displaystyle \underset{n=1}{\overset{2}{}}}\lambda _{H_n^+\stackrel{~}{f}_i\stackrel{~}{f}_j}^2I(m_{H_n^+}^2,m_{\stackrel{~}{f}_i}^2,m_{\stackrel{~}{f}_j}^2),`$ (D.6) where, in the first line of Eq. (D.6), positive and negative signs apply to $`𝒞𝒫`$-even ($`H^0,h^0`$) and odd ($`A^0,G^0`$) Higgs/Goldstone bosons respectively, and in the last line $`\stackrel{~}{f}_i=\{\stackrel{~}{t}_i,\stackrel{~}{b}_i\}`$. We also have, for $`n=1,2,3,4`$, $`D_n=\{s_\alpha ^2,c_\alpha ^2,s_\beta ^2,c_\beta ^2\}`$ and $`D_n^{}=\{c_\alpha ^2,s_\alpha ^2,c_\beta ^2,s_\beta ^2\}`$. The ordering of the Higgs/Goldstone bosons is $`H_n^0=H^0,h^0,G^0`$ and $`A^0`$ for $`n=1,2,3,4`$ and $`H_n^+=G^+,H^+`$ for $`n=1,2`$. Two tests can be applied to check the correctness of the effective potential $`V(h_1,h_2)`$. First, the potential should vanish in the supersymmetric limit (i.e., when all soft-breaking parameters are taken to be zero), and second, the potential $`V(h_1,h_2)`$ should be invariant under changes of the renormalization scale, up to the order of our perturbative calculation. The vanishing of the potential in the supersymmetric limit is proved by simple algebra. In the following we show the invariance of the two-loop potential under a RG transformation. Using the derivatives of $`I`$, $`J`$ and $`L`$ functions with respect to the renormalization scale $`Q`$ $`{\displaystyle \frac{I(m_1^2,m_2^2,m_3^2)}{\mathrm{ln}Q^2}}={\displaystyle \underset{i=1}{\overset{3}{}}}\left[A_0(m_i^2)+m_i^2\right],`$ (D.7) $`{\displaystyle \frac{J(m_1^2,m_2^2)}{\mathrm{ln}Q^2}}=m_1^2A_0(m_2^2)+m_2^2A_0(m_1^2),`$ (D.8) $`{\displaystyle \frac{L(m_1^2,m_2^2,m_3^2)}{\mathrm{ln}Q^2}}=(m_1^22m_2^22m_3^2)A_0(m_1^2)m_2^2A_0(m_2^2)m_3^2A_0(m_3^2)`$ $`+m_1^4(m_2^2+m_3^2)^2,`$ (D.9) and the one-loop MSSM RGEs for top-(s)quark and Higgs boson masses, Eqs. (B.18-B.22), we find that the RG variation $`{\displaystyle \frac{V_2}{\mathrm{ln}Q^2}}𝒟^{(1)}V_1`$ $`=`$ $`{\displaystyle \frac{8g_3^2h_t^2h_2^2}{(16\pi ^2)^2}}\left(M_{\stackrel{~}{Q}}^2+M_{\stackrel{~}{U}}^2+2M_3^2+X_t^22M_3X_t\right)`$ (D.10) $``$ $`{\displaystyle \frac{9h_t^4h_2^2}{(16\pi ^2)^2}}\left\{M_{\stackrel{~}{Q}}^2+M_{\stackrel{~}{U}}^2+{\displaystyle \frac{1}{2}}\left[m_{H_2}^2+(A_t+X_t)^2\right]\right\}`$ modulo terms independent of the Higgs field $`h_2`$. Here $`𝒟^{(1)}V_1`$ represents the one-loop RG variation of the one-loop potential Eq. (D.3). This result agrees exactly with the two-loop RG variation of the tree-level potential $`𝒟^{(2)}V_0`$ \[cf. Eqs. (D.2) and (B.15-B.17)\], so that $$\frac{d}{d\mathrm{ln}Q^2}(V_0+V_1+V_2)𝒟^{(2)}V_0+𝒟^{(1)}V_1+\frac{V_2}{\mathrm{ln}Q^2}=0.$$ (D.11) Note that this is a nontrivial check that all $`\mathrm{ln}Q^2`$ terms cancel with each other between Eq. (D.10) and $`𝒟^{(2)}V_0`$; this cancellation guarantees the correct leading and next-to-leading order logarithmic behavior of the effective potential. To derive the analytical expression of $`\mathrm{\Delta }m_{h^0}^2`$ in Sec. 2, we need to expand the two-loop potential $`V_2`$ in powers of $`m_t/M_S`$ and $`m_tX_t/M_S^2`$; besides many straightforward expansions, we have used (A.25)-(A.28) for the $`t\stackrel{~}{q}\stackrel{~}{h}`$ and $`\stackrel{~}{t}\stackrel{~}{q}h`$ diagrams of (D.6), with $`\stackrel{~}{q}=\stackrel{~}{t}`$ or $`\stackrel{~}{b}`$.
warning/0003/hep-ex0003015.html	ar5iv	text	# References PRESENT SITUATION OF DIFFRACTED X-RAY RADIATION AND RESONANCE (COHERENT) TRANSITION RADIATION INDUCED BY HIGH ENERGY CHARGED PARTICLES IN FREQUENCIES REGION EXIDING ATOMIC ONE. M.L. Ter-Mikayelyan Institute for Physical Research of Armenian Academy of Science 378410, Ashtarak-2, Armenia. ## Abstract The review is devoted to the modern investigations of electromagnetic radiation by relativistic charged particles propagating with constant velocity through the periodic media. Two examples of radiation are considered in this review, the first corresponds to the Bragg scattering of charged particles pseudophotons in crystals, the second one to the Fresnel scattering of pseudophotons in periodic medium. Both examples play essential role in construction new compact tunable sources in X-ray region. 1. Introduction High-energy electromagnetic processes in medium have been summarized in my monograph in 1969 year and translated into English in 1972 . (The numbering of chapters, paragraphs, formulas and total text of English edition are the same as in Russian version.) Since that time this field of high-energy physics developed catastrophically fast in many various directions, including quantum chromodynanic (see and references therein). During that period of time many original publications, review articles and books have been published. The number of original publications is impossible to estimate. Only the physicists from former Soviet Union published in Physics Usphekhi since 1972 ten reviews and eleven monographs in various publishing houses. Numerous conferences, workshops and seminars in different countries were devoted to these problems. For example: traditional conferences in Russia (Tomsk, RREPS1993-1995-1997-1999), (see and references therein), in Germany (Tabarz, 1998) etc. All these problems of high-energy radiation physics are based on the following underlying idea; the length of trajectory (coherence length) along the trajectory of initiating reaction particle increases with energy of incoming particle and the directionality of process (see Appendix). The history of this concept has been published in Russian by E.L. Feinberg. in his well known paper in Priroda , (see translation of into English in Appendix). Among numerous problems from this field I will review in this paper only two; the diffracted X radiation (DXR) and the resonance transition radiation (RTR). 2. Diffracted X Radiation (DXR) DXR has been considered in 1969 and included in Chapt.5 of book . The theoretical consideration can be very simplified, if we follow the corresponding theory of X-ray scattering in crystals. Expanding electrical field of fast moving particle in time and variable part of dielectric susceptibility $`\epsilon _1(r)`$ over vectors $`\stackrel{}{g}`$ of the reciprocal lattice. $$\epsilon =\epsilon _0+\epsilon _1(\stackrel{}{r});\epsilon _1(\stackrel{}{r})=n_\stackrel{}{b}\mathrm{exp}(i\stackrel{}{g}\stackrel{}{r})$$ (1) and using the theory of X ray scattering we get the following expression for scattered electric field at a distance $`R_0`$ $$\stackrel{}{E_\omega ^{^{}}}(\stackrel{}{R_0})=\frac{e}{\epsilon _0R_0}\underset{\stackrel{}{b}}{}n_\stackrel{}{b}\left[\stackrel{}{k^{^{}}}\times \stackrel{}{k^{^{}}}\times \left(\frac{\omega \stackrel{}{v}}{c^2}+\frac{\stackrel{}{g}}{\epsilon }\right)\right]\frac{\delta [\omega (\stackrel{}{k^{^{}}}\stackrel{}{g})\stackrel{}{v}]}{(\stackrel{}{k^{^{}}}\stackrel{}{g})^2(\omega ^2/c^2)\epsilon _0}.$$ (2) Polarization of D.X.R. is linear and is given by expression (2). The radiation angle is determined, assuming the argument of the $`\delta `$-function in (2) to be equal to zero $$\mathrm{cos}\theta =\frac{c}{v\sqrt{\epsilon _0}}+\frac{(\stackrel{}{g}\stackrel{}{v})c}{\omega v\sqrt{\epsilon _0}}$$ (3) where $`\theta `$ is the angle between velocity and scattered photon. Expression (3) often used in practice, for measuring the dependence of emitted photon energy upon the scattering angle $`\theta `$. The equality (3) follows from energy-momentum conservation laws for photon radiation in a crystal if we take into account, that crystal can receive momentum inverse proportional to the length of periodicity. Intensity is given by following expressions $$d\stackrel{}{I}_{\omega ,\stackrel{}{n}}=c\sqrt{\epsilon _0}\|\stackrel{}{E}_\omega ^{^{}}\|^2R_0^2d\omega d\mathrm{\Omega }$$ (4) $$d\stackrel{}{I}_{\omega ,\stackrel{}{n}}=\frac{e^2\omega ^2T}{2\pi \epsilon ^{5/2}c}\underset{\stackrel{}{b}}{}n_\stackrel{}{b}^2\left\|\stackrel{}{k^{^{}}}\times \left(\frac{\omega \epsilon }{c^2}\stackrel{}{v}+\stackrel{}{g}\right)\right\|^2\times \frac{\delta [\omega (\stackrel{}{k^{^{}}}\stackrel{}{g})\stackrel{}{v}]}{[(\stackrel{}{k^{^{}}}\stackrel{}{g})^2(\epsilon _0/c^2)\omega ^2]^2}d\omega d\mathrm{\Omega }.$$ (5) Formulae for polarization (2), angular distribution (3), and the radiated energy in frequency interval $`d\omega `$ (5) are the base of kinematical theory and have been confirmed in numerous theoretical and experimental papers. The first experimental investigations of DXR have been done in Tomsk, by group of A.P. Potylitsyn , the second in Yerevan , third in Kharkov . DXR initiated by protons has been observed in . In papers was shown that kinematical theory may be sufficient to explain modern experimental results on spectral and angular distributions of DXR, as well as absolute differential yield. In papers quantum theory of DXR has been developed and has been shown, that for relativistic particles, if recoil due to the photon emission is small, the quantum expressions coincides with corresponding expressions presented in . I will mention that like dynamical theory of X-ray scattering, the corresponding theory in DXR has been developed in series of publications and included in monographs . This type of radiation is referred in literature under names quasi-cherencov, parametric and even polarization radiation. To simplify the terminology I will use for this type of radiation the name dynamical theory of DXR The name PXR often used in literature for both type of radiation is meaningless and doesn’t correspond to the nature of radiation (I thank H.Nitta for this comment). During the last years new experimental investigation has been accomplished, which improved our knowledge of DXR. The remarkable features of DXR (monochromatic with continuously variable wavelength, propagation direction well separated from the electrons beam one, small energy and angular spread of the order of magnitude one over gamma, and at last the small sizes of setup) suggest that DXR can be an perspective X- ray source in future. A series of paper has been published recently by joint group from Institute fuer Kernphysik in Darmstadt, from Kharkov, Rossendorf and Johannesburg, using superconducting linear accelerator S-DALLINAC with electron energy below 10 MeV. The line width of 8.98 keV DXR for electron energy 6 Mev has been measured applying an absorption technique using a copper foil and tuning the energy of the DXR peak across the K-absorption edge of copper. The spectral density in the peak deduced from experiment was, $`I=0.95\times 10^7`$ photons/(electron sr.eV), and linewidth 48 eV (see Fig.1). In paper upper limits of line width of DXR 1.2 eV and 3.5 eV have been determined for (111) and (022) reflections of silicon at photon energies of 4966 eV and 8332 eV. Investigations of the line width of DXR at the Mainz Microtron MMI a relative energy width ($`\mathrm{\Delta }E/E=10^5`$ should be reached for the silicon (333) reflection . The spectral and angular distribution of DXR have been studied mostly in silicon and diamond crystals over a range from a few MeV up to several GeV of electrons and are in consistent with the theory . The last experiments carried out by physicists from Germany (Werner Heisenrberg Institute and Institute fuer Kernphysik), has shown close to $`100\%`$ linearly polarization at every single point of photons angular distribution, with agreement with the theory . In these investigations polarization has been analyzed by means of novel method of polarimetry exploiting directional information of the photoeffect in a charge coupled device consisting of $`1.3\times 10^6`$ square pixels of 6.8mkm . The advent of such devices opens a promising route towards a universal X-ray detector for simultaneous imaging, spectroscopy and polarimetry. The angular distributions of DXR polarization directions, calculated recently on the base of the theory , are close to the experimental and calculated data presented in for DXR in forward and backward hemispheres, but in disagreement with calculations for DXR polarization at right angle. The disagreement between calculations is due because in the longitudinal density effect has been neglected (private remark from A.V. Schagin). The discrepancy of both calculations with experiment for polarization in forward hemisphere emission remains unsolved yet and new measurement is needed (private communication from R. Kottaus). The physicist from Tomsk investigated influences of temperature on DXR intriducing Debay-Waller factor in expression (5). They obtained good agreement with experimental data . In first publication devoted to the influence of acoustic waves and gradient of temperature on DXR shows that the intensity of DXR may be increased several times (see Fig.2). Nevertheless intensity of DXR attained in laboratories in several keV domains is the same order of magnitude as synchrotron radiation of big accelerarators (see and references therein). Concluding this short review of DXR I will notice that more complicated theoretical and experimental problems remains unsolved. In particular the region of applicability of dynamical DXR theory and its correspondence with experiments has not been investigated seriously. During the International Workshop on Radiation Physics, in Tabartz Prof. Baryshevsky V. affirmed that dynamical version of DXR is necessary to understand experimental data . On the other hand Prof. N. Nasonov maintains opposite statement.I. Feranchuk and A. Ivashin incorporated quantitatively electron multiple scattering and photon absorption in kinematics theory , but more subtle theoretical treatment is necessary. 3. Resonance Transition Radiation (RTR) The well known expression for Transition Radiation (TR) introduced in physics in 1946 by V. Ginzburg and I. Frank received a new development, when it was investigated in radiation frequencies exceeding optical . At that time the longitudinal density effect and coherence length concept introduced in high-energy radiation processes in papers were well known and the results of papers can be easily understood and derived from . The problems arise when the expression for transition radiation, which is valid only for one interface, tried applying for many periodically spaced interfaces and in limiting case for periodic medium. It must be taken into account, that nonrelativistic charged particles propagating through periodic medium will emit photons with frequencies proportional to the frequency of propagation the periodicity of medium (resonance condition). For relativistic particles because of Lorenz transformation we get the following resonance condition $$\mathrm{cos}\theta =\frac{c}{v\sqrt{\epsilon _0}}\frac{2\pi rc}{l\omega \sqrt{\epsilon _0}}\mathrm{cos}\psi $$ (6) where ($`\omega `$-frequency of radiation, v -velocity of particle, ($`\epsilon _0`$ \- effective dielectric susceptibility, ($`\theta `$ \- angle between incoming particle velocity and direction of radiation, ($`\psi `$ \- angle of incidence of charged particle onto the one-dimensional periodic medium, $`l`$-period of medium, $`r`$ \- number of emitting harmonic. The resonance condition (6) can be easily derived using the energy-momentum conservation laws in periodic medium. This kind of TR was termed as RTR. For $`l`$ we get the well-known Tamm-Frank expression for Vavilov-Cherenkov radiation in homogeneous medium. RTR consist of overlapping radiated harmonics, each has its threshold in energy of radiating particle and depends on parameters of medium. Theory of RTR has been published in my article presented for publication by L. Landau to Dokladi Acad. Nauk in 1960 and published in Nuclear Physics in 1961 . For periodic medium radiation on $`r`$-harmonic appears when the particle velocity exceeds the group velocity of corresponding photons.<sup>1</sup><sup>1</sup>1Beginning from I. Frank proposal, many physicists tried to increase the TR intensity from one interface using many foils. The calculations in optical region were cumbersome and negaive, because of neglecting resonance condition. This problem was similar to the corresponding one in saturation problem of ionization losses solved by E. Fermi, who enlarged the I. Tamm calculation for Cherenkov-Vavilov radiation. The most convenient medium for experimental investigation is laminar periodic medium with many plates (Fig.3). For laminar medium consisted with two different plats $`\sqrt{\epsilon }_0`$ has simple form $$\sqrt{\epsilon }_0=\frac{l_1\sqrt{\epsilon _1}+l_2\sqrt{\epsilon _2}}{l}$$ (7) where $`l_1`$ and $`l_2`$ are thickness, $`l=l_1+l_2`$ and $`\epsilon _1`$ and $`\epsilon _2`$ are dielectric susceptibilities of plates. For frequencies much more higher optical frequencies from inequality $`\|\mathrm{cos}\theta \|1`$, for each harmonic we get $$\omega _{max}=\frac{4\pi cr}{l}\left(\frac{E}{mc^2}\right)^2\omega \frac{l\omega _0^2}{4\pi cr}=\omega _{min},$$ (8) where $`\omega _0`$ is the plasma frequency $$\omega _0^2=\frac{4\pi NZe^2}{m_e}$$ (9) and for laminar medium $$NZ=(N_1Z_1l_1+N_2Z_2l_2)/l$$ (10) From (8) we get the threshold energy for radiation harmonic of $`r`$-number. The intensity of RTR is given by following expression (see formula $`(28.92^{})`$ from or paper ). $`dI_{\omega ,\theta }`$ $`=`$ $`{\displaystyle \frac{e^2\theta ^3d\theta d\omega }{2\pi c}}\|{\displaystyle \frac{\epsilon _2\epsilon _1}{\left(1\frac{v}{c}\sqrt{\epsilon _1}\mathrm{cos}\theta \right)\left(1\frac{v}{c}\sqrt{\epsilon _2}\mathrm{cos}\theta \right)}}\|^2\times `$ (11) $`\times \mathrm{sin}^2\left[{\displaystyle \frac{l_1\omega }{2c}}\left(1{\displaystyle \frac{v}{c}}\sqrt{\epsilon _1}\mathrm{cos}\theta \right)\right]{\displaystyle \frac{\mathrm{sin}^2\frac{n\beta }{2}}{\mathrm{sin}^2\beta }}.`$ where the first term corresponds to Ginzburg-Frank transition radiation, the second corresponds to the interference of radiation from two interfaces of one plate and the last term corresponds to the coherent summation of radiation from n-plates. The quantity $`\beta `$ equals $$\beta =\left(1\frac{v}{c}\sqrt{\epsilon _1}\mathrm{cos}\theta \right)\frac{\omega l_1}{2v}+\left(1\frac{v}{c}\sqrt{\epsilon _2}\mathrm{cos}\theta \right)\frac{\omega l_2}{2v}$$ (12) If the number of plates increases the last term can be substituted by delta function and we get the resonance condition (6). The theory of RTR depends dramatically on coherence length, if for example coherence length exceeds the distances between two interfaces in a plate the radiation from two interfaces must be summed coherently. The same result takes place for total periodic medium. Radiation from plates will sum independently if the distances between plates exceed the coherence length. We shall discuss the related experiments later. The RTR including absorption influence and multiple scattering effects has been discussed in . But in that time (sixteenths years) we were interested to apply RTR for construction a new type of counters for very high energies of particles were Cherencov detectors were insensitive. These new detectors has been constructed by group of F.R. Arutunian in 1963 . In following investigations the property of these detectors were improved and reviewed in many publications . They are used now for identification of particles in modern high-energy accelerators (see for example and references therein). Since 1985 RTR received a new impact for developing in different domain of physics. Joint group of physicist from Stanford University and Livermore Lawrence Laboratory investigated RTR using the linear accelerator with energy of electrons equal to 17.2, 25, 54 MeV to produce photons in keV region . Stack of Be, C, Al foils consistent each from18 foils with thickness 1 $`\mu `$m separated in distance 0.75mm (for carbon) and 1.5mm (Be and Al) have been used. Experimental data for RTR intensity angular and spectral distributions presented in confirms the theory of interference at the interfaces of a single foil. The authors assert that an easy-to-tune source of intense polarized monochromatic radiation holds much promise for submicron lithography. For instance, a 0.5 $`\mu `$m resolution was reported in . Though in \[47-49\] no interference effects were observed with radiation from different beryllium foils (see Fig.4), the same authors noticed an unusual interference pattern in . Interesting observations of interference effects in RTR were made by French researchers at the electron accelerator in Sacle (see Fig.5). The achievement of 0.3 $`\mu `$m resolution was reported in . In the soft range of the spectrum (1-3 keV) the RTR spatial distribution for electron energy of 50-228 MeV was observed in . Here attention is drawn to the fact that RTR results through the whole radiator stack, which forms the periodic structure, and concentrates in the solid cone whose angle grows with electron energy. Changes in the periodic structure parameters (e.g. electron energy, structure size) suppress interference effects and give rise to TR for which the emission angle, in coincidence with the TR theory, is inversely proportional to the electron energy (see Fig.6). Teams from universities of Kyoto, Tohoku, Hyrosimo, Tokyo (I. Endo et al.) and Tomsk Institute for Nuclear Physics (A.P. Potylitsin et al.) in cooperation with various Japanese firms have made a lot of research into RTR at accelerators in Japan \[54-58\]. The goal of the research was not only to investigate resonance effects in RTR but also to determine the parameters of the electron beam and active medium best suitable for practical implementations of RTR. 4. The radiation of moving particles on complex structures (DXR+ RTR) The use of complex periodic structures was first suggested in the 90s when it became clear that the development of effective kiloelectron-volt generators requires increased intensities of DXR+ RTR . Russian and Japanese physicists joined efforts to conduct research in this area. Papers present the results of the irradiation of a target consisted of three crystals 16 (m thick with 800-MeV electrons (the synchrotron in Tomsk (Fig.7a)) and with 900-MeV electrons (the linear accelerator in Tokyo). Besides DXR, in the first crystal layer of the stack the electron beam generates RTR, which undergoes Bragg diffraction on the following crystal layers. This gives rise to the effective growth of emission (a 1.7-times increase was observed in the experiment). It should be noted that with the emission angle of the same order the RTR intensity and spectral width is much greater than that of DXR. The difference is that RTR follows the electron path, while DXR propagates along Bragg’s angles of refraction from the corresponding crystal planes in the reciprocal lattice. The authors introduced a new name for this type of emission - parametric (diffracted) RTR (substituting letter P by D (DRTR)), which we will further keep to. The next experiment dealt with a 900-MeV electron beam and a target consisted of ten mylar foils and graphite crystal (Fig. 7b). The authors assert that even with a few foils DRTR follows Bragg’s angles and is much more intensive than DXR. The last joint works of these authors investigates radiation in the keV spectral range in a periodic medium consisting with crystalline plates. A 900 MeV electron beam and a target of 1 to 100 plates of monocrystal silicon were used in the experiment (see Fig.7c). The DRTR intensity of 35.5 keV photons proved to be comparable to that of synchrotron emission caused by a 1.7 GeV electron beam. The paper also considers the relationship between the radiation intensity and the number of plates - the issue that was discussed earlier in . Papers , and (the last citation refers related papers presented at the meeting in Tabarz, Germany, 1998) build a theory that establishes a link between RTR and diffraction radiation that caused by a charged particle flying over a surface with periodic irregularities. In the experiment an electron beam propagated over a GaAs plate whose surface had 300 identical strips which were 10 $`\mu m`$ wide, 100 $`\mu m`$ high and spaced 33 $`\mu m`$ apart. The authors observed radiation that consisted of DXR and DRTR, the intensity of the latter being much higher. A great deal of theoretical papers discussing interference of various kinds of radiation has been published recently. Paper offers a method of separating DXR and DRTR. Paper shows that DXR output in mosaic crystal is the same as in a perfect crystal, and DRTR output is much higher. Diffraction of TR on a crystal structure is considered theoretically in . Several relevant theoretical papers were also presented at the recent international conferences . 5. Conclusion As I have already noted in Introduction, I have elucidated, out of a large set of questions, only DXR and RTR, which are developed recently in the numerous physical community. I hope to focus on other problems of yigh energy electromagnetic processes in medium in my forthcoming reviews. Aouthor will be very grateful for any comments and suggestions to improve this review. 6. Acknowledgments I am very grateful to Organizing Committee of RREPS’99 for the invitation to take part in the traditional symposium on the lake Baikal organized by Tomsk Nuclear Physics Institute. I would like to express a special gratitude to A.P. Potylitsyn, Yu. L. Pivovarov and L. Puzyrevich whose hospitality is very hard to overestimate. The work was performed within the program of the Ministry of Education and Science of the Republic of Armenia, grant $`\mathrm{\#}96772`$ and the INTAS grant $`\mathrm{\#}99392`$. APPENDIX EFFECT CONFIRMED 40 YEARS LATER. ”Nature”, Russian Academy of Sciences, 1994, N1, pp. 30-33 E.L. Feinberg Corresponding member of RAS, P.N. Lebedev Physical Institute of RAS It was recently reported that in the accelerator center of Stanford University (SLAC) direct experimental evidence was obtained of the suppression of bremsstrahlung of relativistic particles in an medium \[ 1\], theoretically developed by L.D. Landau and I.Ya. Pomeranchuk 40 years ago . This experiment confirms already the third of the important effects predicted in a series of works of Soviet theoreticians in 1952-1954. All these effects are bound by a common physical idea (or a basis), although they are displayed in different interactions of high-energy particles, and not only electromagnetic, but, as well, nuclear. This basis was built in 1952 in the Ph.D-thesis of M.L. Ter-Mikayelian , the post-graduate, of that time, in the Theoretical Department of the P.N. Lebedev Physical Institute, AS of USSR. The work was devoted to the investigation of the bremsstrahlung but not on single atom, as was studied before; it was considered in a medium, specifically, in a crystal. The result of this work which seemed at that time paraamorphousdoxical, consisted in a statement that at very high energies, when the wavelength of either the emitted photon or the electron is tens of millions, milliards times shorter than the mean interatomic distance in the medium, the usual radiation pattern changes dramatically. Particularly, if the motion occurs along the crystal axis, this radiation may many times exceed on individual atoms. The process in this case is widely extended in space and includes a domain with characteristic sizes many times exceeding the interatomic separations. All atoms of the crystal in this domain act coherently, and as a result, the radiation is enhanced significantly. This length was, naturally, termed as the coherence length. The work under discussion had forerunners. A possibility of the influence of crystalline structure on the bremsstrahlung of fast particles was discussed by B. Ferretti in 1950, and still earlier, in 1935, by E.Williams, who developed independently the well-known, in theoretical physics, Weizsecker-Willians method. However, either the work of Ferretti or the notation of Williams (who obtained, by the way, an incorrect sign of the effect) remained unpersuasive and did not attract the attention until Ter-Mikayelian succeeded to show that the bases of the process are paradoxical (for that time) physical causes which turned out to have much wider significance than the explanation of radiation in crystals. This lead to the development of a new direction involving various processes in high-energy physics. Now attempts are made to apply these results to processes of of high-energy hadrons inside the nucleus regarded as a material medium. Coherence length Ter-Mikayelian succeeded to show that in the processes at high energies where the particles are scattered at small angles (decreasing with the icrease of the energy), the longitudinal momentum, $`q_{\|\|}`$, transferred to the target, drops, and, consequently, according to the uncertainty relation, $`\mathrm{\Delta }q_{\|\|}\mathrm{\Delta }x_{\|\|}\mathrm{}`$ , the longitudinal distance $`\mathrm{\Delta }x_{\|\|}`$ involved in the process increases with the energy. Therefore, it is the coherence length, $`L_{coh}\mathrm{}/q_{\|\|}`$ , rather than the wavelength of the particle, that can be a measure of the size of domain, relevant for the effect. When speaking not especially of a crystal it would be reasonable to term this length otherwise, namelythe zone or length of process formation. In 1952 this looked incredible and even absurd since it was accepted that characteristic distances of formation of electromagnetic processes are of the same order as the wavelengths of particles involved (or the atomic sizes). In order to illustrate this nontrivial result let us consider the process of bremsstrahlung of a photon with energy $`\mathrm{}\omega `$ and momentum $`\mathrm{}\stackrel{}{k}`$ on a fixed coulombian center. Let $`E_1`$ and $`\stackrel{}{p_1}`$ be the initial energy and momentum of the radiating relativistic particle of mass m while $`E_2`$ and $`\stackrel{}{p_1}`$ are the same quantities in its final state. Let us then use the energy and momentum conservation laws, $`E_1E_2=\mathrm{}\omega (1)`$ $`\stackrel{}{p_1}\stackrel{}{p_2}\mathrm{}\stackrel{}{k}=\stackrel{}{q}(2)`$ Where $`\stackrel{}{q}`$ is the momentum transferred to the nucleus, and project the latter on the initial direction of particle’s motion. For this purpose we multiply Eq. (2) by the initial velocity of the particle, $`\stackrel{}{v_1}`$: $`\stackrel{}{v_1}\delta \stackrel{}{p}\mathrm{}\stackrel{}{k}\stackrel{}{v_1}=\stackrel{}{v_1}\stackrel{}{q}\|\stackrel{}{v_1}\|q_{},`$ Where $`\delta \stackrel{}{p}=\stackrel{}{p_1}\stackrel{}{p_2}`$, and $`q_{}`$ is the longitudinal momentum transferred along the motion of the emitting particle. Since $`\stackrel{}{v}\delta \stackrel{}{p}=\delta E=E_1E_2=\mathrm{}\omega `$, we have for small energy variation, $`\mathrm{}\omega E`$, of the initially ultrarelativistic particle ( $`\|\stackrel{}{v_1}\|\|\stackrel{}{v_2}\|c`$ and $`k=\omega /c`$), $`q_{}=\frac{\mathrm{}\omega }{c}\left(1\frac{v}{c}\mathrm{cos}\theta \right)`$. (3) Where $`\theta `$ is the angle between the emitted and the direction of motion, $`\stackrel{}{v_1}`$ , of the emitting particle. As the radiation at high energies is known to be sharply directed, the obtained formula should be considered as small $`\theta `$ . For $`\theta \sqrt{1v^2/c^2}`$ $`q_{}=\frac{\mathrm{}\omega }{c}\left(1\frac{v}{c}\right)`$. (4) As it was mentioned above, the coherence length (formation zone) along the path of the emitting particle amounts by the order of magnitude to $`L_{coh}\frac{\mathrm{}}{q_{}}\frac{c}{\omega \left(1v/c\right)}\frac{E^2}{m^2c^3\omega }`$. (5) With the time of passing the zone being equal to $`t_{coh}\frac{L_{coh}}{v}\frac{E^2}{m^2c^4\omega }`$. (6) This means that at $`vc`$ the quantity $`L_{coh}`$ may reach macroscopic values. (For large variations of energy of the emitting particle, $`L_{coh}`$ is given by an expression like (5). Paradoxically of this result is up to now being emphasized in review articles , although the results raises no doubts. However, in 1952 it was not at once that one succeeded to convince even Landau and Pomeranchuk of the correctness of these arguments, of the presence of the formation zone increasing with energy, and so on . Nevertheless, already in autumn 1952 Ter-Mikayelian reported at Landau’s seminar the details of his thesis, with complete mutual understanding and approval. It should be added only that the described effect was completely, in all theoretically developed details, checked experimentally in a crystalline medium, ten years later. At present it is used, in particular, to obtain quasimonochromatic and polarized ($`\gamma `$-quanta from electron accelerators . The importance of the arisen conception of the length of formation was at once estimated by Landau and Pomeranchuk, and they (and not only they) began to think to further theoretically develop this phenomenon. Influence of multiple scattering on the bremsstrahlung in amorphous medium. First Pomeranchuk noticed to Ter-Mikayelian that if all what he said about the coherence length in the crystal is correct, then in an amorphous medium as well, the traditional Bethe-Hilter formula for the bremsstrahlung on a single atom shuld have been changed, due to the absorption on a distance termed the radiation length, at $`L_{rad}L_{coh}`$. This statement raised no objections of either Ter-Mikayelian or Landau who advised to evaluate this effect. Soon, however, after examining the problem, Landau came to the conclusion that the influence of multiple scattering will take place rather than the influence of absorption by the emitting particle. Rather soon Landau and Pomeranchuk evaluated this effect and acquainted Ter-Mikayelian with the manuscript of their joint article asking to tell his remarks. A discussion of this work took place, and the article was approved. It happened so that Landau and Pomeranchuk, starting with the formula (5), and explaining, that, in accordance with the formula (6), the time tcoh ”does very sharply increase with the energy and, as a cosequence, those distances between electron and nucleus play a role which significantly exceed atomic sizes”, did not, apparently, by a misunderstanding only, refer to the work of Ter-Mikayelian. But he would not think (felt shy?) to tell them that it had to be done . This lead to such a moving of events that this story seems to be not only quite appropriate here but also instructive from the point of view of the scientific ethics. In that time conditions of isolation of Soviet Science a publication of our works in foreign languages was strongly prohibited as ”cringing to abroad”. Nevertheless, Landau’s name in a published, even in Russian, article, attracted the attention of the famous American physicist Dyson. Having, naturally, not known about the Ter-Mikayelian’s work Dyson suspected that an interesting effect should exist in a crystal, and published (in coauthorship with G. Ueberall) a paper in an american journal presenting a result coinciding with that of Ter-Mikayelian (and refered, of course, not to him but to Landau). Learning this Landau sent urgently the reprints of Ter-Mikayelian’s works to Dyson, USA, showing that the work he published had already been done here. In a reply letter Dyson appraised highly the works of Ter-Mikayelian and recognized that he with Ueberall obtained the same physical results using, however, another calculation technique. Landau there and then acquinted the participants of the next seminar with the Dyson’s letter. Being elegant and clear physically the work of Landau and Pomeranchuk needed some mathematical improvements. This has been done by A.B. Migdal who used a fine and original technique to complete the Landau-Pomeranchuk theory, from quantitative point of view, to a logically closed form, and obtained an expression for the bremsstrahlung in amorphous medium with allowance for the influence of multiple scattering. This expression used sometimes to be termed the Landau-Pomeranchuk-Migdal formula. Our physicists used this formula frequently to calculate the development of broad electromagnetic showers of cosmic rays. It is just this formula that was recently confirmed experimentally in Stanford. Longitudinal density effect. With these works the investigations of pecuiliarities of the radiation of ultrafast particles did not stop. Ter-Mikayelian generalized very soon the work of Landau-Pomeranchuk in he sense that he took into account the role of the dielectric polarization of the amorphous medium . As it turned out, this polarizations affects the radiation of ”soft” quanta with the energy of the order of or less than $`\mathrm{}\omega _{crit}=\mathrm{}\omega _0/\sqrt{1v^2/c^2}.`$ (7) Here $`\omega _0^2=4\pi NZe^2/m`$ is the squared plasma frequency, $`N`$ the number of atoms per $`cm^3`$, $`m`$ and $`e`$ the electronic mass and charge, $`Z`$ the number of electrons in the atom. The estimation of influence of the medium polarization on the formation length can readily be obtained from the above expressions by taking into account that in a medium we have actually $`k=\frac{\omega }{c}\sqrt{\epsilon }`$, with $`\epsilon `$ being the dielectric constant of the medium. For the frequencies considerably exceeding the atomic ones: $`\epsilon =1\omega _0^2/2\omega ^2`$ . Substituting correspondingly k in the expression (2), it is easy to see that the formula for the coherence length takes the form $`L_{coh}=\frac{c}{\omega }\left[1/\left(1\frac{v}{c}+\frac{\omega _0^2}{2\omega ^2}\right)\right]`$ . This length is now ”cut” for the photons at frequencies $`\omega \omega _{crit}`$. In this case, the increase in the energy of the emitting particle ($`vc`$) results in that $`L_{coh}`$ remains constant at a given frequency $`\omega `$, which leads to an essential modification of either the Bethe-Hitler or the Landau-Pomeranchuk formula in the region of very soft quanta . This density effect in the bremsstrahlung is in way similar to the density effect in ionization losses discovered by E. Fermi. The difference is following: in the second case it is the effective impact parameters (i.e. the distances in the direction perpendicular to paths of particles) that are ”cut”, while in the first case it is the longitudinal distances along the path of the emitting particle. In this connection the Stanford physicists term this phenomenon the longitudinal density effect with referring to the work of Ter-Mikayelian. Unfortunately, in the experiment performed in Stanford University, the intensity of emitted photons was measured in dependence on their energy only in the region from 5 to 500 MeV. Since the electron energy was 25 GeV the frequencies of those photons exceeded, and the longitudinal density effect could not still be displayed to the full extent. It would be interesting to conduct corresponding experiments (even at considerably lower energies of the emitting particle) for the emission spectrum of photons at frequencies of the order of or less than (crit. In principle, this method could be employed for measurements of fast particle energies which is important for the experimental physics of ultrahigh energies. Application to hadronic processes. We have already mentioned that the increase of the formation zone with the energy of relativistic particles, revealed by Ter-Mikayelian, was applied also to high-energy hadronic processes. Here three consequent effects can be mentioned. The first, so to say, preliminary, is not of particular importance and has not been checked experimentally. It is valuable mainly from the methodological point of view. With use of this effect it turned to be possible to calculate the emission of photons by a charged pion the plane wave of which is incident no required to know the details of interaction between the pion and nucleus, it is sufficient that such a nucleus cuts a round hole in the plane wave. Then a diffraction of pions occurs. And, as at small diffraction angles the length of zone of formation of emitted photon is very long, all needed integration can be made outside the nucleus and is performed without a detailed knowledge of the laws of interaction between the passing pion and nucleus . However, the further attack in the same direction lead to much more important result. A statement was made that a diffracted pion (like any hadron) can dissociate into other hadrons . For example, a diffracted nucleon can emit a pion. Of course, in this case the probability of such a diffraction dissociation, i.e. of the process of pion emission by a diffracted hadron, can be calculated only by the perturbative theory giving merely a rough estimation of the cross-section of this process. But the cross-section is again determined by the integration over a large, increasing with the energy region outside the hadron or nucleus target. This gives the process some features, which allow distinguishing it among other hadron generation processes. Prediction of the diffraction dissociation of hadrons raised doubts for a long time, but already in sixties it was confirmed experimentally, and is now of great importance in high-energy hadron physics. It was very concretely described in ”Regestic” as an exchange of pomeron, a quasiparticle with zero quantum numbers . However, in that ”preRegestic” age many unclear question arose concerning the nature of the effect which, as we saw, seemed to occur completely outside the nucleus-target. It was questioned: ”Where enters the interaction with the nucleus?” Once Pomeranchuk answered with irritation: ”Well, you can hold that a chaste conception occurred”. In order to clear the mechanism, a special work has been done concerning a similar possible effect, which is much more illustrative and calculable, that is the effect of diffraction splitting (dissociation) of a deuton . It was then developed to a new direction-diffraction splitting of nuclei. Introduction of the concept of coherence length, or formation zone, its use in various physical phenomena essentially changed our ideas about the radiation processes occurring at high energies. These are, we remind, in the first place, three effects that are under discussion here and confirmed experimentally: bremsstrahlung of photons in crystal, diffraction dissociation of hadrons and the Landau-Pomeranchuk effect in amorphous media. It should be noted that related processes have been considered earlier as well, particularly, when V.L. Ginsburg and I.M. Frank predicted the existence of transition radiation. The use of concept of coherence length increasing with energy of emitting particle, in consideration of this phenomenon permitted to enrich its theoretical description, extend it considerably into the ultrahigh energy region, and then to create new detectors of relativistic particles . All these works were done in a new years at the time when our country separated from the world science by an ”iron curtain”. When this curtain raised, the journal ”Nuovo Cimento” ordered our scientists a number of reviews of soviet investigations on various problems. It was found that very much was done originally. As to the above-mentioned questions, reviews on these topics contained already more than a dozen original; publications which resulted actually from the work of Ter-Mikayelian. Footnotes 1. See,, e.g., CERN Courier, 1994, v. 34, N 1, p. 12-13. 2. L.D. Landau , I.Ya. Pomeranchuk , Dokl. AN SSSR 92, 735 (1953). 3. M.L. Ter-Mikayelian , JETP 25, 289 (1953); 25, 296 (1953). 4. I had a pleasure to be his supervisor. 5. See footnote . 6. A.I. Akhiezer , N.F. Shulga , Uspekhi Fiz. Nauk 137, 561 (1982). 7. A colourful discussion with them on this occasion I have described in my memoirs. See: E.L. Feinberg ”Landau et al”, Reminescence of L.D. Landau, Moscow, 1988, p. 253. 8.For details see: M.L. Ter-Mikayelian , ”High Energy Electromagnetic Processes in Medium”, N.Y., 1972. 9. Later this misunderstanding was corrested . See: V.B. Berestetskii , E.M. Lifshits , L.P. Pitaevskii, Quantum Electrodynamics, Moscow, 1980, p. 452. 10. A.B. Migdal , Dokl. AN SSSR 54, (1954); 105, 77 (1955). 11. M.L. Ter-Mikayelian , Dokl. AN SSSR 94, 1033 (1954). 12. By treating this effect the same technique was used as in the work of Landau and Pomeranchuk. Migdal in his final publication of 1955 (see footnote ) introduced also the corresponding changes into his expressions. 13. See Landau , I.Ya. Pomeranchuk , JETP 24, 505 (1953). 15. Pomeranchuk I.Ya, Feinberg E.L., Dokl. AN SSSR 93, 439 (1953). 16. See: P. Lanshof , Pomeron, Priroda, 1994, N 12, p. 17-25. 17. This phenomenon was predicted independently and practically at the same time by different authors. See: A.I. Akhiezer , A.G. Sitenko , JETP 32, 794 (1957); E.L. Feinberg, JETP 29 115 (1955); R. Clauber , Phys. Rev., 88, 30 (19550. 18. See footnote . For the relation between the transition radiation and the bremsstrahlung of ultrasort particles, see: M.L. Ter-Mikayelian, ”Radiation of Particles in Periodic Media”, Priroda, N 12, p. 68-73. 19. E.L. Feinberg , I.Ya Pomeranchuk , Nouvo Cimento, Supplemento, 111 652 (1956); E.L. Feinberg , Usp. Fiz. Nauk, 58 193 (1956). Figure captions Fig.1. DXR spectrum at$`\mathrm{\Theta }=42.9`$ the 6.8 MeV electron beam direction. Fig.2. The spectra of electrons emission in quartz crystal under acoustic waves ($`\mathrm{}`$) and unaffected ($``$). Fig.3. Particle passage through a stack of foils. Fig.4. TR angular distribution in a plate. Thickness of beryllium foil was 1 $`\mu m`$. Experiment demonstrates interference effect in a plate. Fig.5. Integrated from 1 to 10 keV TR and RTR angular distribution from TR (incoherent) ($`l_2`$=1.5 mm) and RTR (coherent) ($`l_2`$=115, 230 and 345 $`\mu m`$) stacks of 8 myler foils ($`l_1=3.8\mu m`$). FIg.6. i) The measured and calculated peak angle for TR (incoherent) and RTR (coherent); ii) The measured and iii) calculated spatial distribution for the (a) RTR, coherent and (b) TR, incoherent in myler stack. Fig.7. Experimental setups (Japan-Russian joint project).
warning/0003/hep-ph0003112.html	ar5iv	text	# Kalb-Ramond axion production in anisotropic string cosmologies ## I Introduction The pre-big-bang (PBB) model of cosmology inspired by the duality properties of string theory, is faced, on the phenomenological side, with the question whether or not it can reproduce the amplitude and slope of the observed temperature anisotropy spectrum and of large-scale density perturbations. Within the PBB model, the inflationary expansion due to the dilaton field in the low-energy effective action of string theory, leads to an amplification of metric fluctuations as well as of quantum fluctuations of all the fields present in PBB cosmology. Such fields, which are not part of the homogeneous background whose perturbations we study, are for example the gauge fields and the pseudo-scalar partner of the dilaton field in the string theory effective action. At first, it was thought that the PBB scenario could not lead to the observed scale-invariant Harrison-Zel’dovich spectrum of perturbations at large-scales. First-order scalar and tensor metric perturbations were found to lead to primordial spectra that grow with frequency , with a normalization imposed by the string cut-off at the shortest amplified scales. These blue spectra have too little power at scales relevant for the observed anisotropies in the cosmic microwave background (CMB). In contrast, the axion energy spectra were found to be diverging at large scales, red spectra, leading to very large CMB anisotropies, in conflict with observations. These results already rule out four-dimensional isotropic PPB cosmology. However, if one allows for internal contracting dimensions in addition to the three expanding ones, the situation is different. The axion field can lead to a flat Harrison-Zel’dovich spectrum of fluctuations for an appropriate relative evolution of the external and the compactified internal dimensions . Thus, it is possible that the amplification of quantum fluctuations of fields which are present in the PBB scenario, can generate via the seed mechanism the observed anisotropy of the CMB radiation. Considering an isotropic PBB model with extra dimensions, the amplification of electromagnetic vacuum fluctuations and of Kalb-Ramond axion vacuum fluctuations lead to interesting observational consequences within the context of primordial magnetic fields and large-scale temperature anisotropies . In particular, massless axions as well as very light axions can exhibit a flat or slightly tilted blue spectrum which may reasonably fit the observational data . (Even though an acoustic peak at $`\mathrm{}350`$ is excluded by experiments published after Ref. was completed, it is possible to shift this peak to $`\mathrm{}220`$ by closing the universe with a cosmological constant. More details about this model can be found in Ref. .) Recently it has been suggested that four-dimensional string cosmology models which expand anisotropically can also lead to blue or flat energy spectra for axionic perturbations . According to Ref. , one can instead of assuming internal extra dimensions , consider an anisotropic four-dimensional background. This has become especially interesting in view of new results which show that the pre-big-bang phase may generically be homogeneous but anisotropic . In Ref. , the axion spectrum is only computed for the part of phase space where the longitudinal component of the wave vector is sufficiently large. In this work we correct the result of Ref. and complete the computation to contain all directions in phase space. We then integrate the obtained spectrum over directions and compare it with the result for the isotropic PBB. We find that the anisotropic spectrum, when averaged over directions agrees roughly with the isotropic one. Therefore, anisotropic expansion during the pre-big-bang phase cannot solve the axion problem of four-dimensional string cosmology. ## II Axion production in the pre-big-bang cosmological model Let us consider a four-dimensional spatially flat anisotropic PBB cosmological model, with metric $$\left(g_{\mu \nu }\right)=\mathrm{diag}[1,a^2(t),b^2(t),b^2(t)];$$ (1) the internal compactified radii (if present) are assumed to be frozen. For simplicity, we assume two directions to expand with the same scale factor $`b`$. Varying the low-energy string theory effective action (in the string frame) $`S=`$ $`{\displaystyle \frac{1}{2\lambda _s^2}}{\displaystyle d^4x\sqrt{g}e^\varphi }`$ (3) $`\times \left[R+g^{\alpha \beta }_\alpha \varphi _\beta \varphi {\displaystyle \frac{1}{12}}H_{\mu \nu \alpha }H^{\mu \nu \alpha }\right]`$ ($`H^{\mu \nu \alpha }`$ denotes the antisymmetric tensor field), with respect to the metric and the dilaton field $`\varphi `$, we obtain the dilaton driven vacuum solutions of the tree-level evolution equations. As derived in Ref. , these solutions read $$a(\eta )=\left[\frac{\eta }{\eta _1}\right]^{\frac{\alpha }{1\alpha }},b(\eta )=\left[\frac{\eta }{\eta _1}\right]^{\frac{\beta }{1\alpha }},$$ (4) and $$\varphi (\eta )=\left(\frac{\alpha +2\beta 1}{1\alpha }\right)\mathrm{log}\left[\frac{\eta }{\eta _1}\right],$$ (5) with $`\alpha `$ and $`\beta `$ satisfying the Kasner condition $$\alpha ^2+2\beta ^2=1.$$ (6) Here $`\eta `$ denotes conformal time with respect to the scale factor $`a`$. It is negative during the pre-big-bang era and $`\eta =\eta _1`$ stands for the transition time from the dilaton driven pre-big-bang era to the radiation dominated post-big-bang era. To obtain the axion evolution equation, we vary the effective action, Eq. (3), with respect to the Kalb-Ramond axion field $`\sigma `$, given by $$H^{\mu \nu \alpha }=e^\varphi \frac{ϵ^{\mu \nu \alpha \rho }}{\sqrt{g}}_\rho \sigma .$$ (7) The evolution equation of the canonical field $`\psi =e^{\varphi /2}b\sigma `$, in Fourier space, reads $$\psi _k^{\prime \prime }+\left[k_L^2+k_T^2\frac{a^2}{b^2}\frac{𝒫^{\prime \prime }}{𝒫}\right]\psi _k=0,𝒫=e^{\varphi /2}b,$$ (8) where $`k_L`$ denotes the modulus of the comoving longitudinal momentum and $`k_T=\sqrt{k_y^2+k_z^2}`$ is the modulus of the transverse momentum. Equation (8) describes the generation of axionic modes, where the anisotropy of the spacetime has been translated into an asymmetry between the longitudinal and transverse momenta. The choice $`\alpha =\beta =1/\sqrt{3}`$ corresponds to the isotropic case, for which $`a=b`$ and the evolution of axionic fluctuations is given by $`\psi _k^{\prime \prime }+\left[k^2{\displaystyle \frac{𝒫^{\prime \prime }}{𝒫}}\right]\psi _k=0,`$ (9) with $`\psi _k=𝒫\sigma _k,𝒫=e^{\varphi /2}a(\eta )^p,`$ (10) so that $`{\displaystyle \frac{𝒫^{\prime \prime }}{𝒫}}={\displaystyle \frac{(\mu ^21/4)}{\eta ^2}},\text{ with }\mu ^2=(p{\displaystyle \frac{1}{2}})^2.`$ (11) The solution of Eq. (9), normalized to an initial vacuum fluctuation spectrum, can be written as $$\psi _k=\eta ^{1/2}H_\mu ^{(2)}(\|k\eta \|),\mu =\left\|p\frac{1}{2}\right\|,\eta \eta _1,$$ (12) with $`\mu =\sqrt{3}`$. $`H_\mu ^{(2)}`$ denotes the Hankel function of second kind (we adopt the conventions of Ref. ). Assuming that the dilaton driven era is followed by a radiation dominated era, the density parameter of produced Kalb-Ramond axions per logarithmic frequency interval is $$\mathrm{\Omega }_\sigma (\omega ,\eta )=\frac{\rho (\omega )}{\rho _c}=\frac{1}{\rho _c}\frac{d\rho _\sigma }{d\mathrm{log}\omega }g_1^2\mathrm{\Omega }_\gamma (\eta )\left[\frac{\omega }{\omega _1}\right]^{32\mu },$$ (13) where $`\rho (\omega )`$ denotes their spectral energy density and $`\rho _c=3M_p^2H^2/(8\pi )`$ stands for the critical energy density. Note that $`\omega _1=k_1/a_1=1/(a_1\|\eta _1\|)`$ represents the maximal amplified frequency, $`g_1=H_1/M_p`$ is the transition scale in units of the Planck mass, $`H_1\omega _1`$ denotes the Hubble scale at which the universe becomes radiation dominated. Hence $`\mathrm{\Omega }_\gamma (\eta )=(H_1/H)^2(a_1/a)^4`$ is the radiation density parameter at a given time $`\eta `$. Clearly a flat spectrum corresponds to $`\mu =3/2`$ and the value $`\mu =\sqrt{3}`$ obtained in a four-dimensional isotropic pre-big-bang model implies a red spectrum, leading to an unacceptable divergence at low frequencies. Let us now go back to the case of a four-dimensional anisotropic background. We first study the evolution of axionic fluctuations and we then calculate the spectral energy density of the axionic inhomogeneities ($`d\rho _\sigma /d\mathrm{log}\omega `$), as they re-enter the horizon during the isotropic radiation dominated era, after being amplified during the anisotropic dilaton driven era. Inserting Eqs. (4) and (5) into Eq. (8), we obtain $$\psi _k^{\prime \prime }+\left(k_L^2+k_T^2\left[\frac{\eta }{\eta _1}\right]^\gamma \frac{\mu ^21/4}{\eta ^2}\right)\psi _k=0,$$ (14) where $`\gamma `$ $`={\displaystyle \frac{2(\alpha \beta )}{1\alpha }},2\mu `$ $`=\|2p1\|,\text{ where}`$ (15) $`p`$ $`={\displaystyle \frac{\alpha +4\beta 1}{2(1\alpha )}},2\mu `$ $`=2{\displaystyle \frac{4\beta }{1\alpha }}.`$ (16) If $`\gamma <0`$, the $`k_T`$-term as well as the $`\eta ^2`$-term go to zero for $`\eta \mathrm{}`$; and initially the parentheses in Eq. (14) is dominated by $`k_L^2`$ (except if $`k_L0`$). If $`k_T`$ is not very large, namely if $$k_T<k_L(k_1/k_L)^{\gamma /2},$$ (17) the scale $`k_L`$ becomes super-horizon, i.e. the parentheses in Eq. (14) is dominated by the $`1/\eta ^2`$-term, before the $`k_T`$-term takes over. In this case, we may entirely neglect the $`k_T`$-term in Eq. (14), which then reduces to Eq. (9) with $`k`$ replaced by $`k_L`$. Therefore, the spectrum for these modes is flat for $`\mu =3/2`$ which corresponds to $$\alpha =7/9,\beta =4/9\text{ and }\gamma =3/8.$$ (18) We also require the solution to expand, i.e. $`\alpha ,\beta <0`$. We first concentrate mainly on these values of the Kasner exponents since they lead to a scale invariant spectrum of fluctuations for directions with a sufficiently large $`k_L`$-component, but we express our results in terms of $`\alpha `$ and $`\beta `$ so that they can then also be applied also to other values of the Kasner indices. In the part of phase-space defined by the inequality given in Eq. (17), the energy density of the produced axions has already been determined in Ref. . Here we correct the result of Ref. and generalize it to the entire phase space. To solve Eq. (14), we distinguish among the following two cases: (I) The modulus of the longitudinal momentum, $`k_L`$, always dominates until $`\eta ^2<1/k_L^2`$ at which point the $`1/\eta ^2`$ term comes to dominate. This is equivalent to the condition given in Eq. (17). (II) At some conformal time $`\eta =\eta _T<\eta _1`$, the modulus of the transverse momentum, $`k_T`$, comes to dominate over $`k_L`$, but the mode is still well within the horizon, i.e. $`\sqrt{k_L^2+k_T^2(\eta _T/\eta _1)^\gamma }>\eta _T^2`$. Equation (14) implies $$\eta _T=\eta _1\left(\frac{k_L}{k_T}\right)^{2/\gamma }.$$ (19) Case (I) : Let us first discuss this case which is also the one studied in Ref. . Here, the inequality given in Eq. (17) holds. For low frequency modes, $`\omega \omega _1`$ this is the case outside a very thin slice around the plane $`k_L=0`$ if $`\gamma <0`$. In this situation we may entirely neglect the second term inside the parentheses of Eq. (14) which yields a Bessel differential equation. Its solution during the pre-big-bang era , is simply $`\psi _k^{\mathrm{PBB}}(k,\eta )=\sqrt{{\displaystyle \frac{\|k_L\eta \|}{k_L}}}`$ $`H_\mu ^{(2)}(k_L\eta ),`$ (21) $`\text{for }\eta \eta _1,`$ After the transition to the radiation dominated era (RD), we assume the dilaton to be frozen and the expansion to have become isotropic. This implies $`𝒫^{\prime \prime }=0`$, $`a/b=1`$ and Eq. (8) reduces to a simple harmonic equation with general solution $`\psi _k^{\mathrm{RD}}(k,\eta )={\displaystyle \frac{1}{\sqrt{k}}}`$ $`\left[c_+e^{ik(\eta +\eta _1)}+c_{}e^{ik(\eta +\eta _1)}\right],`$ (23) $`\text{for }\eta \eta _1.`$ By matching the in-coming solution $`\psi _k^{\mathrm{PBB}}`$ to the out-going one $`\psi _k^{\mathrm{RD}}`$, and by also matching their first derivatives, at the transition time $`\eta =\eta _1`$, we obtain the frequency mixing coefficient $`c_{}(k)`$: $$c_{}=\frac{1}{\sqrt{2\pi }}\sqrt{\frac{1}{(k\eta _1)(k_L\eta _1)^{2\mu }}}.$$ (24) The coefficient $`c_{}`$ determines the occupation numbers of produced axions. The spectral energy density of the produced axions reads $$\rho _L(\omega ,s)=\frac{d\rho _\sigma }{d\mathrm{log}\omega }\frac{\omega ^4}{\pi ^2}\|c_{}(\omega )\|^2.$$ (25) From Eqs. (24), (25) we obtain with $`\omega =k/a`$ for $`\mu =3/2`$ $$\rho _L(\omega ,s)\frac{1}{2\pi ^3}\omega _1^4/s^3,$$ (26) where $`s=k_L/k`$. Thus, if the longitudinal momentum $`k_L`$ dominates, the spectrum of produced Kalb-Ramond axions is flat, i.e. independent of $`\omega `$, but anisotropic. This result generically agrees with the finding of Ref. (up to a factor $`1/s^2`$, which we think is missing in Ref. ). Case (II) : We now assume that the $`k_T`$-term comes to dominate before the perturbation becomes superhorizon. As long as the perturbation is sub-horizon, we may approximate Eq. (14) by $$\psi _k^{\prime \prime }+\left(k_L^2+k_T^2\left[\frac{\eta }{\eta _1}\right]^\gamma \right)\psi _k=0,$$ (27) An approximate solution to this equation is $$\psi \frac{\mathrm{exp}\left(\eta \sqrt{k_L^2+q^2(\eta /\eta _1)^\gamma k_T^2}\right)}{\sqrt{\pi /2}\left[k_L^2+(\eta /\eta _1)^\gamma k_T^2\right]^{1/4}},$$ (28) with $`q=1/(1+\gamma /2)=(1\alpha )/(1\beta )`$. In the regime considered, $`\eta \sqrt{k_L^2+q^2(\eta /\eta _1)^\gamma k_T^2}1`$, this solution becomes exact, if either $`k_L`$ or $`k_T`$ vanishes and it is a good approximation if one of the two terms dominates. If the $`k_L`$-term and the $`k_T`$-term are of the same order, the relative error is about $`\|\gamma /2\|=3/16`$. It is also clear that this represents the correctly normalized incoming vacuum solution. At conformal time $`\eta =\eta _T`$, the transverse momentum $`k_T`$ comes to dominate over the $`k_L`$\- term in Eq. (14). At even later times, the $`\eta ^2`$-term will eventually dominate. After $`\eta _T`$ Eq. (14) can be approximated by $$\psi _k^{\prime \prime }+\left(k_T^2\left[\frac{\eta }{\eta _1}\right]^\gamma \frac{\mu ^21/4}{\eta ^2}\right)\psi _k=0,$$ (29) with general solution $`\psi _k(k_T,\eta )=`$ $`c_T^{(1)}\sqrt{\|k_T\eta \|}H_{\mu q}^{(1)}\left(\|k_T\eta \|q\left[{\displaystyle \frac{\eta }{\eta _1}}\right]^{\gamma /2}\right)`$ (31) $`ic_T^{(2)}\sqrt{\|k_T\eta \|}H_{\mu q}^{(2)}\left(\|k_T\eta \|q\left[{\displaystyle \frac{\eta }{\eta _1}}\right]^{\gamma /2}\right),`$ where $`q`$ is as above, and $`H_{\mu q}^{(1)},H_{\mu q}^{(2)}`$ are Hankel functions of the 1st and 2nd kind of order $`\mu q`$. For large $`k_T\|\eta \|`$ the second term just corresponds to the solution (28) in the limit where $`k_L`$ can be neglected. Therefore, by matching the solutions we find $`c_T^{(1)}`$ $`=`$ $`0`$ (32) $`c_T^{(2)}`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{k_T}}},`$ (33) up to an irrelevant phase. Next, we have to match the solution for the field $`\sigma `$ of the pre-big-bang era to the solution for $`\sigma `$ during the radiation era at the transition time $`\eta =\eta _1`$. As we go from the pre- to the post-big-bang era, we assume the universe to become isotropic and the dilaton field $`\varphi `$ to become frozen. Thus, here the matching of the in-coming to the out-going solution for $`\sigma `$, is not equivalent to matching $`\psi `$. The relation between the the axion field $`\sigma `$ and the canonical field $`\psi `$ at conformal time $`\eta `$ is $`\sigma ^{\mathrm{RD}}(\eta )`$ $`=`$ $`\left[{\displaystyle \frac{\eta }{\eta _1}}\right]^1\psi ^{\mathrm{RD}}(\eta ),`$ (34) $`\sigma ^{\mathrm{PBB}}(\eta )`$ $`=`$ $`\left[{\displaystyle \frac{\eta }{\eta _1}}\right]^\lambda \psi ^{\mathrm{PBB}}(\eta ).`$ (35) The canonical field in Fourier space during RD is given in Eq. (23). Matching the solutions and their first derivatives for $`\sigma `$, as we pass from PBB to RD at time $`\eta =\eta _1`$, we obtain for $`\|k_T\eta _1\|1`$, the Bogoliubov coefficient $`c_{}`$ given by $`\|c_{}\|^2=`$ $`\left[{\displaystyle \frac{\mathrm{\Gamma }^2(\mu q)}{4\pi ^2}}2^{2\mu q}\left({\displaystyle \frac{3}{2}}\mu q\right)^2\right]\left({\displaystyle \frac{k_T}{k_L}}\right)^{2\mu q}`$ (37) $`\times s^{2\mu q}\left({\displaystyle \frac{\omega }{\omega _1}}\right)^{12\mu q}.`$ With Eq. (25) we then obtain that the energy density of the produced Kalb-Ramond axions, in the case where the transverse momentum $`k_T`$ comes to dominate, i.e the inequality given in Eq. (17) is violated: $`\rho _T(\omega ,s)=`$ $`\left[{\displaystyle \frac{\mathrm{\Gamma }^2(\mu q)}{4\pi ^2}}2^{2\mu q}\left({\displaystyle \frac{3}{2}}\mu q\right)^2\right]{\displaystyle \frac{1}{\pi ^2}}`$ (39) $`\times \left({\displaystyle \frac{k}{k_T}}\right)^{2\mu q}\omega _1^{1+2\mu q}\omega ^{32\mu q}.`$ Inserting the values $`\mu =3/2,\alpha =7/9\beta =4/9`$ which lead to a flat spectrum in case I one finds a somewhat blue spectrum in case II, $$\rho _T(\omega ,s)\omega ^{9/13}.$$ (40) Of course this case also gives a finite answer on the plane $`k_L=0`$ for which the result obtained under case I diverges. ## III Results and Conclusion In total we can summarize the calculated spectrum by $$\rho (\omega ,s)\frac{\omega _1^4}{2\pi ^3}\{\begin{array}{c}s^{2\mu }\left(\frac{\omega }{\omega _1}\right)^{32\mu }\text{ if }k_T<k_L(k_1/k_L)^{\gamma /2}\hfill \\ (1s^2)^{\mu q}\left(\frac{\omega }{\omega _1}\right)^{32\mu q}\text{ else,}\hfill \end{array}$$ (41) where $`s=k_L/k`$ , $`\mu =12\beta /(1\alpha )`$ and $`q=(1\alpha )/(1\beta )`$. For our prefered values, $`\alpha =7/9`$ and $`\beta =4/9`$ which imply $`\mu =3/2`$ and $`q=16/13`$, the above result reduces to $$\mathrm{\Omega }_\sigma (\omega ,s,\eta )g_1^2\mathrm{\Omega }_\gamma (\eta )\{\begin{array}{c}s^3\text{ if }k_T<k_L(k_1/k_L)^{\gamma /2}\hfill \\ (1s^2)^{\frac{24}{13}}\left(\frac{\omega }{\omega _1}\right)^{9/13}\text{ else.}\hfill \end{array}$$ (42) In the regime of phase space where the longitudinal mode of the momentum is very small i.e. when the condition given in Eq. (17) is violated, the spectrum of the produced Kalb-Ramond axions is not flat. For a given value of $`\omega `$, this is the case if $`s`$ is smaller than the critical value $`s_c`$ which is well approximated by $$s_c(\omega )\left(\frac{\omega }{\omega _1}\right)^{3/13},\text{ if }\omega \stackrel{<}{}0.1\omega _1,$$ (43) a very small value for cosmologically interesting frequencies. In Figure 1 the the energy density $`\rho (\omega )`$ is shown as a function of $`s`$ for different values of $`\omega `$. For $`s`$ fixed, if the modulus of the longitudinal momentum dominates in Eq. (14), more precisely if it satisfies the condition $`k_L>k_T(k_T/k_1)^{\gamma /2/(1+\gamma /2)}`$, the spectrum of the produced Kalb-Ramond axions is flat. To estimate the total energy density per logarithmic frequency interval we integrate the axion density $`\mathrm{\Omega }_\sigma (\omega ,s)`$ over $`s`$. For this we use $$d^3k=2\pi k_Tdk_Ldk_T=4\pi k^2dsdk,$$ (44) where we have used $`dk_L=kds+sdk`$ and $`dk_T={\displaystyle \frac{s}{\sqrt{1s^2}}}kds+\sqrt{1s^2}dk.`$ Hence, we have $`\mathrm{\Omega }_\sigma (\omega ,\eta )={\displaystyle \mathrm{\Omega }_\sigma (\omega ,s,\eta )𝑑s}`$ (45) $``$ $`{\displaystyle \frac{1}{\rho _c}}\left[{\displaystyle _0^{s_c(\omega )}}\rho _T(\omega ,s)𝑑s+{\displaystyle _{s_c(\omega )}^1}\rho _L(\omega ,s)𝑑s\right]`$ (46) $``$ $`g_1^2\mathrm{\Omega }_\gamma (\eta )\left[s_c\left({\displaystyle \frac{\omega _1}{\omega }}\right)^{9/13}+{\displaystyle \frac{0.5}{s_c^2}}\right].`$ (47) This spectrum is shown in Fig. 2. Using $`s_c(\omega /\omega _1)^{3/13}`$, which is a good approximation as long as $`\omega 0.1\omega _1`$ it can also be seen directly from Eq. (47) that the isotropic energy spectrum is nearly reproduced. The isotropic spectral index, $`32\sqrt{3}0.464`$ is actually replaced by $`6/130.463`$. Inserting reasonable values for the string scale, $`0.01g_1<1`$, we see that also in the anisotropic case axions are over-produced in unacceptable amounts. Even if the spectrum of the axions from wave vectors directed sufficiently far from the plane $`k_L=0`$, is scale-invariant, the enhancement of the spectrum in the vicinity of the plane $`k_L=0`$ leads to a total contribution which agrees with the one obtained in the isotropic case. Therefore, the model is excluded (see Ref. ). So far we have mainly considered the case $`\alpha =7/9`$ and $`\beta =4/9`$, but our results apply quite generically, as long as $`\gamma <0`$ and thus the $`k_L`$-term dominates at sufficiently early times. But also if $`\gamma >0`$, Eq. (28) is an approximate solution on sub-horizon scales. In this situation, however the $`k_T`$-term dominates at sufficiently early times and continues to do so until the perturbation becomes super-horizon if the inequality given in Eq. (17) is violated. For $`\gamma >0`$ this is the case outside a narrow cylinder around the $`k_T=0`$ axis. Therefore, the generic formula given in Eq. (46) always applies, but $`s_c1`$, if $`\gamma <0`$ and $`s_c1`$, if $`\gamma >0`$. For general values of $`\alpha `$ and $`\beta `$ we obtain $`\mathrm{\Omega }_\sigma (\omega ,\eta )={\displaystyle \mathrm{\Omega }_\sigma (\omega ,s,\eta )𝑑s}`$ (48) $``$ $`g_1^2\mathrm{\Omega }_\gamma (\eta )[\left({\displaystyle \frac{\omega }{\omega _1}}\right)^{32\mu q}{\displaystyle _0^{s_c(\omega )}}(1s^2)^{\mu q}ds`$ (50) $`+\left({\displaystyle \frac{\omega }{\omega _1}}\right)^{32\mu }{\displaystyle _{s_c(\omega )}^1}s^{2\mu }ds].`$ The transition value of $`s`$ is given by $$\sqrt{1s_c^2}=s_c^{1+\gamma /2}\left(\frac{\omega }{\omega _1}\right)^{\gamma /2}.$$ (51) If $`\gamma <0`$ (i.e. $`\alpha <\beta `$), the factor $`\left(\frac{\omega }{\omega _1}\right)^{\gamma /2}`$ is very large in most of phase space and hence $`s_c1`$. On the other hand, if $`\gamma >0`$ (i.e. $`\alpha >\beta `$), the above factor is very small for the relevant frequencies, $`\omega \omega _1`$ and $`s_c1`$. A reasonable approximation is $`s_c`$ $`\left({\displaystyle \frac{\omega }{\omega _1}}\right)^{q1}`$ $`\text{if }\gamma <0`$ (52) $`1s_c^2`$ $`\left({\displaystyle \frac{\omega }{\omega _1}}\right)^{2/q2}`$ $`\text{if }\gamma >0,`$ (53) where we have used the relation $`q=1/(1+\gamma /2)`$. Inserting these results in Eq. (50), the integrals can be approximated by $$\mathrm{\Omega }_\sigma (\omega ,\eta )g_1^2\mathrm{\Omega }_\gamma (\eta )\left(\frac{\omega }{\omega _1}\right)^n,\text{ where }$$ (54) $`n`$ $`=2+q2\mu q`$ $`={\displaystyle \frac{1+\alpha +2\beta }{1\beta }}\text{ if }\alpha <\beta ,`$ (55) $`n`$ $`=1+2/q2\mu `$ $`={\displaystyle \frac{1+\alpha +2\beta }{1\alpha }}\text{ if }\alpha >\beta .`$ (56) Clearly, since $`\alpha ^2+2\beta ^2=1`$ and $`\alpha ,\beta 0`$ it is $`\alpha +2\beta 1`$. This shows that the spectrum is never blue and becomes scale invariant only in the degenerate case with two static dimensions, $`\beta =0`$. This is also shown in Fig. 3, where the above approximation for the spectral index plotted as a function of $`\alpha `$: the spectrum always remains red with a spectral index relatively close to the isotropic value, $`n_{\mathrm{iso}}=32\sqrt{3}0.46`$, except in the extremal case, when two dimensions are frozen and $`\alpha =1`$. If one relaxes the condition that both $`a`$ and $`b`$ be expanding and just asks for volume expansion, $`\alpha +2\beta <0`$, there is another pair of values for the Kasner indices leading to a flat spectrum, namely $`\alpha =1/3`$ and $`\beta =2/3`$. However, if we want expansion in all three dimensions the spectrum is always red. To summarize, we find that anisotropic expansion has very little influence on the overall axion production and cannot cure the axion problem of four-dimensional pre-big-bang models. Only by allowing for extra dimensions one can escape this conclusion and obtain a scale invariant spectrum of axions as described in Refs. . A ’realistic’ string cosmology with a Kalb-Ramond axion can therefore be realized only in models with extra dimensions. ###### Acknowledgements. It is a pleasure to thank A. Buonanno, T. Damour, K. Kunze, G. Veneziano and A. Vilenkin for useful discussions. This work is supported by the Swiss National Science Foundation.
warning/0003/cond-mat0003355.html	ar5iv	text	# Calculation of Second Topological Moment ⟨𝑚²⟩ of Two Entangled Polymers ## I The Problem Consider two polymers $`P_1`$ and $`P_2`$ which statistically can be linked with each other any number of times $`m=0,1,2,\mathrm{}`$. The situation is illustrated in Fig. 1 for $`m=2`$. We would like to find the probability distribution of the linking numbers $`m`$ as a function of the lengths of $`P_1`$ and $`P_2`$. As a first contribution to solving this problem we calculate, in this note, the second moment of this distribution, $`m^2`$. The self-entanglements are ignored. The solution of this two-polymer problem may be considered as an approximation to a more interesting physical problem in which a particular polymer is linked to any number $`N`$ of polymers, which are effectively replaced by a single long “effective” polymer . Let $`G_m(𝐱_1,𝐱_2;L_1,L_2)`$ be the configurational probability to find the polymer $`P_1`$ of length $`L_1`$ with fixed coinciding end points at $`𝐱_1`$ and the polymer $`P_2`$ of length $`L_2`$ with fixed coinciding end points at $`𝐱_2`$, entangled with a Gaussian linking number $`m`$. The second moment $`m^2`$ is defined by the ratio of integrals $`m^2={\displaystyle \frac{d^3𝐱_1d^3𝐱_2_{\mathrm{}}^+\mathrm{}𝑑mm^2G_m(𝐱_1,𝐱_2;L_1,L_2)}{d^3𝐱_1d^3𝐱_2_{\mathrm{}}^+\mathrm{}𝑑mG_m(𝐱_1,𝐱_2;L_1,L_2)}}`$ (1) performed for either of the two probabilities. The integrations in $`d^3𝐱_1d^3𝐱_2`$ covers all positions of the end points. The denominator plays the role of a partition function of the system: $`Z{\displaystyle d^3𝐱_1d^3𝐱_2_{\mathrm{}}^+\mathrm{}𝑑mG_m(𝐱_1,𝐱_2;L_1,L_2)}`$ (2) Due to the translational invariance of the system, the probabilities depend only on the differences between the end point coordinates: $`G_m(𝐱_1,𝐱_2;L_1,L_2)=G_m(𝐱_1𝐱_2;L_1,L_2)`$ (3) Thus, after the shift of variables, the spatial double integrals in (1) can be rewritten as $`{\displaystyle d^3𝐱_1d^3𝐱_2G_m(𝐱_1𝐱_2;L_1,L_2)}`$ (4) $`=V{\displaystyle d^3𝐱G_m(𝐱;L_1,L_2)},`$ (5) where $`V`$ denotes the total volume of the system. ## II Polymer Field Theory for Probabilities The linking number for the two polymers is given by the Gauss integral $`I_G(P_1,P_2)={\displaystyle \frac{1}{4\pi }}{\displaystyle _{P_1}}{\displaystyle _{P_2}}[d𝐱_1\times d𝐱_2]{\displaystyle \frac{𝐱_1𝐱_2}{\|𝐱_1𝐱_2\|^3}}.`$ (6) It takes the values $`m=0,\pm 1,\pm 2,\mathrm{}`$. With the help of two vector potentials $`𝐀_1`$ and $`𝐀_2`$, the phase factor $`e^{im\lambda }`$ can be obtained as a result of a local gauge theory of the Chern-Simons type: $`e^{im\lambda }`$ $`=`$ $`{\displaystyle 𝒟A_1^\mu 𝒟A_2^\mu }`$ (7) $`\times `$ $`e^{𝒜_{\mathrm{CS}}\kappa _{P_1}𝑑𝐱_1𝐀_1\lambda _{P_2}𝑑𝐱_2𝐀_2},`$ (8) where $`𝒜_{\mathrm{CS}}`$ is the action $`𝒜_{\mathrm{CS}}=i\kappa {\displaystyle d^3𝐱\epsilon _{\mu \nu \rho }A_1^\mu _\nu A_2^\rho },`$ (9) Indeed, the correlation functions $`D_{ij}^{\mu \nu }(𝐱,𝐱^{})`$ of the gauge fields are $`A_1^\mu (𝐱)A_1^\nu (𝐱^{})`$ $`=`$ $`0,A_2^\mu (𝐱)A_2^\nu (𝐱^{})=0,`$ (10) $`A_1^\mu (𝐱)A_2^\nu (𝐱^{})`$ $`=`$ $`{\displaystyle \frac{d^3p}{(2\pi )^3}e^{i𝐩(𝐱𝐱^{})}\frac{iϵ_{\mu \lambda \nu }k^\lambda }{𝐤^2}}`$ (11) $`=`$ $`{\displaystyle \frac{1}{4\pi }}ϵ_{\mu \lambda \nu }_\lambda {\displaystyle \frac{1}{\|𝐱𝐱^{}\|}}`$ (12) $`=`$ $`{\displaystyle \frac{1}{4\pi }}ϵ_{\mu \nu \kappa }{\displaystyle \frac{(xx^{})^\lambda }{\|𝐱𝐱^{}\|^3}},`$ (13) such that the functional integral on the right-hand side of (8) produces directly the phase factor $`e^{iI_G(P_1,P_2)\lambda }`$ with the eigenvalue $`e^{im\lambda }`$. We can select configurations with a certain linking number $`m`$ from all configurations by forming the integral $`_{\mathrm{}}^{\mathrm{}}𝑑\lambda e^{im\lambda }`$ over this quantity. The most efficient way of describing the statistical fluctuations of the polymers $`P_1`$ and $`P_2`$ is by two complex polymer fields $`\psi _1^{a_1}(𝐱_1)`$ and $`\psi _2^{a_2}(𝐱_2)`$ with $`n_1`$ and $`n_2`$ replica $`(a_1=1,\mathrm{},n_1,a_2=1,\mathrm{},n_2)`$. At the end we shall take $`n_1,n_20`$ to ensure that these fields describe only one polymer each . For these fields we define an auxiliary probability $`G_\lambda (\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z})`$ to find the polymer $`P_1`$ with open ends at $`𝐱_1,𝐱_1^{}`$ and the polymer $`P_2`$ with open ends at $`𝐱_2,𝐱_2^{}`$. The double vectors $`\stackrel{}{𝐱}_1(𝐱_1,𝐱_1^{})`$ and $`\stackrel{}{𝐱}_2(𝐱_2,𝐱_2^{})`$ collect initial and final endpoints of the two polymers $`P_1`$ and $`P_2`$. The auxiliary probability $`G_\lambda (\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z})`$ is given by a functional integral $`G_\lambda (\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z})=\underset{n_1,n_20}{lim}{\displaystyle 𝒟(\mathrm{fields})}`$ (14) $`\times \psi _1^{a_i}(𝐱_1)\psi _1^{a_1}(𝐱_1^{})\psi _2^{a_2}(𝐱_2)\psi _2^{a_2}(𝐱_2^{})e^𝒜`$ (15) where $`𝒟(\text{fields})`$ indicates the measure of functional integration, and $`𝒜`$ the action governing the fluctuations. It consists of kinetic terms for the fields, a quartic interaction of the fields to account for the fact that two monomers of the polymers cannot occupy the same point, the so-called excluded-volume effect, and a Chern-Simons field to describe the linking number $`m`$. Neglecting at first the excluded-volume effect and focusing attention on the linking problem only, the action reads $$𝒜=𝒜_{\mathrm{CS}}+𝒜_{\mathrm{pol}},$$ (16) with a polymer field action $`𝒜_{\mathrm{pol}}={\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle d^3𝐱\left[\|\overline{𝐃}^i\mathrm{\Psi }_i\|^2+m_i^2\|\mathrm{\Psi }_i\|^2\right]}.`$ (17) in which we have omitted a gauge fixing term, which enforces the Lorentz gauge. They are coupled to the polymer fields by the covariant derivatives $`𝐃^i=\mathbf{}+i\gamma _i𝐀^i,`$ (18) with the coupling constants $`\gamma _{1,2}`$ given by $`\gamma _1=\kappa \gamma _2=\lambda .`$ (19) The square masses of the polymer fields contain the chemical potentials $`z_{1,2}`$ of the polymers: $`m_i^2=2Mz_i.`$ (20) They are conjugate variables to the length parameters $`L_1`$ and $`L_2`$, respectively. The symbols $`\mathrm{\Psi }_i`$ collect the replica of the two polymer fields $`\mathrm{\Psi }_i=(\psi _i^1,\mathrm{},\psi _i^{n_i}),\mathrm{\Psi }_i^{}=(\psi _i^1,\mathrm{},\psi _i^{n_i}),`$ (21) and their absolute squares contain the sums over the replica $`\|𝐃^i\overline{\mathrm{\Psi }}_i\|^2={\displaystyle \underset{a_i=1}{\overset{n_i}{}}}\|𝐃^i\psi _i^{a_i}\|^2,\|\mathrm{\Psi }_i\|^2={\displaystyle \underset{a_i=1}{\overset{n_i}{}}}\|\psi _i^{a_i}\|^2.`$ (22) Having specified the fields, we can now write down the measure of functional integration in Eq. (15): $`𝒟(\text{fields})={\displaystyle 𝒟A_1^\mu 𝒟A_2^\mu 𝒟\mathrm{\Psi }_1𝒟\mathrm{\Psi }_1^{}𝒟\mathrm{\Psi }_2𝒟\mathrm{\Psi }_2^{}}.`$ (23) By Eq. (8), the parameter $`\lambda `$ is conjugate to the linking number $`m`$. We can therefore calculate the desired probability $`G_m(\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;L_1,L_2)`$ from the auxiliary one $`G_\lambda (\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z})`$ by the following Laplace integrals $`\stackrel{}{z}=(z_1,z_1)`$: $`G_m(\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;L_1,L_2)=\underset{\genfrac{}{}{0pt}{}{𝐱_1^{}𝐱_1}{𝐱_2^{}𝐱_2}}{lim}{\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}{\displaystyle \frac{dz_1}{2\pi i}}{\displaystyle \frac{dz_2}{2\pi i}}e^{(z_1L_1+z_2L_2)}`$ (24) $`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\lambda e^{im\lambda }G_\lambda (\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z}).`$ (25) ## III Calculating the Partition Function Let us use the polymer field theory to calculate the partition function (2). By Eq. (25), it is given by the integral over the auxiliary probabilities $`Z`$ $`=`$ $`{\displaystyle d^3𝐱_1d^3𝐱_2\underset{\genfrac{}{}{0pt}{}{𝐱_1^{}𝐱_1}{𝐱_2^{}𝐱_2}}{lim}_{ci\mathrm{}}^{c+\mathrm{}}\frac{dz_1}{2\pi i}\frac{dz_2}{2\pi i}e^{(z_1L_1+z_2L_2)}}`$ (27) $`\times {\displaystyle _{\mathrm{}}^+\mathrm{}}dm{\displaystyle _{\mathrm{}}^+\mathrm{}}d\lambda e^{im\lambda }G_\lambda (\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z}).`$ The integration over $`dm`$ is trivial and gives $`2\pi \delta (\lambda )`$, enforcing $`\lambda =0`$, so that $`Z`$ $`=`$ $`{\displaystyle d^3𝐱_1d^3𝐱_2\underset{\genfrac{}{}{0pt}{}{𝐱_1𝐱_1}{𝐱_2^{}𝐱_2}}{lim}_{ci\mathrm{}}^{c+i\mathrm{}}\frac{dz_1dz_2}{2\pi i}e^{(z_1L_1+z_2L_2)}}`$ (29) $`\times G_{\lambda =0}(\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z})`$ To compute $`G_{\lambda =0}(\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z})`$ we observe that the action $`𝒜`$ in Eq. (16) is quadratic in $`\lambda `$. Let us expand $`𝒜`$ as $`𝒜=𝒜_0+\lambda 𝒜_1+\lambda ^2𝒜_2`$ (30) where $`𝒜_0`$ $``$ $`𝒜_{\mathrm{CS}}`$ (31) $`+`$ $`{\displaystyle d^3𝐱\left[\|𝐃_1\mathrm{\Psi }_1\|^2+\|\mathbf{}\mathrm{\Psi }_2\|^2+\underset{i=1}{\overset{2}{}}\|\mathrm{\Psi }_i\|^2\right]},`$ (32) a linear coefficient $`𝒜_1{\displaystyle d^3𝐱𝐣_2(𝐱)𝐀_2(𝐱)}`$ (33) with a “current” of the second polymer field $`𝐣_2(𝐱)=i\mathrm{\Psi }_2^{}(𝐱)\mathbf{}\mathrm{\Psi }_2(𝐱),`$ (34) and a quadratic coefficient $`𝒜_2{\displaystyle \frac{1}{4}}{\displaystyle d^3𝐱𝐀_2^2\|\mathrm{\Psi }_2(𝐱)\|^2}.`$ (35) With these definitions we can rewrite (30) as $`G_{\lambda =0}(\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z})={\displaystyle 𝒟(\mathrm{fields})e^{𝒜_0}}`$ (36) $`\times \psi _1^{a_1}(𝐱_1)\psi _1^{a_1}(𝐱_1^{})\psi _2^{a_2}(𝐱_2)\psi _2^{a_2}(𝐱^{})`$ (37) From Eq. (32) it is clear that $`G_{\lambda =0}(\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z})`$ is the product of the configurational probabilites of two free polymers. Note that the fields $`\mathrm{\Psi }_2,\mathrm{\Psi }_2^{}`$ are free, whereas the fields $`\mathrm{\Psi }_1,\mathrm{\Psi }_1^{}`$ are apparently not free because of the couplings with the Chern-Simons fields through the covariant derivative $`𝐃^1`$. This is, however, an illusion: the fields $`A_\mu ^i`$ have a vanishing diagonal propagators $`A_\mu ^iA_\nu ^i=0`$. integrating out $`A_2^\mu `$ in (37), we find the flatness condition: $$\epsilon ^{\mu \nu \rho }_\nu A_\mu ^i=0.$$ (38) On a flat space with vanishing boundary conditions at infinity this implies $`A_1^\mu =0`$. As a consequence, the functional integral (37) factorizes $`G_{\lambda =0}(\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z})=G_0(𝐱_1𝐱_1^{};z_1)G_0(𝐱_2𝐱_2^{};z_2),`$ (39) where $`G_0(𝐱_i𝐱_i^{};z_i)`$ are the free correlation functions of the polymer fields: $`G_0(𝐱_i𝐱_i^{};z_i)=\psi _i^{a_i}(𝐱_i)\psi _i^{a_i}(𝐱_i^{}).`$ (40) In momentum space, the correlation functions are $`\stackrel{~}{\psi }^{a_i}(𝐤_i)\stackrel{~}{\psi }_i^{a_i}(𝐤_i^{})=\delta ^{(3)}\left(𝐤_i+𝐤_i^{}\right){\displaystyle \frac{1}{𝐤_i^2+m_i^2}}`$ (41) such that $`G_0(𝐱_i𝐱_i^{};z_i)={\displaystyle \frac{d^3k}{(2\pi )^3}e^{i𝐤𝐱}\frac{1}{𝐤_i^2+m_i^2}},`$ (42) and $`G_0(𝐱_i𝐱_i^{};L_i)`$ $`=`$ $`{\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}{\displaystyle \frac{dz_i}{2\pi i}}e^{z_iL_i}G_0(𝐱_i𝐱_i^{};z_i)`$ (44) $`={\displaystyle \frac{1}{4\sqrt{2}M}}\left({\displaystyle \frac{M}{2\pi }}\right)^{3/2}L_i^{3/2}e^{M(𝐱_i𝐱_i^{})/2L_i}.`$ Thus we obtain for (29): $`Z=2\pi {\displaystyle d^3𝐱_1d^3𝐱_2}`$ (46) $`\times \underset{\genfrac{}{}{0pt}{}{𝐱_1^{}𝐱_1}{𝐱_2^{}𝐱_2}}{lim}G_0(𝐱_1𝐱_1^{};L_1)G_0(𝐱_2𝐱_2^{};L_2)`$ The integrals at coinciding end points can easily be performed, and we find $`Z={\displaystyle \frac{2\pi MV^2}{(8\pi )^3}}(L_1L_2)^{3/2}`$ (47) It is important to realize that in Eq. (27), the limits of coinciding end points $`𝐱_i^{}𝐱_i`$ and the inverse Laplace transformations do not commute unless a proper renormalization scheme is chosen to eliminate the divergences caused by the insertion of the composite operators $`\|\psi (r)\|^2`$. This can be seen for a single polymer $`P`$. If we were to commuting the limit of coinciding end points with the Laplace transform, we would obtain $`{\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}{\displaystyle \frac{dz}{2\pi }}e^{zL}\underset{𝐱^{}𝐱}{lim}G_0(𝐱𝐱^{};z)`$ (48) $`=`$ $`{\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}{\displaystyle \frac{dz}{2\pi i}}e^{zL}G_0(\mathrm{𝟎},z),`$ (49) where $`G_0(\mathrm{𝟎};z)=\|\psi (𝐱)\|^2.`$ (50) This expectation value, however, is linearly divergent: $`\|\psi (𝐱_a)\|^2={\displaystyle \frac{d^3k}{k^2+m^2}}\mathrm{}`$ (51) ## IV Calculation of Numerator in Second Moment Let us now turn to the numerator in Eq. (1): $`N{\displaystyle d^2𝐱_1d^3𝐱_2_{\mathrm{}}^{\mathrm{}}𝑑mm^2G_m(𝐱_1,𝐱_2;L_1,L_2)}.`$ (52) We set up a functional integral for $`N`$ in terms of the auxiliary probability $`G_{\lambda =0}(\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z})`$ analogous to Eq. (27): $`N`$ $`=`$ $`{\displaystyle d^3𝐱_1d^3r_2_{\mathrm{}}^{\mathrm{}}𝑑mm^2\underset{\genfrac{}{}{0pt}{}{𝐱_1^{}𝐱}{𝐱^{}𝐱_2}}{lim}_{ci\mathrm{}}^{c_\tau i\mathrm{}}\frac{dz_i}{2\pi i}\frac{dz_2}{2\pi i}}`$ (54) $`e^{(z_1L_1+z_2L_2)}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\lambda e^{im\lambda }G_\lambda (\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z}).`$ The integration in $`dm`$ is easily performed after noting that $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑mm^2e^{im\lambda }G_\lambda (\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z})`$ (56) $`={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑m\left({\displaystyle \frac{^2}{\lambda ^2}}e^{im\lambda }\right)G_\lambda (\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z}).`$ After a double integration by parts in $`\lambda `$, and an integration in $`m`$, we obtain $`N`$ $`=`$ $`{\displaystyle d^3𝐱_1d^3𝐱_2\underset{\genfrac{}{}{0pt}{}{𝐱_1^{}𝐱_1}{𝐱_2^{}𝐱_2}}{lim}(1)_{ci\mathrm{}}^{c+i\mathrm{}}\frac{dz_1}{2\pi i}\frac{dz_2}{2\pi i}e^{(z_1L_1+z_2L_2)}}`$ (58) $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\lambda \delta (\lambda )\left[{\displaystyle \frac{^2}{\lambda ^2}}G_\lambda (\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z})\right]`$ Performing the now trivial integration in $`d\lambda `$ yields $`N`$ $`=`$ $`{\displaystyle d^3𝐱_1d^3𝐱_2\underset{\genfrac{}{}{0pt}{}{𝐱_1^{}𝐱_1}{𝐱_2^{}𝐱_2}}{lim}(1)_{ci\mathrm{}}^{c+i\mathrm{}}\frac{dz_1}{2\pi i}\frac{dz_2}{2\pi i}e^{(z_1L_1+z_2L_2)}}`$ (60) $`\times \left[{\displaystyle \frac{^2}{\lambda ^2}}G_\lambda (\stackrel{}{𝐱}_1,\stackrel{}{𝐱}_2;\stackrel{}{z})\right]_{\lambda =0}`$ To compute the term in brackets, we use again Eq. (30) and Eqs. (32) –(68), and find $`N`$ $`=`$ $`{\displaystyle d^3𝐱_1d^3𝐱_2\underset{\genfrac{}{}{0pt}{}{n_10}{n_20}}{lim}_{ci\mathrm{}}^{c+i\mathrm{}}\frac{dz_1}{2\pi i}\frac{dz_2}{2\pi i}e^{(z_1L_1+z_2L_2)}}`$ (63) $`\times {\displaystyle }𝒟(\text{fields})\mathrm{exp}(𝒜_0)\|\psi _1^{a_1}(𝐱_1)\|^2\|\psi _2^{a_2}(𝐱_2)\|^2`$ $`\times \left[\left({\displaystyle d^3𝐱𝐀_2\mathrm{\Psi }_2^{}\mathbf{}\mathrm{\Psi }_2}\right)^2+{\displaystyle \frac{1}{2}}{\displaystyle d^3𝐱𝐀_2^2\|\mathrm{\Psi }_2\|^2}\right].`$ In this equation we have taken the limits of coinciding endpoint inside the Laplace integral over $`z_1,z_2`$. This will be justified later on the grounds that the potentially dangerous Feynman diagrams containing the insertions of operations like $`\|\mathrm{\Psi }_i\|^2`$ vanish in the limit $`n_1,n_20`$. In order to calculate (63), we decompose the action into a free part $`𝒜_0^0`$ $``$ $`𝒜_{\mathrm{CS}}`$ (64) $`+`$ $`{\displaystyle d^3𝐱\left[\|𝐃^1\mathrm{\Psi }_1\|^2+\|\mathbf{}\mathrm{\Psi }_2\|^2+\underset{i=1}{\overset{2}{}}2\|\mathrm{\Psi }_i\|^2\right]},`$ (65) and interacting parts $`𝒜_1^0{\displaystyle d^3𝐱𝐣_1(𝐱)𝐀_1(𝐱)}`$ (66) with a “current” of the first polymer field $`𝐣_1(𝐱)i\mathrm{\Psi }_1^{}(𝐱)\mathbf{}\mathrm{\Psi }_1(𝐱),`$ (67) and $`𝒜_0^2{\displaystyle \frac{1}{4}}{\displaystyle d^3𝐱𝐀_1^2\|\mathrm{\Psi }_1(𝐱)\|^2}.`$ (68) Expanding the exponential $$e^{𝒜_0}=e^{𝒜_0^0+𝒜_0^1+𝒜_0^2}=e^{𝒜_0}\left[1𝒜_0^1+\frac{(𝒜_0^1)^2}{2}𝒜_0^2+\mathrm{}\right],$$ (69) and keeping only the relevant terms, the functional integral (63) can be rewritten as a purely Gaussian expectation value $`N`$ $`=`$ $`\kappa ^2{\displaystyle d^3𝐱_1d^3𝐱_2\underset{\genfrac{}{}{0pt}{}{n_10}{n_20}}{lim}_{ci\mathrm{}}^{c+i\mathrm{}}\frac{dz_1}{2\pi i}\frac{dz_2}{2\pi i}e^{(z_1L_1+z_2L_2)}}`$ (70) $`\times `$ $`{\displaystyle 𝒟(\text{fields})\mathrm{exp}(𝒜_0^0)\|\psi _1^{a_1}(𝐱_1)\|^2\|\psi _2^{a_2}(𝐱_2)\|^2}`$ (71) $`\times `$ $`\left[\left({\displaystyle d^3𝐱𝐀_1\mathrm{\Psi }_1^{}\mathbf{}\mathrm{\Psi }_1}\right)^2+{\displaystyle \frac{1}{2}}{\displaystyle d^3𝐱𝐀_1^2\|\mathrm{\Psi }_1\|^2}\right]`$ (72) $`\times `$ $`\left[\left({\displaystyle d^3𝐱𝐀_2\mathrm{\Psi }_2^{}\mathbf{}\mathrm{\Psi }_2}\right)^2+{\displaystyle \frac{1}{2}}{\displaystyle d^3𝐱𝐀_2^2\|\mathrm{\Psi }_2\|^2}\right]`$ (73) Only four diagrams shown in Fig. (2) contribute in Eq. (73). Note that the initially asymmetric treatment of polymers $`P_1`$ und $`P_2`$ in the action (16) has led to a completely symmetric expression for the second moment. Only the first diagram in Fig. 2 is divergent due to the divergence of the loop formed by two vector correlation functions. This infinity may be absorbed in the four-$`\mathrm{\Psi }`$ interaction accounting for the excluded volume effect which we do not consider at the moment. No divergence arises from the insertiona of the composite fields $`\|\mathrm{\Psi }_i(𝐱_i)\|^2`$. In this respect, the disconnected diagrams shown in Fig. 3 are potentally dangerous. But these vanish in the limit of zero replica indices $`n_1,n_20`$. ## V Calculation of first Feynman diagrams in Fig. LABEL:@fff From Eq. (73) one has to evaluate the following integral $`N_1`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{4}}\underset{\genfrac{}{}{0pt}{}{n_10}{n_20}}{lim}{\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}{\displaystyle \frac{dz_1}{2\pi i}}{\displaystyle \frac{dz_2}{2\pi i}}e^{(z_1L_1+z_2L_2)}`$ (76) $`{\displaystyle d^3x_1d^3x_2d^3x_1^{}d^3x_2^{}}`$ $`\|\psi _1^{a_1}(𝐱_1)\|^2\|\psi _2^{a_2}(𝐱_2)\|^2\left(\|\mathrm{\Psi }_1\|^2𝐀_1^2\right)_{𝐱_1^{}}\left(\|\mathrm{\Psi }_2\|^2𝐀_2^2\right)_{𝐱_2^{}}.`$ There is an ultraviolet-divergent contribution which should be properly regularized. The system has, of course, a microscopic scale, which is the size of the momomers. This, however, is not the appropriate short-distance scale to be uses here. The model treats the polymers as random chains. However, the momomers of a polymer in the laboratory are usually not freely movable, so that polymers have a certain stiffness. This gives rise to a certain persistence length $`\xi _0`$ over which a polymer is stiff. This length scale is increased to $`\xi >\xi _0`$ by the excluded-volume effects. This is the length scale which should be used as a proper physical short-distance cutoff. We may impose this cutoff by imagining the model as being defined on a simple cubic lattice of spacing $`\xi `$. This would, of course, make analytical calculations quite difficult. Still, as we shall see, it is possible to estimate the dependence of the integral $`N_1`$ and the others in the physically relevant limit in which the lengths of the polymers are much larger than the persistence length $`\xi `$. An alternative and simpler regularization is based on cutting off all ultraviolet-divergent continuum integrals at distances smaller than $`\xi `$. After such a regularization, the calculation of $`N_1`$ is rather straightforward. Replacing the expectation values by the Wick contractions corresponding to the first diagram in Fig. 2, and performing the integrals as shown in the Appendix, we obtain $`N_1={\displaystyle \frac{V}{4\pi }}{\displaystyle \frac{M^2}{(8\pi )^6}}(L_1L_2)^{\frac{1}{2}}`$ (77) $`\times `$ $`{\displaystyle _0^1}𝑑s\left[(1s)s\right]^{\frac{3}{2}}{\displaystyle d^3xe^{M𝐱^2/2s(1s)}}`$ (78) $`\times `$ $`{\displaystyle _0^1}𝑑t\left[(1t)t\right]^{\frac{3}{2}}{\displaystyle d^3ye^{M𝐲^2/2t(1t)}d^3x_1^{\prime \prime }\frac{1}{\|𝐱_1^{\prime \prime }\|^4}}.`$ (79) The variables $`𝐱`$ and $`𝐲`$ have been rescaled with respect to the original ones in order to extract the behavior of $`N_1`$ in $`L_1`$ and $`L_2`$. As a consequence, the lattices where $`𝐱`$ and $`𝐲`$ are defined have now spacings $`\xi /\sqrt{L_1}`$ and $`\xi /\sqrt{L_2}`$ respectively. The $`𝐱,𝐲`$ integrals may be explicitly computed in the physical limit $`L_1,L_2>>\xi `$, in which the above spacings become small. Moreover, it is possible to approximate the integral in $`𝐱_1^{\prime \prime }`$ with an integral over a continuous variable $`\rho `$ and a cutoff in the ultraviolet region: $`{\displaystyle d^3x_1^{\prime \prime }\frac{1}{\|𝐱_1^{\prime \prime }\|^4}}`$ $``$ $`4\pi ^2{\displaystyle _\xi ^{\mathrm{}}}{\displaystyle \frac{d\rho }{\rho ^2}}.`$ (80) After these approximations, we finally obtain $`N_1={\displaystyle \frac{V\pi ^{1/2}}{8}}{\displaystyle \frac{M^1}{(8\pi )^3}}(L_1L_2)^{1/2}\xi ^1.`$ (81) ## VI Calculation of Second an Third Feynman Diagrams in Fig. LABEL:@fff Here we have to calculate $`N_2=\kappa ^2\underset{\genfrac{}{}{0pt}{}{n_10}{n_20}}{lim}{\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}{\displaystyle \frac{dz_1}{2\pi i}}{\displaystyle \frac{dz_2}{2\pi i}}e^{(z_1L_1+z_2L_2)}`$ (82) $`\times {\displaystyle }d^3x_1d^3x_2{\displaystyle }d^3x_1^{}d^3x_1^{\prime \prime }d^3x_2^{}`$ (83) $`\times \|\psi _1^{a_1}(𝐱_1)\|^2\|\psi _2^{a_2}(𝐱_2)\|^2(𝐀_1\mathrm{\Psi }_1^{}\mathbf{}\mathrm{\Psi }_1)_{𝐱_1^{}}`$ (84) $`\times (𝐀_1\mathrm{\Psi }_1^{}\mathbf{}\mathrm{\Psi }_1)_{𝐱_1^{\prime \prime }}(𝐀_2^2\|\mathrm{\Psi }_2\|^2)_{𝐱_2^{}}`$ (85) The above amplitude has no ultraviolet divergence, so that no regularization is required. The Wick contractions pictured in the second Feynman diagrams of Fig. 2 lead to the integral $`N_2=2VL_2^{1/2}L_1^1{\displaystyle \frac{M}{(2\pi )^6}}{\displaystyle _0^1}𝑑t{\displaystyle _0^t}𝑑t^{}C(t,t^{})`$ (86) where $`C(t,t^{})`$ is a function independent of $`L_1`$ and $`L_2`$: $`C(t,t^{})`$ $`=`$ $`\left[(1t)t^{}(tt^{})\right]^{3/2}`$ (90) $`\times {\displaystyle }d^3xd^3yd^3ze^{M(𝐲𝐱)^2/2(1t)}`$ $`\times \left(_𝐲^\nu e^{M𝐲^2/2t^{}}\right)\left(_𝐱^\mu e^{M𝐱^2/2(tt^{})}\right)`$ $`\times {\displaystyle \frac{\left[\delta _{\mu \nu }𝐳(𝐳+𝐱)\left(z+x\right)_\mu z_\nu \right]}{\|𝐳\|^3\|𝐳+𝐱\|^3}}.`$ As in the previous section, the variables $`𝐱,𝐲,𝐳`$ have been rescaled with respect to the original ones in order to extract the behavior in $`L_1`$. If the polymer lengths are much larger than the persistence length one can ignore the fact that the monomers have a finite size and it is possible to compute $`C(t,t^{})`$ analytically, leading to $`N_2`$ $`=`$ $`{\displaystyle \frac{VL_2^{1/2}L_1^1}{4^3(2\pi )^6}}M^{1/2}\sqrt{2}K,`$ (91) where $`K`$ is the constant $`K`$ $``$ $`{\displaystyle \frac{1}{6}}B({\displaystyle \frac{3}{2}},{\displaystyle \frac{1}{2}})+{\displaystyle \frac{1}{2}}B({\displaystyle \frac{5}{2}},{\displaystyle \frac{1}{2}})`$ (93) $`B({\displaystyle \frac{7}{2}},{\displaystyle \frac{1}{2}})+{\displaystyle \frac{1}{3}}B({\displaystyle \frac{9}{2}},{\displaystyle \frac{1}{2}})={\displaystyle \frac{19\pi }{384}}0.154,`$ and $`B(a,b)=\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)/\mathrm{\Gamma }(a+b)`$ is the Beta function. For large $`L_1\mathrm{}`$, this diagram gives a negligible contribution with respect to $`N_1`$. The third diagram in Fig. 2 give the same as the second, except that $`L_1`$ and $`L_2`$ are interchanged. $$N_3=N_2\|_{L_1L_2}.$$ (94) ## VII Calculation of Fourth Feynman Diagram in Fig. LABEL:@fff Here we have the integral $`N_4`$ $`=`$ $`4\kappa ^2{\displaystyle \frac{1}{2}}\underset{\genfrac{}{}{0pt}{}{n_10}{n_20}}{lim}{\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}{\displaystyle \frac{dz_1}{2\pi i}}{\displaystyle \frac{dz_2}{2\pi i}}e^{(z_1L_1+z_2L_2)}`$ (98) $`\times {\displaystyle }d^3x_1d^3x_2{\displaystyle }d^3x_1^{}d^3x_2^{}d^3x_1^{\prime \prime }d^3x_2^{\prime \prime }`$ $`\times \|\psi _1^{a_1}(𝐱_1)\|^2\|\psi _2(𝐱_2^{a_2})\|^2(𝐀_1\mathrm{\Psi }_1^{}\mathbf{}\mathrm{\Psi }_1)_{𝐱_1^{}}(𝐀_1\mathrm{\Psi }_1^{}\mathbf{}\mathrm{\Psi }_1)_{𝐱_1^{\prime \prime }}`$ $`\times (𝐀_2\mathrm{\Psi }_2^{}\mathbf{}\mathrm{\Psi }_2)_{𝐱_2^{}}(𝐀_2\mathrm{\Psi }_2^{}\mathbf{}\mathrm{\Psi }_2)_{𝐱_2^{\prime \prime }}.`$ which has no ultraviolet divergence. After some effort we find $`N_4={\displaystyle \frac{1}{24^6}}{\displaystyle \frac{M^3V}{(2\pi )^{11}}}(L_1L_2)^{1/2}`$ (99) $`\times {\displaystyle _0^1}ds{\displaystyle _0^s}ds^{}{\displaystyle _0^1}dt{\displaystyle _0^t}dt^{}C(s,s^{},t,t^{}),`$ (100) where $`C(s,s^{};t,t^{})=\left[(1s)s^{}(ss^{})\right]^{3/2}\left[(1t)t^{}(tt^{})\right]^{3/2}`$ (101) $`\times {\displaystyle }{\displaystyle \frac{d^3p}{(2\pi )^3}}[ϵ_{\mu \lambda \alpha }{\displaystyle \frac{p^\alpha }{𝐩^2}}ϵ_{\nu \rho \beta }{\displaystyle \frac{p^\beta }{𝐩^2}}+ϵ_{\mu \rho \alpha }{\displaystyle \frac{p^\alpha }{𝐩^2}}ϵ_{\nu \lambda \beta }{\displaystyle \frac{p^\beta }{𝐩^2}}]`$ (102) $`\times [{\displaystyle }d^3x^{}d^3y^{}e^{i\sqrt{L_1}𝐩(𝐱^{}𝐲^{})}e^{M𝐱^{}{}_{}{}^{2}/2(1s)}`$ (103) $`\times \left(_𝐲^{}^\nu e^{M𝐲^{}{}_{}{}^{2}/2t^{}}\right)\left(_𝐱^{}^\mu e^{M(𝐱𝐲)^2/2(ss^{})}\right)]`$ (104) $`\times [{\displaystyle }d^3u^{}d^3v^{}e^{i\sqrt{L_2}𝐩(𝐮^{}𝐯^{})}e^{M𝐯^{}{}_{}{}^{2}/2(1t)}`$ (105) $`\times \left(_𝐮^{}^\rho e^{M𝐮^{}{}_{}{}^{2}/2t^{}}\right)\left(_𝐯^{}^\lambda e^{M(𝐮^{}𝐯^{})^2/2(tt^{})}\right)]`$ (106) and $`𝐱^{},𝐲^{}`$ are scaled variables. To take into account the finite persitence length, they should be defined on a lattice with spacing $`\xi /\sqrt{L_1}`$. Similarly, $`𝐮^{},𝐯^{}`$ should be considered on a lattice with spacing $`\xi /\sqrt{L_2}`$. Without performing the space integrations $`d^3𝐱^{}d^3𝐲^{}d^3𝐮^{}d^3𝐯^{}`$, the behavior of $`N_4`$ as a function of the polymer lengths can be easily estimated in the following limits: 1. $`L_11;L_1L_2`$ $`N_4L_1^1`$ (107) 2. $`L_21;L_2L_1`$ $`N_4L_2^1`$ (108) 3. $`L_1,L_21,L_2/L_1=\alpha =\text{finite}`$ $`N_4L_1^{3/2}`$ (109) Moreover, if the lengths of the polymers are considerably larger than the persistence length, the function $`C(s,s^{},t,t^{})`$ can be computed in a closed form: $`N_4`$ $``$ $`{\displaystyle \frac{1}{(2\pi )^5}}{\displaystyle \frac{1}{(2\pi )^{3/2}}}{\displaystyle \frac{32}{\sqrt{2}}}(L_1L_2)^{1/2}M^{1/2}V`$ (110) $`\times `$ $`{\displaystyle _0^1}𝑑s{\displaystyle _0^1}𝑑t(1s)(1t)(st)^{1/2}`$ (111) $`\times `$ $`\left[L_1t(1s)+L_2(1t)s\right]^{1/2}.`$ (112) It is simple to check that this expression has exactly the above behaviors. ## VIII Final Result Collecting all contributions we obtain the result for the second topological moment: $`m^2={\displaystyle \frac{N_1+N_2+N_3+N_4}{Z}},`$ (113) with $`N_1,N_2,N_3,N_4,Z`$ given by Eqs. (47), (81), (91), (94), and (112). In all formulas, we have assumed that the volume $`V`$ of the system is much larger than the size of the volume occupied by a single polymer, i.e., $`VL_1^3`$ To discuss the physical content of the result (113), we assume $`P_2`$ to be a long effective polymer representing all polymers in a uniform solution. We introduce the polymer concentration $`\rho `$ as the average mass density of the polymers per unit volume: $`\rho ={\displaystyle \frac{M}{V}}`$ (114) where $`M`$ is the total mass of the polymers $`M={\displaystyle \underset{i=1}{\overset{N_p}{}}}m_a{\displaystyle \frac{L_k}{a}}.`$ (115) Here $`m_a`$ is the mass of a single monomer of length $`a`$, $`L_k`$ is the length of polymer $`P_k`$, and $`N_p`$ is the total number of polymers. Thus $`L_k/a`$ is the number of monomers in the polymer $`P_k`$. The polymer $`P_1`$ is singled out as any of the polymers $`P_k`$, say $`P_{\overline{k}}`$, of length $`L_1=L_{\overline{k}}`$. The remaining ones are replaced by a long effective polymer $`P_2`$ of length $`L_2=\sigma _{k\overline{k}}L_k`$. From the above relations we may also write $`L_2{\displaystyle \frac{aV\rho }{m_a}}`$ (116) In this way, the length of the effective molecule $`P_2`$ is expressed in terms of physical parameters, the concentration of polymers, the monomer length, and the mass and volume of the system. Inserting (116) into (113), with $`N_1,N_2,N_3,N_4,Z`$ given by Eqs. (47), (81), (91), (94), and (112). and keeping only the leading terms for $`V1`$, we find for the second topological moment $`m^2`$ the approximation $`m^2{\displaystyle \frac{N_1+N_2}{Z}},`$ (117) and this has the approximate form $`m^2={\displaystyle \frac{a\rho }{m_a}}\left[{\displaystyle \frac{\xi ^1L_i}{16\pi ^{1/2}M^2}}{\displaystyle \frac{K\sqrt{2}L_1^{1/2}}{(2\pi )^4M^{3/2}}}\right],`$ (118) with $`K`$ of (93). ## IX Summary We have set up a topological field theory to describe two fluctuating polymers $`P_1`$ and $`P_2`$, and calculated the second topological moment for the linking number $`m`$ between $`P_1`$ and $`P_2`$. The result is used to calculate the second moment for a single polymer with respect to all others in a solution of many polymers. In forthcoming work we shall calculate the effect of the excluded volume. ## X Appendix In this appendix we present the computations of the amplitudes $`N_1,\mathrm{},N_4`$. We shall need the following simple tensor formulas involving two completely antisymmetric tensors $`\epsilon ^{\mu \nu \rho }`$: $`\epsilon _{\mu \nu \rho }\epsilon ^{\mu \alpha \beta }=\delta _\nu ^\alpha \delta _\rho ^\beta \delta _\nu ^\beta \delta _\rho ^\alpha `$ (119) $`\epsilon _{\mu \nu \rho }\epsilon ^{\mu \nu \beta }=2\delta _\rho ^\beta .`$ (120) The Feynman diagrams shown in Fig. 2) corresponds to a product of four correlation functions $`G_0`$ of Eq. (51), which have to be integrated over space and Laplace transformed. For the latter we make use of the convolution property of the integral over two Laplace transforms $`\stackrel{~}{f}(z)`$ and $`\stackrel{~}{g}(z)`$ of the functions $`f,g`$: $`{\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}{\displaystyle \frac{dz}{2ni}}e^{zL}\stackrel{~}{f}(z)\stackrel{~}{g}(z)={\displaystyle _0^L}𝑑sf(s)g(LS)`$ (121) All spatial integrals are Gaussian of the form $`{\displaystyle d^3xe^{a𝐱+2b𝐱𝐲}}=(2n)^{3/2}a^{3/2}e^{b^2y^2/a},a>0.`$ (122) We are now ready to evaluate $`N_1`$ in Eq. (76). Taking the limit of vanishing replica indices we find with the help of Eqs. (119)-(121): $`N_1`$ $`=`$ $`{\displaystyle d^3x_1},d^3x_2{\displaystyle _0^{L_1}}𝑑s{\displaystyle _0^{L_2}}𝑑t{\displaystyle d^3x_1^{}d^3x_2^{}}`$ (123) $`\times `$ $`G_0(𝐱_1𝐱_1^{};s)G_0(𝐱_1^{}𝐱_1;L_1S)`$ (124) $`\times `$ $`G_0(𝐱_2𝐱_2^{};t)G_0(𝐱_2^{}𝐱_2;L_2t){\displaystyle \frac{l}{\|𝐱_1^{}𝐱_2^{}\|s}}.`$ (125) Performing the changes of variables $`s^{}={\displaystyle \frac{s}{L_1}}t^{}={\displaystyle \frac{t}{L_2}}𝐱={\displaystyle \frac{𝐱_1𝐱_1^{}}{\sqrt{4}}}𝐲={\displaystyle \frac{𝐱_2𝐱_2^{}}{(\sqrt{L_2})}}`$ (126) and setting $`𝐱_1^{\prime \prime }𝐱_1^{}𝐱_2^{}`$, we easily derive (78). For small $`\xi /\sqrt{L_1}`$ and $`\xi /\sqrt{L_2}`$, we d use the approximation (80), the space integrals can be done the formula (122). After some calculation one finds the final result of Eq. (108). The amplitude $`N_2`$ calculated quite similarly. Contracting the fields in Eq. (85), and keeping only the contributions which do not vanish in the limit of zero replica indices, we arrive at $`N_2`$ $`=`$ $`{\displaystyle d^3x_1d^3x_2d^3x_1^{}d^3x_1^{\prime \prime }d^3x_2^{}}`$ (127) $`\times `$ $`[{\displaystyle _0^{L_1}}ds{\displaystyle _0^S}ds^{}G_0(𝐱_1^{}𝐱_1;L_1s)`$ (128) $`\times `$ $`_{x_1^{\prime \prime }}^\nu G_0(𝐱_1𝐱_1^{\prime \prime };s^{})_{x_1^{}}^\mu G_0(𝐱_1^{\prime \prime }𝐱_1^{};ss^{})]`$ (129) $`\times `$ $`D_{\mu \lambda }(𝐱_1^{}𝐱_2)D_{\nu \lambda }(𝐱_1^{\prime \prime }𝐱_2^{})`$ (130) $`\times `$ $`\left[{\displaystyle _0^{L_2}}𝑑tG_0(𝐱_2𝐱_2^{};L_2t)G_0(𝐱_2^{}𝐱_2;t)\right].`$ (131) where $`D_{\mu \nu }(𝐱,𝐱^{})`$ are the correlation functions (11)–(13) of the vector potentials. Setting $`𝐱_2\sqrt{L_2}𝐮+𝐱_2^{}`$ and supposing that $`\xi /\sqrt{L_2}`$ is small, the integral over $`𝐮`$ can be easily evaluated with the help of the Gaussian integral (122). After the substitutions $`𝐱_1^{\prime \prime }=\sqrt{L_1}𝐲+𝐱_1`$ $`𝐱_1^{}=\sqrt{L_1}(𝐲𝐱)+𝐱_1`$, $`𝐱_2^{}=\sqrt{L_1}(𝐲𝐱𝐳)+𝐱_1`$ and a rescaling of the variables $`s,s^{}`$ by a factor $`L_1^1`$, we derive Eq. (86) with (90). For small $`\xi /\sqrt{L_1},\frac{3}{\sqrt{L_2}}`$, the spatial integrals are easily evaluated leading to: $`N_2`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}VL_2^{1/2}L_1^1M^{1/2}}{(4n)^6}}`$ (133) $`{\displaystyle _0^1}𝑑t{\displaystyle _0^t}𝑑t^{}t^{}(1t)\sqrt{{\displaystyle \frac{tt^{}}{1t+t^{}}}}`$ After the change of variable $`t^{}t^{\prime \prime }=tt^{}`$, the double integral is reduced to a sum of integrals the type $`c(n,m)={\displaystyle _0^1}𝑑tt^m{\displaystyle _0^t}𝑑t^{}t^{}{}_{}{}^{n}\sqrt{{\displaystyle \frac{t^{}}{1t^{}}}},m,n=\mathrm{integers}.`$ (134) These can be simplified by replacing $`t^m`$ by $`dt^{m+1}/dt(m+1)`$, and doing the integrals by parts. In this way, we end up with a linear combination of integrals of the form: $`{\displaystyle _0^1}𝑑t{\displaystyle \frac{t^{\mu +\frac{1}{2}}}{\sqrt{1t}}}=B(\mu +{\displaystyle \frac{3}{2}},{\displaystyle \frac{1}{2}}).`$ (135) The calculations of $`N_3`$ and $`N_4`$ are very similar, and may be omitted here.
warning/0003/cond-mat0003147.html	ar5iv	text	# Evidence of quantum criticality in the doped Haldane system Y2BaNiO5 \[ ## Abstract Experimental bulk susceptibility $`\chi (T)`$ and magnetization $`M(H,T)`$ of the S=1 Haldane chain system doped with nonmagnetic impurities, Y<sub>2</sub>BaNi<sub>1-x</sub>Zn<sub>x</sub>O<sub>5</sub> ($`x0.08`$), are analyzed. A numerical calculation for the low-energy spectrum of non-interacting open segments describes very well experimental data above 4 K. Below 4 K, we observe power-law behaviors, $`\chi (T)`$$``$ $`T`$ and $`M(H,T)`$ $`T`$<sup>1-α</sup>$`f`$<sub>α</sub>$`(H/T)`$, with $`\alpha `$ $`(<1)`$ depending on the doping concentration $`x`$. This observation suggests the appearance of a gapless quantum phase due to a broad distribution of effective couplings between the dilution-induced moments. \] The past two decades have seen a resurgence of interest devoted to macroscopic quantum phenomena in Heisenberg antiferromagnets (HAF). A current issue in this field concerns the interplay between quantum spin fluctuations and quenched (i.e., time-independent) disorder in one dimension (1D) . A simple nontrivial model to address this issue is the 1D-HAF with spins $`S=1`$, which has a Haldane gap $`\mathrm{\Delta }0.4J`$ and a short correlation length $`\xi 6`$ in its spin-liquid ground state ($`J`$ is the nearest-neighbor exchange parameter) . In this system, the disorder in the form of site depletion leads to interesting quantum phenomena due to the creation of two effective $`S=1`$/2 spins on the opposite edges of a segment, near the vacancies . The inclusion of weak magnetic bonds across the spin-vacancies induces bond disorder and causes departure from simple finite-size behaviors. In this case, some recent theoretical works, for the $`S=1`$ chain with variable nonmagnetic doping, found a quantum (zero-temperature) phase transition from the Haldane phase to a 1D random-singlet (RS) phase in which effective spins are coupled into singlets over all length scales . On the experimental side, however, this problem has been very little explored from the point of view of quantum criticality. While it is known that some $`S=1`$ diluted compounds can sustain paramagnetism down to low temperature, most of the existing experimental studies aimed to confirm the existence of $`S=1/2`$ end-chain states or to reveal the internal structure of these end states . A noticeable exception is the recent observation of a 3D long-range order (LRO) caused by nonmagnetic doping in PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> . This Haldane system shows, however, substantial interchain interactions, so it is almost critical towards the formation of a LRO in the absence of intentional disorder. In this work, we examine the effects of nonmagnetic doping on the bulk magnetic responses in the $`S=1`$ quasi-1D-HAF Y<sub>2</sub>BaNiO<sub>5</sub>. We chose this well characterized system because of the rather high value of its main energy scale, $`J`$$``$280 K, and the very small value of the interchain coupling, $`\|J_{}/J\|<10^3`$ . Furthermore, it has a simple orthorhombic crystal structure and a small single-ion anisotropy ($`D/J`$$``$-0.04, $`E/J`$$``$-0.01) . A weak nonmagnetic doping in Y<sub>2</sub>BaNiO<sub>5</sub>, which is realized by substituting Zn<sup>2+</sup> or Mg<sup>2+</sup> ions for the $`S=1`$ Ni<sup>2+</sup> ions, yields new magnetic states below the Haldane gap ($`\mathrm{\Delta }`$$``$100 K) with no sign of LRO down to very low temperature . Recently, the low-temperature specific heat of Y<sub>2</sub>BaNi<sub>1-x</sub>Zn<sub>x</sub>O<sub>5</sub> ($`x=0.04`$) and the ESR spectra of a series of Mg-doped compounds have been quantitatively explained by an effective model which describes the low-energy spectrum of non-interacting open $`S=1`$ chains . These results and the recent NMR experiments in Ref. 8 are strong indications of the existence of $`S=1/2`$ end-chain states. In the present work, we analyze new susceptibility and magnetization data for Y<sub>2</sub>BaNi<sub>1-x</sub>Zn<sub>x</sub>O<sub>5</sub> with 0.04$``$x$``$0.08. Above 4 K, the simple model of non-interacting segments used in Refs. describes the data very well. However, at very low temperature we observe a new regime with divergent power law behavior, which is similar to that observed in some doped semiconductors . Our results support the existence of a zero-temperature gapless phase due to randomness in the effective bond distribution that is reinforced by the interchain interactions. The low-energy effective model for this phase resembles a 2D or 3D RS state. Figure 1 shows the linear susceptibilities, $`\chi (T)`$, measured with a commercial SQUID magnetometer for a series of polycrystalline samples of Y<sub>2</sub>BaNi<sub>1-x</sub>Zn<sub>x</sub>O<sub>5</sub> ($`x=0`$, 0.04, 0.06 and 0.08) that were prepared through standard solid-state reactions . The results obtained for the nominally pure sample compare well with those in the published literature . At low temperatures, $`\chi (T)`$ is well fitted by the sum of a $`T`$-independent term $`\chi `$<sub>0</sub>, a Curie law and a thermally activated term with $`\mathrm{\Delta }`$$``$100 K (see Fig.1). The sum of the activated part and the constant $`\chi _010^4`$ cm<sup>3</sup>/mol is negligible below 15 K. The low-temperature susceptibility (below 10-15 K) for $`x=0`$ is mainly due to the Curie contribution that corresponds to 1.4 percent of free $`S=1/2`$ spins per formula unit. This Curie behavior can be interpreted as coming from natural crystal defects and a weak excess oxygen. The main effect of the intentional nonmagnetic dilution is to increase drastically the low-temperature upturn in $`\chi (T)`$, as is evident from Fig. 1. This indicates the existence of paramagnetic moments whose amount is related to the doping concentration. The behavior of $`\chi (T)`$ is however qualitatively the same, regardless of the level of dilution. The data for the doped compounds were first analyzed with a modified Curie-Weiss law, $`\chi (T)=\chi _0+C/(T\theta )`$, between 1.9 and 15 K. A small phenomenological $`\theta `$$``$-0.6 K was required to describe the data. The dependence of the Curie constant $`C`$ as a function of the Zn doping is shown in the inset of Fig. 1. At low doping, the observed $`C`$ value is that for two $`S=1/2`$ spins per nonmagnetic impurity. At high $`x`$, however, the Curie constant should be that for one half of a spin-1 per Zn since a short segment with an odd (even) number of spins behaves as a $`S=1`$ ($`S=0`$) ’molecular’ entity at low temperature and the number of odd chains is half the total number of chains. This tendency seems to be observed (see Fig. 1). We have also measured the magnetization, $`M(H,T)`$, for applied field up to 5 T. For all samples, $`M(H,T)`$-$`\chi `$<sub>0</sub>$`H`$ shows a non-Brillouin behavior and is smaller than twice the Brillouin function for one $`S=1/2`$ spin per Zn. To go further in our analysis, we have compared the experimental data with a theoretical model based on the picture of non-interacting segments, as described hereafter. The energy spectrum of finite segments described by the $`S=1`$ Heisenberg Hamiltonian is characterized by four low-energy states, a singlet and a triplet . It corresponds to the existence of a coupling between the two $`S=1/2`$ edge spins, which may be AF or ferromagnetic (F). The energy difference among these levels decreases exponentially with the length of the chain $`N`$, as $`(1)^Ne^{N/\xi }`$, and already for chains of a few lattice sites, these four states are separated from the rest of the spectrum by an energy of the order of the Haldane gap . Therefore at temperature $`T`$ and magnetic energies $`\mu _BH`$ much lower than the Haldane gap, the magnetic and thermal properties of a system composed of segments of several lengths are completely described by the energy of these four states and the matrix elements of the magnetization operator in this reduced subspace. We have calculated these quantities using the Density Matrix Renormalization Group (DMRG) method. The resulting Hamiltonian which describes the low-energy properties of a segment including the triplet $`\|1,S_z`$ and the singlet state $`\|0`$ is: $`H_{eff}`$ $`=`$ $`E_0(N)+(J\alpha (N)+D\beta (N))\|00\|+D\gamma (N)S_z^2`$ (2) $`+E\gamma (N)(S_x^2S_y^2)\mu _B{\displaystyle \underset{\nu \alpha }{}}H^\alpha g^{\alpha \nu }S_t^\nu ,`$ where $`E_0(N),\alpha (N)`$, $`\beta (N)`$and $`\gamma (N)`$ are functions of the chain length $`N`$ (determined from the DMRG data). The validity of the last term has been verified explicitly by calculating the matrix elements of $`S_t^+`$, and $`S_t^{}`$ for all chains. The magnetic moment operator is -$`_HH_{eff}`$, and to compare with experiments in polycrystalline samples, we have averaged its expectation value over all possible orientations of the crystal. Also, a distribution of chain segments corresponding to a random distribution of defects was assumed . Figure 2 shows the comparison between the measured $`\chi (T)T`$ products for $`T<12`$ K and the numerical solution. The adjustable parameters in the calculations were the doping $`x`$ and $`\chi `$<sub>0</sub>. All the other parameters were held fixed at the values deduced from previous works : $`J=280`$ K, $`D/J=0.039`$, $`E/J=0.013`$, $`g_x=g_y=2.17`$ and $`g_z=2.20`$. The fitted values of $`x`$ ($`x=0.041`$, 0.060 and 0.074) are close to the nominal concentrations, and the fitted $`\chi `$<sub>0</sub> remains in the expected range ($`\chi _0=1.0\times 10^4`$ to $`1.9\times 10^4`$ cm<sup>3</sup>/Ba-mol). Taking into account that uncertainties in the above parameters might affect the theoretical curve, the success of the anisotropic model of non-interacting segments to explain the data above 4 K is remarkable. Note that for the cases with nominal concentration of Zn $`x=0.06`$ and $`x=0.08`$, a decrease of the order of 5% or less in the actual $`x`$ used, would improve considerably the fitting above 5 K, but increase the differences between theory and experiment in the low-temperature part. Instead, because of the different rate of change of both curves, it is not possible to obtain a good fit for $`T<4`$ K by small changes of $`x`$ or other parameters. A good agreement between the experimental magnetization ($`M(H,T)`$-$`\chi `$<sub>0</sub>$`H`$) and the numerical data above 4 K (not shown) is also achieved using the parameters required to describe the susceptibility data. This noticeable agreement confirms the existence of the $`S=1/2`$ end-chain excitations. Note that an accurate description of the experimental data needs to consider both, a random distribution of defects and the effect of the single-ion anisotropy , which is stronger for shorter segments. However, below 4 K, the experimental data deviates from the predictions of the above described model. This is seen in Fig. 2 for the susceptibility. It is also observed for our magnetization data $`M(H,T)`$-$`\chi `$<sub>0</sub>$`H`$. In previous studies of specific heat and electron spin resonance, while excellent agreement was obtained above 4 K, a discrepancy was also found below this temperature and interpreted as the onset of interchain interactions not taken into account in this model . An analysis of the susceptibility shows that, in the vicinity of 4 K, there is a crossover to a sub-Curie power law regime, $`[\chi (T)\chi _0]T^\alpha `$ with $`\alpha =0.83`$, 0.79 and 0.76 for $`x=0.04`$, 0.06 and 0.08, respectively (see Fig. 3). The susceptibility data for the Mg-doped compound Y<sub>2</sub>BaNi<sub>0.959</sub>Mg<sub>0.041</sub>O<sub>5</sub> in Ref. also obey a power-law form with $`\alpha =0.73`$ (see Fig. 3). This behavior, which cannot be reproduced by the model of non-interacting segments, is reminiscent of some random exchange Heisenberg systems for which the AF couplings between nearest-neighbor $`S=1/2`$ spins are random in their magnitude, such as insulating phosphorous-doped silicon (Si:P) . Bhatt and Lee have proposed a theoretical method to explain the properties of Si:P . Within this approach, the spin pairs with the strongest AF bonds freeze into inert singlets, leaving behind a new ensemble of active spins with a renormalized distribution of exchange, and the quantum fluctuations drive the whole system into a ground state consisting of random local singlets (RS phase). The behavior of the susceptibility, $`\chi (T)`$$``$$`T`$ with $`\alpha <1`$, follows from a divergent power-law distribution of the renormalized exchange, $`P(J)`$$``$$`J`$. A scaling of the magnetization has been also predicted and experimentally observed , $`M(H,T)/T^{1\alpha }=f_\alpha (H/T)`$. A similar scaling is, in fact, very well obeyed by our data below (but not above) 4 K using the $`\alpha `$ exponents determined from the susceptibility. Fig. 3 shows a plot of $`(M(H,T)\chi _0H)/T^{1\alpha }`$ versus $`H/T`$ for each of the three samples studied. All the data for each sample, for different magnetic fields and temperatures, lie on a single scaling curve, which is however different from the expression of $`f_\alpha (H/T)`$ calculated on the basis of the Bhatt and Lee’s solution for Si:P . This is illustrated in Fig. 3 for Y<sub>2</sub>BaNi<sub>0.94</sub>Zn<sub>0.06</sub>O<sub>5</sub>. These power law behaviors show that the picture of the ideal 1D anisotropic Heisenberg non-interacting chain breaks down below 4 K, and that additional interactions play a role. From simple arguments based on perturbation theory, one expects that the dominant inter-chain interaction is a coupling $`J`$<sub>b</sub> between $`S=1`$ spins lying in nearest-neighbor chains along the $`b`$ direction, perpendicular to the chain axis $`a`$. A calculation using the cell perturbation method gives $`J`$<sub>b</sub>$``$0.2 K . This interaction should be at least an order of magnitude larger than that between two spins on the same atomic chain, with a Zn atom in between. As discussed in more detail in Ref. $`J`$<sub>b</sub> induces an effective interaction $`J^{}`$ between $`S=1/2`$ end-states of neighboring chains along the $`b`$ direction, the magnitude and sign of which depends on the detailed position of the defects. The maximum absolute value of $`J^{}`$ is of the order of $`10J_b`$ and it should decay exponentially with the distance along the chain between the defects. Since the interaction between $`S=1/2`$ defects lying in different $`ab`$ planes is expected to be at least an order of magnitude smaller, the physics in the range 0.1 K$`<T<`$4 K seems to correspond to a 2D array of $`S=1/2`$ spins, with random F and AF interactions coming from exponentially decaying intra- and interchain couplings. In agreement with this, the spin glass behavior observed below 3 K in Ca doped systems Y<sub>2-y</sub>Ca<sub>y</sub>BaNiO<sub>5</sub> indicates a dimensionality higher than one in the system at low temperature . For a 2D or 3D $`S=1/2`$ model with random AF interactions, Bhatt and Lee obtained a power law behavior with $`\alpha `$ $`(<1)`$ depending on the doping concentration . This result is consistent with our observations but, as mentioned above, we obtain a different scaling function for the magnetization. This might be due to the presence of effective F coupling in Y<sub>2</sub>BaNi<sub>1-x</sub>Zn<sub>x</sub>O<sub>5</sub>. Our results can also be compared to a recent theoretical work in which the spin-1 chain with an AF coupling across nonmagnetic impurities is mapped onto a bond disordered spin-1/2 chain . Within this 1D approach, the Haldane state is found to be stable to weak dilution. For sufficiently strong randomness, however, the system flows towards a 1D RS phase with diverging correlation length. In this gapless phase, the form of the low energy asymptotic behavior of $`\chi (T)`$ is independent of the level of the disorder, $`\chi (T)1/(T\mathrm{ln}^2T)`$. In the Haldane phase, there is a Griffiths gapless region where $`\chi (T)`$ varies as $`T`$ with $`\alpha `$ $`(<1)`$ which is dependent of the details of the randomness. It is unclear however whether a similar physics should be valid or not when interchain coupling is turned on . Furthermore, the form of the asymptotic behavior of $`\chi (T)`$ is not known for a 2D or 3D RS phase . If there were no interchain coupling $`J`$, the low-energy model for Y<sub>2</sub>BaNi<sub>1-x</sub>Zn<sub>x</sub>O<sub>5</sub> with $`0.04x0.08`$ would be the 1D RS phase depicted in Ref. since the bond across a Zn is expected to be much weaker than the average coupling between the $`S=1/2`$ edges within a segment . From our observations, it is tempting to say that the $`H/T`$ variable in the scaling of $`M(H,T)`$ indicates a $`T_c=0`$ critical point and that Y<sub>2</sub>BaNi<sub>1-x</sub>Zn<sub>x</sub>O<sub>5</sub> has a quantum critical ground state even for the smallest doping $`x=0.04`$. From the theoretical point of view, the detailed explanation of the observed power-law behaviors remains open. Interestingly, a similar singular behavior has been disclosed for two $`S=1`$ 1D-HAF compounds, namely AgVP<sub>2</sub>S<sub>6</sub> and NENP, on nominally pure samples . A broad distribution of couplings between the dilution-induced moments and quantum fluctuations drives Y<sub>2</sub>BaNi<sub>1-x</sub>Zn<sub>x</sub>O<sub>5</sub>, and perhaps other real systems as well, towards a novel ground state which resembles a random valence bond state. This is in contrast with the LRO observed in PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> . Therefore there is no universal behavior under nonmagnetic doping in real $`S=1`$ 1D-HAF substances. Details of the topology of the magnetic interaction lattice are probably crucial in determining the way the Haldane phase is destabilized by the quenched disorder. K.H. and C.D.B. are supported by CONICET, Argentina. A.A.A. is partially supported by CONICET. C.P. is partially supported by the ’Institut Universitaire de France’. C.P. thanks H. Mutka for helpful discussions and continued support. K. H. is grateful to the Physics Department of the University of Buenos Aires for the hospitality during this work. This work was supported by PICT 03-00121-02153 of ANPCyT and PIP 4952/96 of CONICET.
warning/0003/hep-th0003051.html	ar5iv	text	# 1 Introduction ## 1 Introduction Superfield representations of super-Poincaré and superconformal algebras have been proved to be useful tools since the early development of supersymmetry for several reasons. They provide the natural framework to formulate supersymmetric field theories in a “covariant fashion” and allow one, in many cases, to achieve a simple understanding of the softening of “quantum divergences”. This milder quantum behaviour of supersymmetric field theories is at the basis of the so-called “non-renormalization theorems” which are one of the striking features of supersymmetric quantum theories . In modern language, which applies to generic supersymmetric theories, these non-renormalization theorems are due to the fact that supersymmetric field theories have some “field representations” that are short, namely, the component field of highest dimension (which is not a total derivative) lies at a lower $`\theta `$ level than what is naively expected from a generic superfield. Examples of such “short” superfields already appear in $`N=1`$ 4d supersymmetry and they are called “chiral” . In the case of superconformal algebras chiral primaries have a “ring structure” under multiplication and their conformal dimension is quantized in terms of the R $`U(1)`$ charge. In $`N`$-extended supersymmetry in $`d=4`$ as well as in other dimensions one needs to generalize the notion of “chiral superfields”. The point is that the shortening is often due to an interplay between the conformal dimension and the (non-Abelian) R-symmetry quantum numbers. The latter, in $`d=3`$ and 6 are related to the Dynkin labels of the $`SO(N)`$ and $`USp(2N)`$ groups while in $`d=4`$ for $`N2`$, to the Dynkin labels of $`SU(N)`$. Extended superspaces, enlarged with coordinates on $`G/H`$ where $`G`$ is the R-symmetry of the superconformal algebra and $`H`$ is a maximal subgroup (with rank of $`H`$ = rank of $`G`$) are called harmonic superspaces . They provide the suitable framework in which the notion of chirality is generalized to Grassmann analyticity . For these “short” superfields the superconformal algebra is realized in a subspace of the full superspace which contains a reduced number of the original anticommuting Grassmann variables. In the spirit of the AdS/CFT correspondence where boundary “conformal operators” of $`CFT_d`$ are mapped onto “bulk states” in $`AdS_{d+1}`$, multiplet shortening translates into a BPS condition on massive (and massless) particle states in anti-de Sitter space (see, for instance, ). Superconformal algebras in $`d`$ dimensions appear as vacuum symmetries of string or M-theory compactified on $`AdS_{d+1}`$. Massive BPS saturated UIR’s of these algebras should therefore be relevant to classify solitons preserving different fractions of supersymmetry, as it happens in the corresponding flat space limit. The general analysis of multiplet shortening is related to the so-called “unitary bounds” of UIR’s of superconformal algebras. For the $`d=4`$ case the latter was obtained in the 80’s in Ref. for $`N=1`$ and in Ref. for arbitrary $`N`$. The relation with the multiplet shortening and the $`AdS_5/CFT_4`$ correspondence was recently spelled out in . The superfield analysis in $`CFT_d`$ is “dual” to the “state” analysis on $`AdS_{d+1}`$ since the same superalgebra acts on these representation spaces. However, the superfield approach is more powerful not only because it allows one to treat quantum field theories but because it leads to a simpler classification of “massive representations” in the language of composite operators. The different BPS conditions in $`AdS_d`$ are rephrased to the different Grassmann analytic operators (generalizations of “chiral operators”) which exist in extended harmonic superspace. The full classification of all BPS conditions was carried out for $`d=4,6`$ superconformal algebras in Refs. and it is extended to the $`d=3`$ $`N=8`$ superconformal algebra in the present paper. The appropriate superconformal algebra is in this case $`OSp(8/4,)`$ which is a different non-compact form of the superalgebra which occurs in the $`(2,0)`$ theory in $`d=6`$. The latter is related to M-theory on $`AdS_7\times S^4`$. The former is appropriate to the $`AdS_4\times S_7`$ compactification of M-theory and some of its representations, both massless and massive, have been widely considered in the literature (see, e.g., ). The purpose of this paper is to extend the harmonic superspace analysis to the $`d=3`$ $`N=8`$ case in order to obtain all BPS states which may occur in $`AdS_4`$. These are the $`AdS`$ analogues of the 1/2, 1/4 and 1/8 BPS states of Poincaré supersymmetry which occur in the classification of extremal black holes in supergravity theories . Therefore BPS states in $`AdS_4`$ correspond, in particular, to anti-de Sitter black holes of $`N=8`$ gauged $`SO(8)`$ supergravity . The paper is organized as follows. In Section 2 we carry out a general analysis of the short highest weight UIR’s of $`OSp(8/4,)`$. To this end we consider $`OSp(8/4,)`$ as the $`N=8`$ 3d superconformal algebra and study the conditions on the HWS’s which are annihilated by all the $`S`$ (conformal supersymmetry) generators and by a fraction (1/2, 3/8, 1/4 or 1/8) of the $`Q`$ (Poincaré supersymmetry) ones. As a result we find that the Lorentz spin of these HWS’s must vanish and that their conformal dimension should be related to their $`SO(8)`$ Dynkin labels. Such HWS’s generate series of representations exhibiting 1/2, 3/8, 1/4 and 1/8 BPS shortening. The simplest multiplets of maximal shortening (1/2 BPS) are the two distinct “supersingletons”. In Section 3 and 4 we realize the $`N=8`$ supersingletons first as constrained superfields in ordinary superspace and then as Grassmann analytic superfields in harmonic superspace. The latter have the advantage that their analyticity properties are preserved by multiplication. This allows us, in Section 5, to construct all composite operators obtained by multiplying supersingleton superfields and undergoing different shortenings corresponding to different BPS states in the $`AdS_4`$ bulk interpretation. We show that by tensoring only one type of supersingletons we can only construct 1/2, 1/4 and 1/8 BPS states, but by mixing the two types we can reproduce the complete classification of short multiplets from Section 2. In this way we also give an indirect proof that all the representations found in Section 2 are indeed unitary. ## 2 Short highest weight UIR’s of $`OSp(8/4,)`$ In this section we shall derive the general conditions on the highest weight state (HWS) of a short representation of $`OSp(8/4,)`$. The superalgebra $`OSp(8/4,)`$ is the $`N=8`$ superconformal algebra in three dimensions (only the part of the algebra relevant to our argument is shown): $`\{Q_\alpha ^i,Q_\beta ^j\}=2\delta ^{ij}\mathrm{\Gamma }_{\alpha \beta }^\mu P_\mu ,\{S_\alpha ^i,S_\beta ^j\}=2\delta ^{ij}\mathrm{\Gamma }_{\alpha \beta }^\mu K_\mu ,`$ (1) $`\{Q_\alpha ^i,S_\beta ^j\}=\delta ^{ij}M_{\alpha \beta }+2ϵ_{\alpha \beta }(T^{ij}+\delta ^{ij}D),`$ (2) $`[D,Q_\alpha ^i]={\displaystyle \frac{i}{2}}Q_\alpha ^i,[D,S_\alpha ^i]={\displaystyle \frac{i}{2}}S_\alpha ^i,`$ (3) $`[M_{\alpha \beta },Q_\gamma ^i]=i(ϵ_{\gamma \alpha }Q_\beta ^i+ϵ_{\gamma \beta }Q_\alpha ^i),[M_{\alpha \beta },S_\gamma ^i]=i(ϵ_{\gamma \alpha }S_\beta ^i+ϵ_{\gamma \beta }S_\alpha ^i),`$ (4) $`[T^{ij},Q_\alpha ^k]=i(\delta ^{ki}Q_\alpha ^j\delta ^{kj}Q_\alpha ^i),[T^{ij},S_\alpha ^k]=i(\delta ^{ki}S_\alpha ^j\delta ^{kj}S_\alpha ^i),`$ (5) $`[T^{ij},T^{kl}]=i(\delta ^{ik}T^{jl}+\delta ^{jl}T^{ik}\delta ^{jk}T^{il}\delta ^{il}T^{jk}).`$ (6) Here we find the following generators: $`Q_\alpha ^i`$ of $`N=8`$ Poincaré supersymmetry carrying a 3d spinor Lorentz index $`\alpha =1,2`$ and an $`SO(8)`$ vector <sup>1</sup><sup>1</sup>1$`SO(8)`$ has three 8-dimensional representations, $`8_v`$, $`8_s`$ and $`8_c`$. Since these three representations are related by $`SO(8)`$ triality, the choice which one to ascribe to the supersymmetry generators is purely conventional. In order to be consistent with the other $`N`$-extended 3d supersymmetries where the odd generators always belong to the vector representation, we prefer to put an $`8_v`$ index $`i`$ on the supercharges. index $`i=1,\mathrm{},8`$; $`S_\alpha ^i`$ of conformal supersymmetry; $`P_\mu `$, $`\mu =0,1,2`$, of translations; $`K_\mu `$ of conformal boosts; $`M_{\alpha \beta }=M_{\beta \alpha }`$ of the 3d Lorentz group $`SO(2,1)SL(2,)`$; $`D`$ of dilations; $`T^{ij}=T^{ji}`$ of $`SO(8)`$. The definition of a short representation we adopt requires that its HWS is annihilated by part of the Poincaré supersymmetry generators $`Q_\alpha ^i`$. Since the latter are irreducible under the Lorentz and R symmetries, the only way to achieve shortening is to break one of them. Postponing the possibility of dealing with the Lorentz group for a future investigation, here we choose to break $`SO(8)`$ down to $`[SO(2)]^4[U(1)]^4`$ and decompose the $`SO(8)`$ vector $`Q_\alpha ^i`$ into eight independent projections carrying different $`U(1)`$ charges. The first two such projections are: $$Q_\alpha ^{\pm \pm }=\frac{1}{2}(Q_\alpha ^1\pm Q_\alpha ^2)$$ (7) and the corresponding charge generator is $`H_1=2iT^{12}`$, so that $$[H_1,Q_\alpha ^{\pm \pm }]=\pm 2iQ_\alpha ^{\pm \pm }.$$ (8) Note the unusual units of charge, which are spinorial rather than vectorial. Let us write down one of the projections of eq. (2) which will be needed in what follows: $$\{Q_\alpha ^{++},S_\beta ^{}\}=\frac{1}{2}M_{\alpha \beta }+ϵ_{\alpha \beta }(D\frac{1}{2}H_1).$$ (9) Similarly, we introduce the second charge $$Q_\alpha ^{(\pm \pm )}=\frac{1}{2}(Q_\alpha ^3\pm Q_\alpha ^4)$$ (10) with generator $`H_2=2iT^{34}`$, so that $$[H_2,Q_\alpha ^{(\pm \pm )}]=\pm 2iQ_\alpha ^{(\pm \pm )}$$ (11) and $$\{Q_\alpha ^{(++)},S_\beta ^{()}\}=\frac{1}{2}M_{\alpha \beta }+ϵ_{\alpha \beta }(D\frac{1}{2}H_2).$$ (12) The third and fourth charges will be introduced in a different way. The components $`\underset{¯}{i}=5,6,7,8`$ of the $`8_v`$ of $`SO(8)`$ form an $`SO(4)`$ vector. Since $`SO(4)SU(2)\times SU(2)`$, we can rewrite it in spinor notation with the help of the Pauli matrices, e.g., $`Q^{\underset{¯}{i}}Q^{\underset{¯}{a}\underset{¯}{a}^{}}=Q^{\underset{¯}{i}}(\sigma ^{\underset{¯}{i}})^{\underset{¯}{a}\underset{¯}{a}^{}}`$. Doing this in eq. (2) we obtain $$\{Q_\alpha ^{\underset{¯}{a}\underset{¯}{a}^{}},S_\beta ^{\underset{¯}{b}\underset{¯}{b}^{}}\}=\frac{1}{2}ϵ^{\underset{¯}{a}\underset{¯}{b}}ϵ^{\underset{¯}{a}^{}\underset{¯}{b}^{}}M_{\alpha \beta }\frac{1}{2}ϵ_{\alpha \beta }(t^{\underset{¯}{a}\underset{¯}{b}}ϵ^{\underset{¯}{a}^{}\underset{¯}{b}^{}}+ϵ^{\underset{¯}{a}\underset{¯}{b}}t^{\underset{¯}{a}^{}\underset{¯}{b}^{}}2ϵ^{\underset{¯}{a}\underset{¯}{b}}ϵ^{\underset{¯}{a}^{}\underset{¯}{b}^{}}D)$$ (13) where the $`SU(2)`$ generators $`t`$ commute with the supersymmetry ones as follows: $$[t^{\underset{¯}{a}\underset{¯}{b}},Q^{\underset{¯}{c}\underset{¯}{c}^{}}]=i(ϵ^{\underset{¯}{c}\underset{¯}{a}}Q^{\underset{¯}{b}\underset{¯}{c}^{}}+ϵ^{\underset{¯}{c}\underset{¯}{b}}Q^{\underset{¯}{a}\underset{¯}{c}^{}}),[t^{\underset{¯}{a}^{}\underset{¯}{b}^{}},Q^{\underset{¯}{c}\underset{¯}{c}^{}}]=i(ϵ^{\underset{¯}{c}^{}\underset{¯}{a}^{}}Q^{\underset{¯}{c}\underset{¯}{b}^{}}+ϵ^{\underset{¯}{c}^{}\underset{¯}{b}^{}}Q^{\underset{¯}{c}\underset{¯}{a}^{}}).$$ (14) In this notation the two remaining charges are given by $$H_3=t^{\underset{¯}{1}\underset{¯}{2}},H_4=t^{\underset{¯}{1}^{}\underset{¯}{2}^{}},$$ (15) and by denoting $`\underset{¯}{1}[+],\underset{¯}{2}[]`$ and $`\underset{¯}{1}^{}\{+\},\underset{¯}{2}^{}\{\}`$, we find $$[H_3,Q^{[\pm ]\{\pm \}}]=[H_4,Q^{[\pm ]\{\pm \}}]=\pm iQ^{[\pm ]\{\pm \}}.$$ (16) The two relevant projections of eq. (2) now are $$\{Q_\alpha ^{[+]\{+\}},S_\beta ^{[]\{\}}\}=\frac{1}{2}M_{\alpha \beta }+ϵ_{\alpha \beta }(D\frac{1}{2}H_3\frac{1}{2}H_4),$$ (17) $$\{Q_\alpha ^{[+]\{\}},S_\beta ^{[]\{+\}}\}=\frac{1}{2}M_{\alpha \beta }ϵ_{\alpha \beta }(D\frac{1}{2}H_3+\frac{1}{2}H_4).$$ (18) Besides the four $`SO(2)`$ charges, the algebra of $`SO(8)`$ contains $`284=24`$ generators which can be arranged into 12 “step-up” operators (positive roots): $$\{𝒯\}_+=\{\begin{array}{c}T^{++(++)},T^{++()},T^{++[\pm ]\{\pm \}};\hfill \\ T^{(++)[\pm ]\{\pm \}};\hfill \\ T^{[++]}T^{[+]\{+\}[+]\{\}},T^{\{++\}}T^{[+]\{+\}[]\{+\}}\hfill \end{array}$$ (19) and their complex conjugates (negative roots). Among them only 4 (= rank of $`SO(8)`$) are independent (simple roots), namely, $`T^{[++]},T^{\{++\}},T^{++()},T^{(++)[]\{\}}`$. Above we have given the decomposition of two of the basic representations of $`SO(8)`$ under the particular embedding of $`[SO(2)]^4`$ that we are using here. These are the $`8_v`$ (the supersymmetry generators $`Q^i`$) and the adjoint $`28`$ (the $`SO(8)`$ generators $`T^{ij}`$). For future reference we also give the decomposition of the two spinor representations, $`8_s`$ ($`\varphi ^a`$, $`a=1,\mathrm{},8`$) and $`8_c`$ ($`\psi ^{\dot{a}}`$, $`\dot{a}=1,\mathrm{},8`$): $`\varphi ^a`$ $``$ $`\varphi ^{+(+)[\pm ]},\varphi ^{()[\pm ]},\varphi ^{+()\{\pm \}},\varphi ^{(+)\{\pm \}}`$ (20) $`\sigma ^{\dot{a}}`$ $``$ $`\sigma ^{+(+)\{\pm \}},\sigma ^{()\{\pm \}},\sigma ^{+()[\pm ]},\sigma ^{(+)[\pm ]}`$ (21) This has been obtained by successive reductions: $`SO(8)SO(2)\times SO(6)U(1)\times SU(4)[SO(2)]^2\times SO(4)[U(1)]^2\times SU(2)\times SU(2)[SO(2)]^4[U(1)]^4`$. Now we turn to the discussion of the representations of $`OSp(8/4,)`$. Let us denote a generic (quasi primary) superconformal field of the $`OSp(8/4,)`$ algebra by the quantum numbers of its HWS: $$𝒟(\mathrm{},J;d_1,d_2,d_3,d_4)$$ (22) where $`\mathrm{}`$ is the conformal dimension, $`J`$ is the Lorentz spin and $`d_1,d_2,d_3,d_4`$ are the Dynkin labels (see, e.g., ) of the $`SO(8)`$ R symmetry. In fact, in our scheme the natural labels are the four charges $`q_1,q_2,q_3,q_4`$ (the eigenvalues of $`H_1,\mathrm{},H_4`$). So, we can alternatively denote the HWS $`\|\mathrm{},J,q_i`$. The Dynkin labels $`[d_1,d_2,d_3,d_4]`$ are related to the $`U(1)`$ charges $`(q_1,q_2,q_3,q_4)`$ as follows: $$d_1=\frac{1}{2}(q_1q_2),d_2=\frac{1}{2}(q_2q_3q_4),d_3=q_3,d_4=q_4.$$ (23) The above relations can be most easily derived <sup>2</sup><sup>2</sup>2We are grateful to L. Castellani for suggesting this to us. by comparing the Dynkin labels and the charges of the HWS of the following four irreps: $`8_v:[1,0,0,0](2,0,0,0)`$, $`28:[0,1,0,0](2,2,0,0)`$, $`8_s:[0,0,1,0](1,1,1,0)`$, $`8_c:[0,0,0,1](1,1,0,1)`$. Note that (23) implies restrictions on the allowed values of the charges of a HWS: $$q_1q_2=2n0,q_2q_3q_4=2k0,q_30,q_40.$$ (24) A general HWS is defined by a subset of generators of the algebra which annihilate it. These include all the conformal supersymmetry generators: $$S_\alpha ^i\|\mathrm{},J,q_i=0$$ (25) (and, consequently, the boosts $`K_\mu `$) as well as the $`SO(8)`$ “step-up” operators (19): $$\{𝒯\}_+\|\mathrm{},J,q_i=0.$$ (26) The second condition defines $`\|\mathrm{},J,q_i`$ as the HWS of a UIR of $`SO(8)`$. A similar condition ensures irreducibility under the Lorentz group. Further, $`\|\mathrm{},J,q_i`$ should be an eigenstate of the generators $`D,M^2,H_i`$ fixing its dimension $`\mathrm{}`$, spin $`J`$ and charges $`q_i`$. Now, what makes a multiplet “short” is the additional requirement that part of the supersymmetry charges $`Q_\alpha ^i`$ also annihilate the HWS. When choosing this subset of $`Q`$’s we have to make sure that it is compatible with the rest of the conditions and with the algebra (1)-(6). First of all, these $`Q`$’s must anticommute among themselves, otherwise the first of eqs. (1) will yield restrictions on the momentum $`P_\mu `$. Secondly, eq. (26) implies that they must form a closed algebra (a Cauchy-Riemann structure) with all the $`SO(8)`$ step-up operators $`\{𝒯\}_+`$. It is easy to see that such a subset can at most involve four supercharges. In the AdS language such multiplets are called 1/2 BPS ($`4=\frac{1}{2}\mathrm{\hspace{0.33em}8}`$ generators annihilate the HWS). There exist two possible choices: type I 1/2 BPS: $`Q^{++}\|\mathrm{},J,q_i=Q^{(++)}\|\mathrm{},J,q_i=Q^{[+]\{+\}}\|\mathrm{},J,q_i=`$ (27) $`Q^{[+]\{\}}\|\mathrm{},J,q_i=0`$ or type II 1/2 BPS: $`Q^{++}\|\mathrm{},J,q_i=Q^{(++)}\|\mathrm{},J,q_i=Q^{[+]\{+\}}\|\mathrm{},J,q_i=`$ (28) $`Q^{[]\{+\}}\|\mathrm{},J,q_i=0.`$ Finally, conditions (27) or (28) should be consistent with (25). Using the projections (9), (12), (17) and (18) of eq. (2), we obtain the following constraint on the charges, conformal weight and spin of the HWS <sup>3</sup><sup>3</sup>3Such relations have been known from the very beginning of supersymmetry, see .: $$\text{type I 1/2 BPS:}q_1=q_2=q_3=2\mathrm{},q_4=0,J=0;$$ (29) $$\text{type II 1/2 BPS:}q_1=q_2=q_4=2\mathrm{},q_3=0,J=0,$$ (30) where $`2\mathrm{}m`$ is a non-negative integer. Computing the Dynkin labels from (23), we can say that the 1/2 BPS multiplets above correspond to $$\text{type I 1/2 BPS:}𝒟(m/2,0;0,0,m,0);$$ (31) $$\text{type II 1/2 BPS:}𝒟(m/2,0;0,0,0,m).$$ (32) Besides the 1/2 BPS conditions there exist weaker shortening conditions. Thus, we can require that a subset of only three supercharges annihilate the HWS. Once again, the choice must be consistent with condition (26), and this gives only one possibility: $$\text{3/8 BPS:}Q^{++}\|\mathrm{},J,q_i=Q^{(++)}\|\mathrm{},J,q_i=Q^{[+]\{+\}}\|\mathrm{},J,q_i=0.$$ (33) This is a 3/8 BPS multiplet in the AdS language. This time the condition on the weight, spin and charges is $$q_1=q_2=q_3+q_4=2\mathrm{},J=0.$$ (34) Denoting $`q_3=m`$, $`q_4=n`$ where $`m,n`$ are non-negative integers and computing the Dynkin labels, we find that this type of multiplet corresponds to $$\text{3/8 BPS:}𝒟(1/2(m+n),0;0,0,m,n).$$ (35) The next step will be to take a subset of two supercharges compatible with (26), which is $$\text{1/4 BPS:}Q^{++}\|\mathrm{},J,q_i=Q^{(++)}\|\mathrm{},J,q_i=0.$$ (36) This is a 1/4 BPS multiplet in the AdS language. This time the condition is $$q_1=q_2=2\mathrm{},J=0,$$ (37) $`q_3`$ and $`q_4`$ being only restricted by (24). Denoting $`q_1=q_2=m+n+2k`$, $`q_3=m`$, $`q_4=n`$ where $`m,n,k`$ are non-negative integers, we find that this type of multiplet corresponds to $$\text{1/4 BPS:}𝒟(1/2(m+n)+k,0;0,k,m,n).$$ (38) Finally, the weakest shortening condition is obtained by retaining only one supercharge (the HWS among the eight projections of $`Q^i`$): $$\text{1/8 BPS:}Q^{++}\|\mathrm{},J,q_i=0.$$ (39) This is a 1/8 BPS multiplet in the AdS language. The condition in this case is $$q_1=2\mathrm{},J=0,$$ (40) $`q_2,q_3`$ and $`q_4`$ satisfying (24). Denoting $`q_1=m+n+2k+2l`$, $`q_2=m+n+2k`$, $`q_3=m`$, $`q_4=n`$ where $`m,n,k,l`$ are non-negative integers, we find $$\text{1/8 BPS:}𝒟(1/2(m+n)+k+l,0;l,k,m,n).$$ (41) This concludes our abstract analysis of the possible short representations of $`OSp(8/4,)`$. Note that we are not directly addressing the question of whether these representations are unitary or not. However, in the rest of the paper we shall show that all of them can be realized by tensoring two elementary building blocks, the so-called supersingleton representations. Since the latter are known to be UIR’s of $`OSp(8/4,)`$, this also answers the above question affirmatively. ## 3 Supersingletons Let us consider the simplest $`OSp(8/4,)`$ representations of the type (31) or (32). They are obtained by setting $`m=1`$, so they correspond to $`𝒟(1/2,0;0,0,1,0)`$ or $`𝒟(1/2,0;0,0,0,1)`$. Such representations are called “supersingletons” . Each of them is just a collection of 8 Dirac supermultiplets made up of “Di” and “Rac” singletons . We observe that in the framework of the AdS/CFT correspondence the supersingleton describes the microscopic degrees of freedom of an M-2 brane with the scalars being the coordinates transverse to the brane which are then in the $`8_v`$ of $`SO(8)`$. The existence of two distinct types of $`N=8`$ 3d supersingletons has first been noted in Ref. . Our task now will be to realize the supersingleton in $`N=8`$ 3d superspace. Consider first type I. Noting that the HWS in the multiplet $`𝒟(1/2,0;0,0,1,0)`$ has spin 0 and the Dynkin labels of the $`8_s`$ of $`SO(8)`$, we take a scalar superfield $`\mathrm{\Phi }_a(x^\mu ,\theta _i^\alpha )`$ carrying an external $`8_s`$ index $`a`$. The superfield $`\mathrm{\Phi }_a`$ is a reducible representation of $`N=8`$ Poincaré supersymmetry. This can be seen from the fact that the first fermion field in its decomposition, $$\mathrm{\Phi }_a(x^\mu ,\theta _i^\alpha )=\varphi _a(x)+\theta _i^\alpha \psi _{\alpha ia}(x)+\mathrm{},$$ (42) is reducible under $`SO(8)`$: $`\psi _{\alpha ia}8_v8_s=8_c56_s`$. The way to achieve irreducibility is to impose a constraint on the superfield which removes the $`56_s`$ part of $`\psi _{\alpha ia}`$: $$\text{type I:}D_\alpha ^i\mathrm{\Phi }_a=\frac{1}{8}\gamma _{a\dot{b}}^i\stackrel{~}{\gamma }_{\dot{b}c}^jD_\alpha ^j\mathrm{\Phi }_c.$$ (43) Here $`D_\alpha ^i`$ are the covariant spinor derivatives satisfying the supersymmetry algebra $$\{D_\alpha ^i,D_\beta ^j\}=2i\delta ^{ij}(\mathrm{\Gamma }^\mu )_{\alpha \beta }_\mu .$$ (44) The $`SO(8)`$ gamma matrices $`\gamma _{a\dot{b}}^i`$ and $`\stackrel{~}{\gamma }_{\dot{a}b}^i=(\gamma ^{iT})_{\dot{a}b}`$ satisfy the Clifford algebra relations $$\gamma _{a\dot{b}}^i\stackrel{~}{\gamma }_{\dot{b}c}^j+\gamma _{a\dot{b}}^j\stackrel{~}{\gamma }_{\dot{b}c}^i=2\delta ^{ij}\delta _{ac},\stackrel{~}{\gamma }_{\dot{a}b}^i\gamma _{b\dot{c}}^j+\stackrel{~}{\gamma }_{\dot{a}b}^j\gamma _{b\dot{c}}^i=2\delta ^{ij}\delta _{\dot{a}\dot{c}}.$$ (45) Using (44) one can show that the constraint (43) eliminates all the components of the superfield but two: $`\mathrm{\Phi }_a(x^\mu ,\theta _i^\alpha )`$ $`=`$ $`\varphi _a(x)+\theta _i^\alpha (\gamma _i)_{a\dot{b}}\psi _{\alpha \dot{b}}(x)`$ (46) $`+\theta _i^\alpha \theta _j^\beta (\gamma _{ij})_{ab}i_{\alpha \beta }\varphi _b`$ $`+\theta _i^\alpha \theta _i^\beta \theta _k^\gamma (\gamma _{ijk})_{a\dot{b}}i_{(\alpha \beta }\psi _{\gamma )\dot{b}}`$ $`+\theta _i^\alpha \theta _i^\beta \theta _k^\gamma \theta _l^\delta (\gamma _{ijkl})_{ab}_{(\alpha \beta }_{\gamma \delta )}\varphi _b`$ where $`_{\alpha \beta }=_{\beta \alpha }=(\mathrm{\Gamma }^\mu )_{\alpha \beta }_\mu `$ and $`\gamma _{ij\mathrm{}}`$ are the antisymmetrized products of the $`SO(8)`$ gamma matrices. In addition, the constraint (43) puts these fields on shell: $$\mathrm{}\varphi _a=0,^{\alpha \beta }\psi _{\beta \dot{a}}=0.$$ (47) Thus, the content of the constrained superfield is a massless multiplet of Poincaré supersymmetry consisting of a scalar in the $`8_s`$ and a spinor in the $`8_c`$ UIR’s of $`SO(8)`$. <sup>4</sup><sup>4</sup>4Superfield representations of other $`OSp(N/4)`$ have been considered in the literature . Note that the field equations (47) can be obtained from a supersymmetric action . Consequently, the physical fields $`\varphi _a`$ and $`\psi _{\alpha \dot{a}}`$ have canonical dimensions $`1/2`$ and $`1`$, respectively. This implies that the superfield $`\mathrm{\Phi }_a`$ has dimension $`1/2`$, in accord with the abstract representation $`𝒟(1/2,0;0,0,1,0)`$. Finally, the alternative supersingleton representation of type II can be realized in terms of a superfield $`\mathrm{\Sigma }_{\dot{a}}`$ carrying an $`8_c`$ external index and satisfying the constraint $$\text{type II:}D_\alpha ^i\mathrm{\Sigma }_{\dot{a}}=\frac{1}{8}\stackrel{~}{\gamma }_{\dot{a}b}^i\gamma _{b\dot{c}}^jD_\alpha ^j\mathrm{\Sigma }_{\dot{c}}.$$ (48) It describes a massless multiplet consisting of a scalar $`\sigma _{\dot{a}}(x)`$ and a spinor $`\chi _{\alpha a}(x)`$ in the $`8_c`$ and $`8_s`$, correspondingly. The problem we want to address now is how to tensor supersingletons. Doing it directly in terms of constrained superfields is quite difficult. Our alternative approach consists in first rewriting the constraints (43) or (48) as analyticity conditions in harmonic superspace, after which the tensor multiplication becomes straightforward. ## 4 The supersingletons as harmonic analytic superfields The harmonic space suitable for our purposes is given by the coset <sup>5</sup><sup>5</sup>5A formulation of the above multiplet in harmonic superspace has been proposed in Ref. (see also and for a general discussion of three-dimensional harmonic superspaces). The harmonic coset used in is $`Spin(8)/U(4)`$. Although the supersingleton itself does indeed live on this smaller coset, the residual symmetry $`U(4)`$ will turn out too big when we start tensoring different realizations of the supersingleton. For this reason we prefer from the very beginning to use the coset (49) with a minimal residual symmetry (see also for a discussion of this point). $$\frac{SO(8)}{[SO(2)]^4}\frac{Spin(8)}{[U(1)]^4}.$$ (49) This is a $`284=24`$-dimensional compact manifold. Instead of trying to introduce explicit coordinates on it, the harmonic method prescribes to use the entire matrices of the fundamental representation of the group to parametrize the coset. The complication in the case of $`SO(8)`$ is that one has three inequivalent fundamental representations, $`8_s,8_c,8_v`$. The solution to this problem has been found in Ref. . One introduces three sets of harmonic variables: $$u_a^A,w_{\dot{a}}^{\dot{A}},v_i^I$$ (50) where $`A`$, $`\dot{A}`$ and $`I`$ denote the decompositions of an $`8_s`$, $`8_c`$ and $`8_v`$ index, correspondingly, into sets of four $`U(1)`$ charges, according to the coset denominator $`[U(1)]^4`$ (see Section 2 for details). Each of the $`8\times 8`$ real matrices (50) is a matrix of the corresponding representation of $`SO(8)Spin(8)`$. This implies that all of them are orthogonal matrices (this is a peculiarity of $`SO(8)`$ due to triality): $$u_a^Au_a^B=\delta ^{AB},w_{\dot{a}}^{\dot{A}}w_{\dot{a}}^{\dot{B}}=\delta ^{\dot{A}\dot{B}},v_i^Iv_i^J=\delta ^{IJ}$$ (51) (and similarly with small and capital indices interchanged). These matrices supply three copies of the group space (i.e., three sets of 28 real variables each), and we only need one to parametrize the coset (49). The condition which identifies the three sets of harmonic variables is $$u_a^A(\gamma ^I)_{A\dot{A}}w_{\dot{a}}^{\dot{A}}=v_i^I(\gamma ^i)_{a\dot{a}}.$$ (52) This relation just expresses the transformation properties of the gamma matrices under $`SO(8)`$. The reader can convince him(her)self that the conditions (51), (52) leave just one set of 28 independent parameters by taking the infinitesimal form of the above matrices. Note that eq. (52) can be viewed as the expression of the vector harmonics in terms of the two types of spinor ones. Therefore we shall choose $`u,w`$ as our harmonic variables. <sup>6</sup><sup>6</sup>6Although each of the three sets of harmonic variables depends on the same 28 parameters, we need at least two sets to be able to reproduce all possible representations of $`SO(8)`$. The idea of the harmonic description of the coset (49) is to consider harmonic functions defined as functions of the above sets of variables modulo transformations of $`[U(1)]^4`$. In other words, a harmonic function always carries a set of four $`U(1)`$ charges. These functions are then given by their “harmonic expansions” in terms of all the products of harmonic variables having the same charges. Take, for instance, the function $`\varphi ^{+(+)[+]}(u,w)`$ $`=`$ $`\varphi _au_a^{+(+)[+]}`$ (53) $`+\varphi _{abc}u_a^{+(+)[+]}u_b^{+(+)[+]}u_c^{()[]}`$ $`+\varphi _{a\dot{b}\dot{c}}u_a^{+(+)[+]}w_{\dot{b}}^{+(+)\{+\}}w_{\dot{c}}^{()\{\}}+\mathrm{}.`$ Although the harmonic function only transforms under $`[U(1)]^4`$, the coefficients in its expansion are representations of $`SO(8)Spin(8)`$. Thus, a harmonic function is a collection of an infinite set of irreps of $`SO(8)`$. In order to make the harmonic functions irreducible we have to impose differential constraints on them. To this end we introduce harmonic derivatives (the covariant derivatives on the coset (49)): $$D^{IJ}=u_a^A(\gamma ^{IJ})^{AB}\frac{}{u_a^B}+w_{\dot{a}}^{\dot{A}}(\gamma ^{IJ})^{\dot{A}\dot{B}}\frac{}{w_{\dot{a}}^{\dot{B}}}+v_i^{[I}\frac{}{v_i^{J]}}.$$ (54) They respect the algebraic relations (51), (52) among the harmonic variables. Moreover, these derivatives form the algebra of $`SO(8)`$ realized on the $`[U(1)]^4`$ projected indices $`A,\dot{A},I`$ of the harmonics. Four of them just count the four $`U(1)`$ charges, i.e. the harmonic functions are their eigenfunctions: $$H_nf^{(q_1,q_2,q_3,q_4)}(u,w)=q_nf^{(q_1,q_2,q_3,q_4)}(u,w),n=1,2,3,4.$$ (55) The remaining 24 ones are the true covariant derivatives on the coset. In our complex $`[U(1)]^4`$ notation these are $$\{𝒟\}_+=\{\begin{array}{c}D^{++(++)},D^{++()},D^{++[\pm ]\{\pm \}};\hfill \\ D^{(++)[\pm ]\{\pm \}};\hfill \\ D^{[++]}D^{[+]\{+\}[+]\{\}},D^{\{++\}}D^{[+]\{+\}[]\{+\}}\hfill \end{array}$$ (56) and their complex conjugates. It is clear that the 12 derivatives (56) correspond to the step-up operators of $`SO(8)`$, see (19). Therefore we can make a harmonic function irreducible by demanding that all of the derivatives (56) annihilate it. In other words, this differential condition reduces the harmonic function to a polynomial corresponding to a highest weight of an $`SO(8)`$ irrep. For example, the constraint $$\{𝒟\}_+\varphi ^{+(+)[+]}(u,w)=0\varphi ^{+(+)[+]}(u,w)=\varphi _au_a^{+(+)[+]},$$ (57) reduces the function (53) to a $`8_s`$. This can easily be generalized to any function of the type (55) satisfying the constraint $$\{𝒟\}_+f^{(q_1,q_2,q_3,q_4)}(u,w)=0.$$ (58) This is the defining condition of the HWS of a UIR of $`SO(8)`$ given by the Dynkin labels from eq. (23). The function satisfying (58) is thus reduced to a polynomial of the harmonic variables: $`f^{(q_1,q_2,q_3,q_4)}(u,w)=f^{(2d_1+2d_2+d_3+d_4,2d_2+d_3+d_4,d_3,d_4)}(u,w)=`$ (59) $`f_{a\mathrm{}b\mathrm{}c\mathrm{}\dot{d}\mathrm{}}(u_a^{+(+)[+]})^{d_2+d_3}(u_b^{+(+)[]})^{d_2}(u_c^{+()\{\}})^{d_1}(w_{\dot{d}}^{+(+)\{+\}})^{d_1+d_4}.`$ Concluding the discussion of the harmonic coset (49) we can say that if one introduces complex coordinates on it, the conditions (58) take the form of (covariant) analyticity conditions. For this reason we can call eqs. (58) “harmonic analyticity” conditions. The purpose of introducing harmonic variables is to be able to project the supersingleton defining constraint (43) (or (48)) in an $`SO(8)`$ covariant way. This means to convert the indices $`i`$ and $`a`$ into $`U(1)`$ charges with the help of the corresponding harmonics: $`D_\alpha ^iD_\alpha ^I=v_i^ID_\alpha ^i`$ and $`\mathrm{\Phi }_a\mathrm{\Phi }^A=u_a^A\mathrm{\Phi }_a`$. Then, using the relation (52) it is easy to show that, e.g., the projection $`\mathrm{\Phi }^{+(+)[+]}`$ satisfies the following constraints: $$D^{++}\mathrm{\Phi }^{+(+)[+]}=D^{(++)}\mathrm{\Phi }^{+(+)[+]}=D^{[+]\{\pm \}}\mathrm{\Phi }^{+(+)[+]}=0.$$ (60) We see that half of the spinor derivatives annihilate the superfield $`\mathrm{\Phi }^{+(+)[+]}`$. This is the superspace realization of the 1/2 BPS shortening condition (27). Since these spinor derivatives anticommute among themselves (as follows from (44) after the appropriate projections), there exists a basis in superspace where $`\mathrm{\Phi }^{+(+)[+]}`$ becomes just a function of half of the odd variables as well as of the harmonic variables: $$\text{type I:}\mathrm{\Phi }^{+(+)[+]}=\mathrm{\Phi }^{+(+)[+]}(x_A,\theta ^{++},\theta ^{(++)},\theta ^{[+]\{\pm \}},u,w)$$ (61) where $$x_{A\alpha \beta }=x_{\alpha \beta }+i\theta _{(\alpha }^{++}\theta _{\beta )}^{}+i\theta _{(\alpha }^{(++)}\theta _{\beta )}^{()}+i\theta _{(\alpha }^{[+]\{+\}}\theta _{\beta )}^{[]\{\}}+i\theta _{(\alpha }^{[+]\{\}}\theta _{\beta )}^{[]\{+\}}.$$ (62) We can say that $`\mathrm{\Phi }^{+(+)[+]}`$ is a “Grassmann analytic” or a “short” superfield. So far eqs. (60) have been derived as a corollary of the defining constraint (43). In order to make the latter equivalent to the former we have to eliminate the harmonic dependence in the superfield (61). This is done by imposing another set of constraints, namely, the harmonic analyticity conditions (58): $$\{𝒟\}_+\mathrm{\Phi }^{+(+)[+]}(x,\theta ^{++},\theta ^{(++)},\theta ^{[+]\{\pm \}},u,w)=0.$$ (63) Note that these new constraints are compatible with (60) since the two sets of derivatives form a closed algebra (a Cauchy-Riemann structure in the terminology of Ref. ). It should be stressed that eq. (63) now has implications other than just restricting the harmonic dependence. The reason is that in the superspace basis (62) where Grassmann analyticity becomes manifest some of the harmonic derivatives from the set $`\{𝒟\}_+`$ acquire torsion terms, e.g., $`D^{++(++)}=_{u,w}^{++(++)}+i\theta ^{++}\mathrm{\Gamma }^\mu \theta ^{(++)}_\mu `$, $`D^{++[+]\{\pm \}}=_{u,w}^{++[+]\{\pm \}}+i\theta ^{++}\mathrm{\Gamma }^\mu \theta ^{[+]\{\pm \}}_\mu `$, etc. This yields space-time derivative constraints on the components of the superfield $`\mathrm{\Phi }^{+(+)[+]}`$. All this amounts to $`\mathrm{\Phi }^{+(+)[+]}`$ becoming “ultrashort”: $`\mathrm{\Phi }^{+(+)[+]}`$ $`=`$ $`u_a^{+(+)[+]}\varphi _a(x)`$ (64) $`+(\theta ^{[+]\{\}\alpha }w_{\dot{a}}^{+(+)\{+\}}\theta ^{[+]\{+\}\alpha }w_{\dot{a}}^{+(+)\{\}}`$ $`\theta ^{++\alpha }w_{\dot{a}}^{(+)[+]}\theta ^{(++)\alpha }w_{\dot{a}}^{+()[+]})\psi _{\dot{a}\alpha }(x)`$ $`+\text{ derivative terms}`$ where the fields are massless. In this way we recover the content (46), (47) of the ordinary constrained superfield describing the supersingleton multiplet. It is instructive to comment on the structure of the two terms in eq. (64). The first one is the component at level 0 in the $`\theta `$ expansion. It is a harmonic function of the type (57), i.e. a harmonic-projected $`8_s`$. The situation at level 1 is more complicated. Originally, one finds a collection of spinor fields with a variety of charges. In order to find out which one among them is the HWS of an $`SO(8)`$ representation, we have to look at the accompanying $`\theta `$’s. It is easy to see that $`\theta ^{[+]\{\}}`$ can serve as a starting point for obtaining the rest by successive applications of the harmonic derivatives $`\{𝒟\}_+`$ (the step-up operators of $`SO(8)`$): $$\theta ^{[+]\{\}}\stackrel{D^{\{++\}}}{}\theta ^{[+]\{+\}}\stackrel{D^{(++)[]\{\}}}{}\theta ^{(++)}\stackrel{D^{++()}}{}\theta ^{++}.$$ (65) At the same time, $`\theta ^{[+]\{\}}`$ cannot be obtained from any other of the projections available in the Grassmann analytic superspace. As a consequence, the harmonic analyticity condition (63) mixes up the corresponding spinor fields (coefficients at level 1 in the $`\theta `$ expansion), with the exception of the one in the term $`\theta ^{[+]\{\}\alpha }\psi _\alpha ^{+(+)\{+\}}(x,u,w)`$. The latter must satisfy the condition $`\{𝒟\}_+\psi _\alpha ^{+(+)\{+\}}=0`$. This means that we are dealing with the HWS of the representation $`(1,1,0,1)[0,0,0,1]`$, i.e. with a $`8_c`$. The remaining level 1 coefficients are related to this HWS by harmonic equations like, e.g., $`D^{\{++\}}\psi _\alpha ^{+(+)\{\}}=\psi _\alpha ^{+(+)\{+\}}`$, etc. In other words, they correspond to different projections (“lower weights”) of this $`8_c`$. The same argument explains why there are no new fields beyond level 1. Indeed, among all the level 2 $`\theta `$ structures we find two which cannot be obtained by acting with the step-up operators on any other structure: $$\theta ^{[+]\{\}\alpha }\theta _\alpha ^{[+]\{\}}A^{(1,1,1,2)},\theta ^{[+]\{\}\alpha }\theta ^{[+]\{+\}\beta }B_{(\alpha \beta )}^{(1,1,1,0)}$$ (66) corresponding to a scalar and a vector fields. Now, harmonic analyticity again implies that these fields should be highest weights of $`SO(8)`$ irreps, but their charges do not satisfy the restrictions (24). The conclusion is that there are no such independent fields in the expansion of the analytic superfield $`\mathrm{\Phi }^{+(+)[+]}`$ (more precisely, $`A^{(1,1,1,2)}=0`$ and $`B_{(\alpha \beta )}^{(1,1,1,0)}=i_{\alpha \beta }\varphi _au_a^{+(+)[]}`$; such terms are denoted as “derivative terms” in (64)). In conclusion we note that the alternative form of the supersingleton (48) is described by the superfield $$\text{type II:}\mathrm{\Sigma }^{+(+)\{+\}}(\theta ^{++},\theta ^{(++)},\theta ^{[\pm ]\{+\}})$$ (67) satisfying the same harmonic constraints (63) but depending on a different set of four odd variables. Also, the charges and Dynkin labels of the first component are those of an $`8_c`$ instead of $`8_s`$. This is the superspace realization of the 1/2 BPS shortening condition (28). ## 5 Short multiplets as supersingleton “composite operators” In the preceding section, with the help of the harmonic variables, we have been able to equivalently rewrite the supersingleton as an ultrashort superfield satisfying both conditions of Grassmann (eq. (60) or eq. (67)) and harmonic (eq. (63)) analyticity. The main advantage of this new analytic form of the supersingleton is the possibility to tensor copies of it in a straightforward way and thus to obtain series of short composite multiplets. As we shall show in this section, this procedure allows us to realize all the abstract short $`OSp(8/4,)`$ multiplets of Section 2. We observe that in the AdS/CFT correspondence the supersingleton multiplet describing the dynamics of many M-2 branes is endowed with an internal symmetry index and composite operators are further restricted to be singlets under the invariance group . The simplest example of a tensor product is obtained by taking $`p`$ identical copies of type I supersingletons, $`(\mathrm{\Phi }^{+(+)[+]})^p`$. Clearly, it satisfies the same constraints of Grassmann and harmonic analyticity. However, the latter is not as strong as before. The reason is that the external charges of the superfield have changed, and the consequences of harmonic analyticity strongly depend on the charges, as the argument at the end of the preceding section has shown. So, for generic $`p4`$ the $`\theta `$ expansion goes up to the maximal level 8: $`(\mathrm{\Phi }^{+(+)[+]})^p`$ $`=`$ $`\varphi ^{[0,0,p,0]}`$ (68) $`+`$ $`\theta ^{[+]\{\}\alpha }\psi _\alpha ^{[0,0,p1,1]}+\mathrm{}`$ $`+`$ $`(\theta ^{[+]\{\}})^2A^{[0,0,p2,2]}+\mathrm{}`$ $`+`$ $`\theta ^{[+]\{\}\alpha }\theta ^{[+]\{+\}\beta }B_{(\alpha \beta )}^{[0,1,p2,0]}+\mathrm{}`$ $`+`$ $`(\theta ^{[+]\{\}})^2\theta ^{[+]\{+\}\alpha }\chi _\alpha ^{[0,1,p3,1]}+\mathrm{}`$ $`+`$ $`\theta ^{[+]\{\}\alpha }\theta ^{[+]\{+\}\beta }\theta ^{(++)\gamma }\rho _{(\alpha \beta \gamma )}^{[1,0,p2,0]}+\mathrm{}`$ $`+`$ $`(\theta ^{[+]\{\}})^2(\theta ^{[+]\{+\}})^2C^{[0,2,p4,0]}+\mathrm{}`$ $`+`$ $`(\theta ^{[+]\{\}})^2\theta ^{[+]\{+\}\alpha }\theta ^{(++)\beta }D_{(\alpha \beta )}^{[1,0,p3,1]}+\mathrm{}`$ $`+`$ $`\theta ^{[+]\{\}\alpha }\theta ^{[+]\{+\}\beta }\theta ^{(++)\gamma }\theta ^{++\delta }E_{(\alpha \beta \gamma \delta )}^{[0,0,p2,0]}+\mathrm{}`$ $`+`$ $`(\theta ^{[+]\{\}})^2(\theta ^{[+]\{+\}})^2\theta ^{(++)\alpha }\sigma _\alpha ^{[1,1,p4,0]}+\mathrm{}`$ $`+`$ $`(\theta ^{[+]\{\}})^2\theta ^{[+]\{+\}\alpha }\theta ^{(++)\beta }\theta ^{++\gamma }\omega _{(\alpha \beta \gamma )}^{[0,0,p3,1]}+\mathrm{}`$ $`+`$ $`(\theta ^{[+]\{\}})^2(\theta ^{[+]\{+\}})^2(\theta ^{(++)})^2F^{[2,0,p4,0]}+\mathrm{}`$ $`+`$ $`(\theta ^{[+]\{\}})^2(\theta ^{[+]\{+\}})^2\theta ^{(++)\alpha }\theta ^{++\beta }G_{(\alpha \beta )}^{[0,1,p4,0]}+\mathrm{}`$ $`+`$ $`(\theta ^{[+]\{\}})^2(\theta ^{[+]\{+\}})^2(\theta ^{(++)})^2\theta ^{++\alpha }\tau _\alpha ^{[1,0,p4,0]}+\mathrm{}`$ $`+`$ $`(\theta ^{[+]\{\}})^2(\theta ^{[+]\{+\}})^2(\theta ^{(++)})^2(\theta ^{++})^2H^{[0,0,p4,0]}+\mathrm{}`$ $`+`$ derivative terms Here we have shown only the leading term at each level and of each Lorentz structure. This is the term whose coefficient is the HWS of an $`SO(8)`$ irrep. The other terms of the same type contain different harmonic projections of the same component field. Further, instead of the charges we have directly indicated the corresponding Dynkin labels of each component field. Note that the level in the expansion also determines the conformal dimension of the components (given the fact that the dimension of the first component is $`p/2`$ and that of a $`\theta `$ is $`1/2`$). We see that $`(\mathrm{\Phi }^{+(+)[+]})^p`$ is a short superfield (it depends on half of the odd variables) of the type (31), but not an ultrashort one, unlike the supersingleton itself (the case $`p=1`$). Still, for $`p=2,3`$ certain terms in the expansion (68) are absent if conditions (24) are not satisfied. In addition, for $`p=2`$ one finds conservation conditions for the fields of spins 2, 3/2 and 1, $`^{\alpha \beta }E_{(\alpha \beta \gamma \delta )}^{[0,0,0,0]}=^{\alpha \beta }\rho _{(\alpha \beta \gamma )}^{[1,0,0,0]}=^{\alpha \beta }B_{(\alpha \beta )}^{[0,1,0,0]}=0`$. This is most easily seen for the top spin 2 which is the only $`SO(8)`$ singlet in the expansion and hence its divergence cannot be matched by any other component. The expansion (68) reproduces (up to triality) the content of the short multiplets of $`OSp(8/4,)`$ found in Refs. , . Further short multiplets can be obtained by tensoring different analytic superfields describing the type I supersingleton. The point is that in Section 4 we chose a particular projection of the defining constraint (43) which lead to the analytic superfield $`\mathrm{\Phi }^{+(+)[+]}`$. In fact, we could have done this in a variety of ways, each time obtaining superfields depending on different halves of the total number of odd variables. If we decide to always leave out the lowest weight $`\theta ^{}`$ in the $`8_v`$ formed by the $`\theta `$’s, we can have four (as many as the rank of $`SO(8)`$) distinct but equivalent analytic descriptions of the type I supersingleton: $`\mathrm{\Phi }^{+(+)[+]}(\theta ^{++},\theta ^{(++)},\theta ^{[+]\{+\}},\theta ^{[+]\{\}}),`$ $`\mathrm{\Phi }^{+(+)[]}(\theta ^{++},\theta ^{(++)},\theta ^{[]\{+\}},\theta ^{[]\{\}}),`$ $`\mathrm{\Phi }^{+()\{+\}}(\theta ^{++},\theta ^{()},\theta ^{[+]\{+\}},\theta ^{[]\{+\}}),`$ $`\mathrm{\Phi }^{+()\{\}}(\theta ^{++},\theta ^{()},\theta ^{[+]\{\}},\theta ^{[]\{\}}).`$ (69) Then we can tensor them in the following way: $`(\mathrm{\Phi }^{+(+)[+]})^{p+q+r+s}(\mathrm{\Phi }^{+(+)[]})^{q+r+s}(\mathrm{\Phi }^{+()\{+\}})^{r+s}(\mathrm{\Phi }^{+()\{\}})^s`$ $`=\varphi ^{[r+2s,q,p,r]}+\mathrm{}`$ $`+\theta _{\alpha _1}^{[+]\{\}}\theta _{\alpha _2}^{[+]\{+\}}\theta _{\alpha _3}^{(++)}\theta _{\alpha _4}^{++}A^{[r+2s,q,p2,r](\alpha _1\mathrm{}\alpha _4)}+\mathrm{}`$ (70) $`+\theta _{\alpha _1}^{[+]\{\}}\theta _{\alpha _2}^{[+]\{+\}}\theta _{\alpha _3}^{(++)}\theta _{\alpha _4}^{++}\theta _{\alpha _5}^{[]\{+\}}\theta _{\alpha _6}^{[]\{\}}B^{[r+2s,q1,p,r](\alpha _1\mathrm{}\alpha _6)}+\mathrm{}`$ $`+\theta _{\alpha _1}^{[+]\{\}}\theta _{\alpha _2}^{[+]\{+\}}\theta _{\alpha _3}^{(++)}\theta _{\alpha _4}^{++}\theta _{\alpha _5}^{[]\{+\}}\theta _{\alpha _6}^{[]\{\}}\theta _{\alpha _7}^{()}\chi ^{[r+2s1,q,p,r](\alpha _1\mathrm{}\alpha _7)}+\mathrm{}`$ Here we have shown the first component which belongs to the $`SO(8)`$ UIR $`[r+2s,q,p,r]`$ and has conformal dimension $`\mathrm{}=\frac{1}{2}(p+2q+3r+4s)`$ (this follows from the fact that the basic supersingleton has dimension $`1/2`$). In (70) one can also see the top spin of each particular series: $`J_{top}=2`$ if $`q=r=s=0`$, $`J_{top}=3`$ if $`r=s=0`$ or $`J_{top}=7/2`$ if either $`r0`$ or $`s0`$. The dimension of the top spin is $`\mathrm{}[J_{top}]=\frac{1}{2}(p+2q+3r+4s)+J_{top}`$ (since each $`\theta `$ carries dimension $`1/2`$). Note the absence of a series with top spin $`J=5/2`$: the reason is that the tensor product of the different realizations (69) of the type I supersingleton can depend on 4, 6 or 7 $`\theta `$’s but not on 5. The above result can be summarized as follows. By considering composite operators made out of type I supersingletons we have constructed the following series of $`OSp(8/4,)`$ UIR’s exhibiting $`1/8`$, $`1/4`$ or $`1/2`$ BPS shortening: $`{\displaystyle \frac{1}{8}}\text{ BPS:}`$ $`𝒟(d_1+d_2+{\displaystyle \frac{1}{2}}(d_3+d_4),0;d_1,d_2,d_3,d_4),d_1d_4=2s0;`$ $`{\displaystyle \frac{1}{4}}\text{ BPS:}`$ $`𝒟(d_2+{\displaystyle \frac{1}{2}}d_3,0;0,d_2,d_3,0);`$ (71) $`{\displaystyle \frac{1}{2}}\text{ BPS:}`$ $`𝒟({\displaystyle \frac{1}{2}}d_3,0;0,0,d_3,0).`$ We see that tensoring only one type of supersingletons cannot reproduce the general result of Section 2 for all possible short multiplets. Most notably, in (71) there is no 3/8 series. The latter can be obtained by mixing the two types of supersingletons: $$[\mathrm{\Phi }^{+(+)[+]}(\theta ^{++},\theta ^{(++)},\theta ^{[+]\{\pm \}})]^{p+q}[\mathrm{\Sigma }^{+(+)\{+\}}(\theta ^{++},\theta ^{(++)},\theta ^{[\pm ]\{+\}})]^q,$$ (72) or the same with $`\mathrm{\Phi }`$ and $`\mathrm{\Sigma }`$ exchanged. Counting the charges and the dimension, we find exact matching with the series (35). Further, mixing two realizations of type I and one of type II supersingletons, we can construct the 1/4 series $$[\mathrm{\Phi }^{+(+)[+]}]^{m+k}[\mathrm{\Phi }^{+(+)[]}]^k[\mathrm{\Sigma }^{+(+)\{+\}}]^n$$ (73) which corresponds to (38). Finally, the full 1/8 series (41) (i.e., without the restriction $`d_1d_4=2s`$ in (71)) can be obtained in a variety of ways. ## 6 Conclusions In this paper we have analyzed all short highest weight UIR’s of the $`OSp(8/4,)`$ superalgebra whose HWS’s are annihilated by part of the super-Poincaré odd generators. In the field theory language, highest weight reps correspond to conformal quasi primary superfields. Short reps correspond to superfields which do not depend on some of the odd coordinates, a concept generalizing the notion of chiral superfields of $`N=1`$ 4d field theories. The number of distinct possibilities have been shown to correspond to different BPS conditions on the HWS. When the algebra is interpreted on the $`AdS_4`$ bulk, for which the 3d superconformal field theory corresponds to the boundary M-2 brane dynamics, these states appear as BPS massive excitations, such as K-K states or AdS black holes, of M-theory on $`AdS_4\times S^7`$. Since in M-theory there is only one type of supersingleton related to the M-2 brane transverse coordinates , according to our analysis massive states cannot be 3/8 BPS saturated, exactly as it happens in M-theory on $`M^4\times T^7`$. Indeed, the missing solution was also noticed in Ref. by studying $`AdS_4`$ black holes in gauged $`N=8`$ supergravity. Curiously, in the ungauged theory, which is in some sense the flat limit of the former, the 3/8 BPS states are forbidden by the underlying $`E_{7(7)}`$ symmetry of $`N=8`$ supergravity . ## Acknowledgements We would like to thank F. Delduc, M. Günaydin, L. Castellani, A. Sciarrino and P. Sorba for enlightening discussions. E.S. is grateful to the TH Division of CERN for its kind hospitality. The work of S.F. has been supported in part by the European Commission TMR programme ERBFMRX-CT96-0045 (Laboratori Nazionali di Frascati, INFN) and by DOE grant DE-FG03-91ER40662, Task C.
warning/0003/math-ph0003037.html	ar5iv	text	# Combinatorial Identities from the Spectral Theory of Quantum Graphs ## I Introduction In the present paper we show that some questions arising in the study of spectral correlations for quantum graphs can be cast as combinatorial problems. Solving these problems for a particular system, we discovered the following novel combinatorial identities: (i) Let $`n,q`$ be arbitrary integers with $`1q<n`$ and $$F_{\nu ,\nu ^{}}(n,q)=\frac{(n1)n}{2}\frac{(1)^{\nu +\nu ^{}}}{\nu \nu ^{}}\left(\genfrac{}{}{0pt}{}{n}{\nu +\nu ^{}}\right)^1\left(\genfrac{}{}{0pt}{}{q1}{\nu 1}\right)\left(\genfrac{}{}{0pt}{}{q1}{\nu ^{}1}\right)\left(\genfrac{}{}{0pt}{}{nq1}{\nu 1}\right)\left(\genfrac{}{}{0pt}{}{nq1}{\nu ^{}1}\right).$$ Then $$S(n,q)=\underset{\nu ,\nu ^{}=1}{\overset{\mathrm{min}(q,nq)}{}}F_{\nu ,\nu ^{}}(n,q)=1.$$ (1) (ii) Let $`s,t`$ be arbitrary positive integers and $$𝒩(s,t)=\underset{\nu =1}{\overset{\mathrm{min}(s,t)}{}}(1)^{t\nu }\left(\genfrac{}{}{0pt}{}{t}{\nu }\right)\left(\genfrac{}{}{0pt}{}{s1}{\nu 1}\right)=(1)^{s+t}\left(\genfrac{}{}{0pt}{}{s+t1}{s}\right)^{1/2}P_{s+t1,s}(t),$$ (2) where $`P_{N,k}(x)`$ are the Kravtchouk polynomials to be defined in Eq. (55). Further, let $`x,y`$ be complex with $`\|x\|,\|y\|<1/\sqrt{2}`$. Then we have the generating functions $$G_1(x)=\underset{s,t=1}{\overset{\mathrm{}}{}}𝒩^2(s,t)x^{s+t}=\frac{x}{2x1}\left(\frac{1}{\sqrt{4x^2+1}}\frac{1}{1x}\right),$$ (3) $$G_2(x)=\underset{s,t=1}{\overset{\mathrm{}}{}}𝒩(s,t)𝒩(t,s)x^{s+t}=\frac{1}{2}\frac{4x^2+2x+1}{(2x+1)\sqrt{4x^2+1}}\frac{1}{2}$$ (4) and $$g(x,y)=\underset{s,t=1}{\overset{\mathrm{}}{}}𝒩(s,t)x^sy^t\frac{xy}{(1+y)(1x+y2xy)}.$$ (5) (iii) Let $`m`$ be any positive integer. Then $$4m^2\underset{q=1}{\overset{2m1}{}}\left(\frac{𝒩(s,t)}{q}\right)^2=2^{2m+1}+(1)^m\left(\genfrac{}{}{0pt}{}{2m}{m}\right)2$$ (6) and $$(2m+1)^2\underset{q=1}{\overset{2m}{}}\left(\frac{𝒩(s,t)}{q}\right)^2=2^{2m+2}2(1)^m\left(\genfrac{}{}{0pt}{}{2m}{m}\right)2.$$ (7) (iv) Let $`0qn`$, and define $$A(n,q)=\frac{1}{\sqrt{2^n}}\{\begin{array}{cc}1\hfill & \mathrm{for}q=0,n\hfill \\ (1)^q(n/q)𝒩(nq,q)\hfill & \mathrm{for}\mathrm{\hspace{0.33em}0}<q<n.\hfill \end{array}$$ (8) Then for any positive integers $`0\kappa \nu `$ and an arbitrary integer $`n_0`$, $$\underset{ϵ0}{lim}ϵ\underset{n=n_0}{\overset{\mathrm{}}{}}\mathrm{e}^{nϵ}\underset{q=0}{\overset{n}{}}A(n+\nu ,q+\kappa )A(n,q)=A(\nu ,\kappa ).$$ (9) These identities establish a new connection between combinatorics and the theory of random ensembles of $`2\times 2`$ matrices. The physical background and motivations are described in a few recent publications , to which the interested reader is referred. An immediate application of (1) is also given in . Here, we shall provide the minimum background necessary for the understanding of the combinatorial aspects of the problem, and for a self-contained exposition of our results. It will also enable us to propose a conjecture which generalizes the combinatorial approach to random matrix theory for matrices of large dimensions. We start with a few definitions: Graphs consist of $`V`$ vertices connected by $`B`$ bonds (or edges). The valency $`v_i`$ of a vertex $`i`$ is the number of bonds meeting at that vertex. We shall assume that two vertices can be connected by a single bond at most. We denote the bonds connecting the vertices $`i`$ and $`j`$ by $`b=[i,j]`$. The notation $`[i,j]`$ will be used whenever we do not need to specify the direction on the bond. Hence $`[i,j]=[j,i]`$. Directed bonds will be denoted by $`d=(i,j)`$, and we shall always use the convention that the bond is directed from the first to the second index. If $`d=(i,j)`$ we use the notation $`\widehat{d}=(j,i)`$. Let $`g^{(i)}`$ be the set of directed bonds $`(i,j)`$ which emanate from the vertex $`i`$, and $`\widehat{g}^{(i)}`$ the set of directed bonds $`(j^{},i)`$ which converge at $`i`$. The vertices $`i`$ and $`j`$ are connected if $`g^{(i)}\widehat{g}^{(j)}\mathrm{}`$. The bond $`d^{}`$ is connected to the bond $`d`$ if there is some vertex $`i`$ with $`dg^{(i)}`$ and $`d^{}\widehat{g}^{(i)}`$. $`d`$ and $`\widehat{d}`$ are always connected. The Schrödinger operator on the graph is defined after the natural metric is assigned to the bonds, and the solutions of the one-dimensional Schrödinger equation on each bond $`(\text{i}\mathrm{d}_xA)^2\psi (x)=k^2\psi (x)`$ is given as a linear combination of counter-propagating waves. $`A`$ stands here for a magnetic flux. The (complex) amplitudes of the counter propagating waves are denoted by $`a_d`$, where the subscript $`d`$ stands for the directed bond along which the wave propagates, $`d=1,..,2B`$. Appropriate boundary conditions at the vertices are imposed, and the spectrum of the Schrödinger operator on the graph is determined as the (infinite, discrete) set of energies $`k_n^2`$, for which there exists a non-trivial set of $`a_d`$ which is consistent with the boundary conditions. The condition of consistency can be expressed by the requirement that $$det(IS(k))=0$$ (10) where the bond-scattering matrix $`S(k)`$ is a unitary operator in the Hilbert space of $`2B`$ dimensional vectors of coefficients $`a_d`$. The unitarity of $`S(k)`$ ensures that the spectrum of the Schrödinger operator is real. The matrix $`S(k)`$, which is the object of our study is defined as $$S_{(i,j),(l,m)}(k)=\delta _{j,l}\mathrm{e}^{\text{i}\varphi _{(i,j)}(k)}\sigma _{i,m}^{(j)}(k).$$ (11) The matrix elements $`S_{d,d^{}}(k)`$ vanish if the bonds are not connected. As a consequence, the unitarity of $`S`$ implies also the unitarity of the $`v_j`$-dimensional vertex-scattering matrices $`\sigma _{i,m}^{(j)}`$ and vice versa. The phases $`\varphi _{(i,j)}`$, are given in terms of the bond length $`L_{[i,j]}`$, and the magnetic flux $`A_{(i,j)}=A_{(j,i)}`$ , $$\varphi _{(i,j)}(k)=(k+A_{(i,j)})L_{[i,j]}.$$ (12) The two phases pertaining to the same bond $`\varphi _d`$ and $`\varphi _{\widehat{d}}`$ are equal when $`A_d=A_{\widehat{d}}=0`$. In this case $`S`$ is symmetric and the Schrödinger operator on the graph is invariant under time reversal. Time-reversal symmetry is violated when some magnetic fluxes do not vanish. $`S`$ can also be interpreted as a quantum time evolution operator describing the scattering of waves with wave number $`k`$ between connected bonds. The wave gains the phase $`\varphi _{(i,j)}(k)`$ during the propagation along the bond $`(i,j)`$, while the $`\sigma _{i,m}^{(j)}`$ describe the scattering at the vertices. In this picture, the unitarity of $`S`$ guarantees the conservation of probability during the time evolution. We will avoid unnecessary technical difficulties and consider the matrices $`\sigma _{i,m}^{(j)}`$ to be $`k`$ independent constants. One may find explicit expressions for $`\sigma _{i,m}^{(j)}`$ by requiring besides unitarity that the wave function is continuous at the vertices. The resulting expression is $$\sigma _{j,j^{}}^{(i)}=\frac{2}{v_i}\delta _{j,j^{}}\text{(Neumann b. c.)}.$$ (13) Note that back-scattering ($`j=j^{}`$) is singled out both in sign and in magnitude. In all nontrivial cases ($`v_i>2`$) the back-scattering amplitude is negative, and the back-scattering probability $`\|\sigma _{j,j}^{(i)}\|^2`$ approaches $`1`$ as the valency $`v_i`$ increases. Finally, a “classical analogue” of the quantum dynamics can be defined as a random walk on the directed bonds, in which the transition probability between bonds $`(i,j)`$, $`(j,l)`$ connected at vertex $`j`$ is $`\|\sigma _{i,l}^{(j)}\|^2`$. The resulting classical evolution operator with matrix elements $$U_{d,d^{}}=\|S_{d,d^{}}\|^2$$ (14) is probability conserving, since unitarity implies $`_i\|\sigma _{i,l}^{(j)}\|^2=1`$. ## II Spectral Two-Point Correlations and Periodic-Orbit Sums The spectrum of $`S`$ consists of $`2B`$ eigenvalues $`\mathrm{e}^{i\theta _l}`$ which are confined to the unit circle. Their distribution is given in terms of the spectral density $$d(\theta )\underset{l=1}{\overset{2B}{}}\delta _{2\pi }(\theta \theta _l)=\frac{2B}{2\pi }+\frac{1}{2\pi }\underset{n=1}{\overset{\mathrm{}}{}}s_n\text{e}^{\text{i}\theta n}+\mathrm{c}.\mathrm{c}.,$$ (15) where $`\delta _{2\pi }`$ denotes the $`2\pi `$ periodic delta function. The first term on the r.h.s. is the average density $`\overline{d}=\frac{2B}{2\pi }`$. The coefficients of the oscillatory part $$s_n=\mathrm{tr}S^n$$ will play an important rôle in the following. $`s_n=\mathrm{tr}S^n`$ is a sum over products of $`n`$ matrix elements of $`S`$, and because of (11) the bond indices of each summand describe a connected $`n`$-cycle ($``$ $`n`$-periodic orbit) on the graph $$s_n=\underset{p𝒫_n}{}𝒜_p\mathrm{e}^{i\mathrm{\Phi }_p}.$$ (16) In (16) $`𝒫_n`$ denotes the set of all $`n`$-periodic orbits (PO’s) on the graph. Note that for the convenience of presentation we will consider cycles differing only by a cyclic permutation as different PO’s. The phases $`\mathrm{\Phi }_p=_{j=0}^{n1}\varphi _{d_j}`$ can be interpreted as the action along the PO $`p`$. The amplitudes $`𝒜_p`$ are given by $$𝒜_p=\underset{j=0}{\overset{n1}{}}S_{d_{j+1},d_j},$$ (17) where $`j`$ is understood mod $`n`$. Sometimes it is useful to split the amplitude in its absolute value and a phase $`\mathrm{e}^{\text{i}\mu _p\pi }`$. For example, in the case of Neumann b. c. (13) $`\mu _p`$ is an integer counting the number of back-scatterings along $`p`$. In complete analogy to (16) we can represent also the traces of powers of the classical evolution operator $$u_n=\mathrm{tr}U^n$$ (18) as sums of periodic orbits of the graph. The two-point correlations in the spectrum of $`S`$ (15) can be expressed in terms of the average excess probability density $`R_2(r;\beta )`$ of finding two phases at a distance $`r`$, where $`r`$ is measured in units of the mean spacing $`\frac{2\pi }{2B}`$, $`R_2(r;\beta )`$ $`=`$ $`{\displaystyle \frac{1}{2B}}{\displaystyle _\pi ^{+\pi }}d\theta d\left(\theta \right)d\left(\theta [\pi /B]r\right)_\beta {\displaystyle \frac{B}{\pi }}`$ (19) $`=`$ $`{\displaystyle \frac{2}{2\pi }}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{cos}\left({\displaystyle \frac{\pi }{B}}nr\right){\displaystyle \frac{1}{2B}}\|s_n\|^2_\beta .`$ (20) The bond scattering matrix depends parametrically on the phases $`\varphi _d`$ (12). We shall define two statistical ensembles for $`S`$ in the following way. The ensemble for which time-reversal symmetry is broken consists of $`S`$ matrices for which the $`\varphi _d`$ are all different, and we consider them as independent variables distributed uniformly on the $`2B`$ torus. Invariance under time reversal implies $`\varphi _d=\varphi _{\widehat{d}}`$ and the corresponding ensemble is defined in terms of $`B`$ independent and uniformly distributed phases. We shall distinguish between these ensembles by the value of the parameter $`\beta =\{`$number of independent phases$`\}/B`$. Expectation values with respect to these measures are denoted in (20) by triangular brackets, $$\mathrm{}_\beta \underset{d}{\overset{\beta B}{}}\left(\frac{1}{2\pi }_\pi ^{+\pi }d\varphi _d\right)\mathrm{}.$$ (21) The Fourier transform of $`R_2(r;\beta )`$ is the form factor $$K(n/2B;\beta )=\frac{1}{2B}\|s_n\|^2_\beta $$ (22) on which our interest will be focussed. If the eigenvalues of the $`S`$ were statistically independent and uniformly distributed on the unit circle, $`K(n/2B)=1`$ for all $`n`$. Any deviation of the form factor from unity implies spectral correlations. Using (16) the form factor (22) is expressed as a double sum over PO’s $`K(n/2B;\beta )`$ $`=`$ $`{\displaystyle \frac{1}{2B}}\left\|{\displaystyle \underset{p𝒫_n}{}}𝒜_p\text{e}^{\text{i}\mathrm{\Phi }_p}\right\|^2_\beta `$ (23) $`=`$ $`{\displaystyle \frac{1}{2B}}{\displaystyle \underset{p,p^{}𝒫_n}{}}𝒜_p𝒜_p^{}\text{e}^{\text{i}(\mathrm{\Phi }_p\mathrm{\Phi }_p^{})}_\beta `$ (24) In order to perform the average over all the phases $`\varphi _d`$ in (23) we write $`\mathrm{\Phi }_p={\displaystyle \underset{d}{}}n_d^{(p)}\varphi _d,`$ (25) where $`n_d^{(p)}`$ counts the number of traversals of each directed bond such that $`_dn_d^{(p)}=n`$. According to (21) we have $`\text{e}^{\text{i}(n\varphi _d+n^{}\varphi _d^{})}_{\beta =1}`$ $`=`$ $`\delta _{n,0}\delta _{n^{},0},`$ (26) $`\text{e}^{\text{i}(n\varphi _d+\widehat{n}\varphi _{\widehat{d}})}_{\beta =2}`$ $`=`$ $`\delta _{n+\widehat{n},0}.`$ (27) Thus, the double sum in (23) can be restricted to families of orbits. For $`\beta =2`$, let $``$ be the family of isometric PO’s which have the same integers $`n_d^{()}`$. That is, the family consists of all the PO’s which traverse the same directed bonds the same number of times, but not necessarily in the same order. In the case $`\beta =1`$, $``$ consists of all PO’s sharing $`n_b^{()}n_d^{()}+n_{\widehat{d}}^{()}`$. That is, the family contains all PO’s which traverse the same set of undirected bonds the same number of times, irrespective of direction or order. We find $`K(n/2B;\beta )`$ $`=`$ $`{\displaystyle \frac{1}{2B}}{\displaystyle \underset{(\beta )_n}{}}\|{\displaystyle \underset{p}{}}𝒜_p\|^2.`$ (28) $`(\beta )_n`$ denotes the set of all vectors $`=[n_d]`$ for $`\beta =2`$ ($`=[n_b]`$ for $`\beta =1`$) of $`\beta B`$ non-negative integers summing to $`n`$, for which at least one PO exists. For Neumann b. c. (13), e. g., (28) amounts to counting the PO’s in a given set $`=[n_d]`$ taking into account the number of back-scatterings along the orbit. The problem of spectral statistics is now reduced to a counting (combinatorial) problem which is, however, very complicated in general. Even the determination of the number of families $``$ for a given $`n`$ is difficult. For $`\beta =2`$ an obvious necessary condition for the existence of a PO with a given set of bond traversals $`=[n_d]`$ is that at any vertex the number of incoming and outgoing bonds is the same, i. e. $$\underset{dg^{(i)}}{}n_d=\underset{d\widehat{g}^{(i)}}{}n_d(i=1,\mathrm{}V).$$ For $`\beta =1`$, the analogous condition reads $$\underset{dg^{(i)}}{}n_d\underset{d\widehat{g}^{(i)}}{}n_d\text{ mod }2=0(i=1,\mathrm{}V),$$ i. e. the total number of traversals of adjacent bonds should be even at each vertex. However, it is not so easy to formulate a sufficient condition for the existence of a PO given a set of numbers $`n_d`$. In particular one must take care to exclude cases, in which the set of traversed bonds is a union of two or more disconnected groups (“composite orbits”). Extensive numerical work revealed that for fully connected graphs ($`v_jV1`$), and for $`V1`$, the form-factor (22) is well reproduced by the predictions of random matrix theory for the Circular Orthogonal Ensemble (COE) ($`\beta =1`$) or the Circular Unitary Ensemble (CUE) ($`\beta =2`$). This leads us to expect that (28) approaches the corresponding random matrix prediction in the limit $`V\mathrm{}`$. This conjecture is proposed as a challenge to asymptotic combinatorial theory. ## III The Ring Graph In the following we will evaluate explicitly the quantities introduced in the previous section for one of the simplest quantum graphs. It consists of a single vertex on a loop (see fig. 1). There are two directed bonds $`d=1`$ and $`\widehat{d}=2`$ with $`\varphi _1\varphi _2`$, i. e. time-reversal symmetry is broken. Since this graph would be trivial for Neumann b.c. the vertex-scattering matrix at the only vertex is chosen as $$\sigma (\eta )=\left(\begin{array}{cc}\mathrm{cos}\eta \hfill & \mathrm{i}\mathrm{sin}\eta \hfill \\ \mathrm{i}\mathrm{sin}\eta \hfill & \mathrm{cos}\eta \hfill \end{array}\right),$$ (29) with $`0\eta \pi /2`$. The corresponding bond-scattering matrix is $$S(\eta )=\left(\begin{array}{cc}\mathrm{e}^{\varphi _1}\hfill & 0\hfill \\ 0\hfill & \mathrm{e}^{\varphi _2}\hfill \end{array}\right)\left(\begin{array}{cc}\mathrm{cos}\eta \hfill & \mathrm{i}\mathrm{sin}\eta \hfill \\ \mathrm{i}\mathrm{sin}\eta \hfill & \mathrm{cos}\eta \hfill \end{array}\right).$$ (30) We shall compute the form factor for two ensembles. The first is defined by a fixed value of $`\eta =\pi /4`$, and the only average is over the phases $`\varphi _d`$ according to (26). The second ensemble includes an additional averaging over the parameter $`\eta `$. We will show that the measure for the integration over $`\eta `$ can be chosen such that the model yields exactly the CUE form factor for 2$`\times `$ 2 random matrices . ### A Periodic Orbit Representation of $`u_n`$ We will first illustrate our method of deriving combinatorial results from the ring graph in a case where a known identity is obtained. Consider the classical evolution operator $`U`$ of the ring graph. According to (14) we have $$U(\eta )=\left(\begin{array}{cc}\mathrm{cos}^2\eta \hfill & \mathrm{sin}^2\eta \hfill \\ \mathrm{sin}^2\eta \hfill & \mathrm{cos}^2\eta \hfill \end{array}\right).$$ (31) The spectrum of $`U`$ consists of $`\{1,\mathrm{cos}2\eta \}`$, such that $$u_n(\eta )=1+\mathrm{cos}^n2\eta .$$ (32) We will now show how this result can be obtained from a sum over the periodic orbits of the system, grouped into families of orbits as in (28). In the classical calculation it is actually not necessary to take the families into account, but we would like to stress the analogy to the quantum case considered below. The periodic orbit expansion of the classical return probability can easily be obtained from (31) by expanding all matrix products in (18). We find $`u_n`$ $`=`$ $`{\displaystyle \underset{i_1=1,2}{}}\mathrm{}{\displaystyle \underset{i_n=1,2}{}}{\displaystyle \underset{j=0}{\overset{n1}{}}}U_{i_j,i_{j+1}}(\eta ),`$ (33) where $`j`$ is again taken mod $`n`$. In the following the binary sequence $`[i_j]`$ ($`i_j\{1,2\}`$; $`j=0,\mathrm{},n1`$) is referred to as the code of the orbit. We will now sort the terms in the multiple sum above into families of isometric orbits. In the present case a family is completely specified by the integer $`qq_1`$ which counts the traversals of the loop $`1`$, i.e., the number of letters $`1`$ in the code word. Each of these $`q`$ letters is followed by an uninterrupted sequence of $`t_j0`$ letters $`2`$ with the restriction that the total number of letters $`2`$ is given by $$\underset{j=1}{\overset{q}{}}t_j=nq.$$ (34) We conclude that each code word in a family $`0<q<n`$ which starts with $`i_1=1`$ corresponds to an ordered partition of the number $`nq`$ into $`q`$ non-negative integers, while the words starting with $`i_1=2`$ can be viewed as partition of $`q`$ into $`nq`$ summands. To make this step very clear, consider the following example: All code words of length $`n=5`$ in the family $`q=2`$ are $`11222`$, $`12122`$, $`12212`$, $`12221`$ and $`22211`$, $`22121`$, $`21221`$, $`22112`$, $`21212`$, $`21122`$. The first four words correspond to the partitions $`0+3=1+2=2+1=3+0`$ of $`nq=3`$ into $`q=2`$ terms, while the remaining $`5`$ words correspond to $`2=0+0+2=0+1+1=1+0+1=0+2+0=1+1+0=2+0+0`$. In the multiple products in (33) a backward scattering along the orbit is expressed by two different consecutive symbols $`i_ji_{j+1}`$ in the code and leads to a factor $`\mathrm{sin}^2\eta `$, while a forward scattering contributes a factor $`\mathrm{cos}^2\eta `$ . Since the sum is over periodic orbits, the number of back scatterings is always even and we denote it with $`2\nu `$. It is then easy to see that $`\nu `$ corresponds to the number of positive terms in the partitions introduced above, since each such term corresponds to an uninterrupted sequence of symbols $`2`$ enclosed between two symbols $`1`$ or vice versa and thus contributes two back scatterings. For the codes starting with a symbol $`1`$ there are $`\left(\genfrac{}{}{0pt}{}{q}{\nu }\right)`$ ways to choose the $`\nu `$ positive terms in the sum of $`q`$ terms, and there are $`\left(\genfrac{}{}{0pt}{}{nq1}{\nu 1}\right)`$ ways to decompose $`nq`$ into $`\nu `$ positive summands. After similar reasoning for the codes starting with the symbol $`2`$ we find for the periodic orbit expansion of the classical return probability $`u_n(\eta )`$ $`=`$ $`2\mathrm{cos}^{2n}\eta +{\displaystyle \underset{q=1}{\overset{n1}{}}}{\displaystyle \underset{\nu }{}}\left[\left({\displaystyle \genfrac{}{}{0pt}{}{q}{\nu }}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{nq1}{\nu 1}}\right)+\left({\displaystyle \genfrac{}{}{0pt}{}{nq}{\nu }}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{q1}{\nu 1}}\right)\right]\mathrm{sin}^{4\nu }\eta \mathrm{cos}^{2n4\nu }\eta `$ (35) $`=`$ $`2\mathrm{cos}^{2n}\eta +{\displaystyle \underset{q=1}{\overset{n1}{}}}{\displaystyle \underset{\nu }{}}{\displaystyle \frac{n}{\nu }}\left({\displaystyle \genfrac{}{}{0pt}{}{q1}{\nu 1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{nq1}{\nu 1}}\right)\mathrm{sin}^{4\nu }\eta \mathrm{cos}^{2n4\nu }\eta `$ (36) The summation limits for the variable $`\nu `$ are implicit since all terms outside vanish due to the properties of the binomial coefficients. From the equivalence between (32) and (36) the combinatorial identity $$\underset{q=1}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{q1}{\nu 1}\right)\left(\genfrac{}{}{0pt}{}{nq1}{\nu 1}\right)=\left(\genfrac{}{}{0pt}{}{n1}{2\nu 1}\right)=\frac{2\nu }{n}\left(\genfrac{}{}{0pt}{}{n}{2\nu }\right).$$ (37) could be deduced which indeed reduces (36) to a form $`u_n(\eta )`$ $`=`$ $`2{\displaystyle \underset{\nu }{}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{2\nu }}\right)\mathrm{sin}^{4\nu }\eta \mathrm{cos}^{2n4\nu }\eta `$ (38) $`=`$ $`(\mathrm{cos}^2\eta +\mathrm{sin}^2\eta )^n+(\mathrm{cos}^2\eta \mathrm{sin}^2\eta )^n,`$ (39) which is obviously equivalent to (32). (37) can also be derived by some straight forward variable substitutions from the identity $$\underset{k=l}{\overset{nm}{}}\left(\genfrac{}{}{0pt}{}{k}{l}\right)\left(\genfrac{}{}{0pt}{}{nk}{m}\right)=\left(\genfrac{}{}{0pt}{}{n+1}{l+m+1}\right)$$ (40) which is found in the literature . ### B Quantum Mechanics: Spacing Distribution and Form Factor In the following two subsections we derive novel combinatorial identities by applying the reasoning which led to (37) to the quantum evolution operator (30) of the ring graph. We can write the eigenvalues of $`S(\eta )`$ as $`\mathrm{e}^{\text{i}(\varphi _1+\varphi _2)/2}\mathrm{e}^{\pm \text{i}\lambda /2}`$ with $$\lambda =2\mathrm{arcos}\left[\mathrm{cos}\eta \mathrm{cos}\left(\frac{\varphi _1\varphi _2}{2}\right)\right]$$ (41) denoting the difference between the eigenphases. For the two-point correlator we find $`R_2(r,\eta )`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _\pi ^{+\pi }}d\theta d\left(\theta +{\displaystyle \frac{\pi r}{2}}\right)d\left(\theta {\displaystyle \frac{\pi r}{2}}\right)_{\varphi _{1,2}}{\displaystyle \frac{1}{\pi }}`$ (42) $`=`$ $`{\displaystyle \frac{\delta _2(r)1}{\pi }}{\displaystyle \frac{\delta _{2\pi }\left(\pi r+\lambda \right)+\delta _{2\pi }\left(\pi r\lambda \right)}{2}}_{\varphi _{1,2}}`$ (43) $`=`$ $`{\displaystyle \frac{\delta _2\left(r\right)1}{\pi }}+{\displaystyle \frac{\mathrm{sin}\|\pi r/2\|}{2\pi }}{\displaystyle \frac{\mathrm{\Theta }(\mathrm{cos}^2\eta \mathrm{cos}^2(\pi r/2))}{\sqrt{\mathrm{cos}^2\eta \mathrm{cos}^2(\pi r/2)}}})`$ (44) Here, $`\delta _2\left(r\right)`$ is the $`2`$-periodic $`\delta `$ function. In particular for equal transmission and reflection probability ($`\eta =\pi /4`$) we have $`R_2(r,\pi /4)`$ $`=`$ $`{\displaystyle \frac{\delta _2\left(r\right)1}{\pi }}+{\displaystyle \frac{1}{2\pi }}\sqrt{{\displaystyle \frac{\mathrm{cos}(\pi r)1}{\mathrm{cos}(\pi r)}}}\mathrm{\Theta }\left({\displaystyle \frac{1}{2}}\|r1\|\right)`$ (45) and, by a Fourier transformation, we can compute the form factor $`K(n,\pi /4)`$ $`=`$ $`\pi {\displaystyle _0^2}dr\mathrm{cos}\left(n\pi r\right)R_2(r,\pi /4)`$ (46) $`=`$ $`1+{\displaystyle \frac{(1)^{m+n}}{2^{2m+1}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2m}{m}}\right){\displaystyle \frac{3}{2}}\delta _{n,0}`$ (47) $``$ $`1+{\displaystyle \frac{(1)^{m+n}}{2\sqrt{\pi n}}}(n1),`$ (48) where $`m=[n/2]`$ and $`[]`$ stands for the integer part. Next we consider the ensemble for which transmission and reflection probabilities are uniformly distributed between $`0`$ and $`1`$. For the parameter $`\eta `$ this corresponds to the measure $`\mathrm{d}\mu (\eta )=2\|\mathrm{cos}\eta \mathrm{sin}\eta \|\mathrm{d}\eta `$. The main reason for this choice is that upon integrating (44) one gets $`R_2^{(\mathrm{av})}(r)`$ $`=`$ $`{\displaystyle \frac{\delta _2\left(r\right)1}{\pi }}+{\displaystyle \frac{\mathrm{sin}^2(\pi r/2)}{\pi }}`$ (49) which coincides with the CUE result for $`2\times 2`$ matrices. A Fourier transformation results in $$K_2(n)=\{\begin{array}{cc}\frac{1}{2}\hfill & \mathrm{for}n=1\hfill \\ 1\hfill & \mathrm{for}n2\hfill \end{array}.$$ (50) The form factors (47), (48) and (50) are displayed in Fig. 1 below. ### C Periodic Orbit Expansion of the Form Factor An explicit formulation of (28) for the ring graph is found by labelling and grouping orbits as explained in the derivation of (36). We obtain $`K_2(n;\eta )`$ $`=`$ $`\mathrm{cos}^{2n}\eta +{\displaystyle \frac{n^2}{2}}{\displaystyle \underset{q=1}{\overset{n1}{}}}\left[{\displaystyle \underset{\nu }{}}{\displaystyle \frac{(1)^\nu }{\nu }}\left({\displaystyle \genfrac{}{}{0pt}{}{q1}{\nu 1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{nq1}{\nu 1}}\right)\mathrm{sin}^{2\nu }\eta \mathrm{cos}^{n2\nu }\eta \right]^2,`$ (51) where $`q`$ denotes the number of traversals of the ring in positive direction and $`2\nu `$ is the number of backward scatterings along the orbit. The inner sum over $`\nu `$ can be written in terms of Kravtchouk polynomials as $`K_2(n;\eta )`$ $`=`$ $`\mathrm{cos}^{2n}\eta +{\displaystyle \frac{1}{2}}{\displaystyle \underset{q=1}{\overset{n1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{nq}}\right)\mathrm{cos}^{2q}\eta \mathrm{sin}^{2(nq)}\eta \left[{\displaystyle \frac{n}{q}}P_{n1,nq}^{(\mathrm{cos}^2\eta ,\mathrm{sin}^2\eta )}(q)\right]^2,`$ (52) and the Kravtchouk polynomials are defined as in by $`P_{N,k}^{(u,v)}(x)=\left[\left({\displaystyle \genfrac{}{}{0pt}{}{N}{k}}\right)(uv)^k\right]^{1/2}{\displaystyle \underset{\nu =0}{\overset{k}{}}}(1)^{k\nu }\left({\displaystyle \genfrac{}{}{0pt}{}{x}{\nu }}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{Nx}{k\nu }}\right)u^{k\nu }v^\nu \left(\begin{array}{c}0kN\hfill \\ u+v=1\hfill \end{array}\right).`$ (55) These functions form a complete system of orthogonal polynomials of integers $`x`$ with $`0xN`$. They have quite diverse applications ranging from the theory of covering codes to the statistical mechanics of polymers . The same functions appear also as a building block in our periodic-orbit theory of Anderson localization on graphs . Unfortunately, we were not able to reduce the above expression any further by using the known sum rules and asymptotic representations for Kravtchouk polynomials. The main obstacle stems from the fact that in our case the three numbers $`N,k,x`$ in the definition (55) are constrained by $`N=k+x1`$. We will now consider the special case $`\eta =\pi /4`$ for which we obtained in the previous subsection the solution (47). The result can be expressed in terms of Kravtchouk polynomials with $`u=v=1/2`$ which is also the most important case for the applications mentioned above. We adopt the common practice to omit the superscript $`(u,v)`$ in this special case and find $`K_2(n;\pi /4)`$ $`=`$ $`{\displaystyle \frac{1}{2^n}}+{\displaystyle \frac{1}{2^{n+1}}}{\displaystyle \underset{q=1}{\overset{n1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{nq}}\right)\left[{\displaystyle \frac{n}{q}}P_{n1,nq}(q)\right]^2.`$ (56) It is convenient to introduce at this point the symbol $`𝒩(s,t)`$ defined in Eq. (2). It allows to rewrite (56) with the help of some standard transformations of binomial coefficients as $`K_2(n;\pi /4)`$ $`=`$ $`{\displaystyle \frac{1}{2^n}}+{\displaystyle \frac{1}{2^{n+1}}}{\displaystyle \underset{q=1}{\overset{n1}{}}}\left[{\displaystyle \frac{n}{q}}𝒩(q,nq1)\right]^2`$ (57) $`=`$ $`{\displaystyle \frac{1}{2^n}}+{\displaystyle \frac{1}{2^{n+1}}}{\displaystyle \underset{q=1}{\overset{n1}{}}}\left[𝒩(q,nq)+(1)^n𝒩(nq,q)\right]^2`$ (58) This expression is displayed in Fig. 1 together with (47) in order to illustrate the equivalence of the two results for the form factor. An independent proof for the equivalence of (47), (56) can be given by comparing the generating functions of $`K_2(n;\pi /4)`$ in the two representations . We define $`G(x)`$ $`=`$ $`{\displaystyle \underset{x=1}{\overset{\mathrm{}}{}}}K_2(n;\pi /4)(2x)^n(\|x\|<1/2)`$ (59) and find from (47) $`G(x)`$ $`=`$ $`{\displaystyle \frac{2x}{12x}}{\displaystyle \frac{1}{2}}+{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{2}}\left({\displaystyle \genfrac{}{}{0pt}{}{2m}{m}}\right)x^{2m}(12x)`$ (60) $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{12x}{\sqrt{1+4x^2}}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{16x}{12x}}.`$ (61) On the other hand we have from (57) $`G(x)={\displaystyle \frac{x}{1x}}+G_1(x)+G_2(x)`$ (62) with $`G_{1,2}(x)`$ defined in the introduction. A convenient starting point to obtain the r.h.s. of (3), (4) is the integral representation $`𝒩(s,t)={\displaystyle \frac{(1)^t}{2\pi \text{i}}}{\displaystyle dz(1+z^1)^t(1z)^{s1}},`$ (63) where the contour encircles the origin. With the help of (63) we find $`g(x,y)`$ $`=`$ $`{\displaystyle \underset{s,t=1}{\overset{\mathrm{}}{}}}𝒩(s,t)x^sy^t`$ (64) $`=`$ $`{\displaystyle \frac{1}{2\pi \text{i}}}{\displaystyle \underset{s,t=1}{\overset{\mathrm{}}{}}}{\displaystyle dz\underset{s,t=1}{\overset{\mathrm{}}{}}(1+z^1)^t(1z)^{s1}x^s(y)^t}`$ (65) $`=`$ $`{\displaystyle \frac{xy}{2\pi \text{i}}}{\displaystyle \underset{s,t=0}{\overset{\mathrm{}}{}}}{\displaystyle dz\frac{1}{1x(1z)}\frac{1+z}{z+y(1+z)}}`$ (66) $`=`$ $`{\displaystyle \frac{xy}{(1+y)(1x+y2xy)}}(\|x\|,\|y\|<1/\sqrt{2}),`$ (67) which was already stated in (5). The contour $`\|1+z^1\|=\|1z\|=\sqrt{2}`$ has been chosen such that both geometric series converge everywhere on it. Now we have $`G_1(x^2)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi \text{i})^2}}{\displaystyle \frac{\mathrm{d}z\mathrm{d}z^{}}{zz^{}}\underset{s,t=1}{\overset{\mathrm{}}{}}\underset{s^{},t^{}=1}{\overset{\mathrm{}}{}}𝒩(s,t)𝒩(s^{},t^{})(xz)^s(x/z)^s^{}(xz^{})^t(x/z^{})^t^{}}`$ (68) $`=`$ $`{\displaystyle \frac{x^4}{(2\pi \text{i})^2}}{\displaystyle dzdz^{}\frac{1}{(1+xz^{})(1+x[z^{}z]2x^2zz^{})}\frac{z^{}}{(z^{}+x)(zz^{}+x[zz^{}]2x^2)}},`$ (69) where $`\|x\|<1/\sqrt{2}`$ and the contour for $`z,z^{}`$ is the unit circle. We perform the double integral using the residua inside the contour and obtain (3) and in complete analogy also (4) such that finally $$G(x)=\frac{x}{1x}+\frac{x}{2x1}\left(\frac{1}{\sqrt{4x^2+1}}\frac{1}{1x}\right)+\frac{1}{2}\frac{4x^22x+1}{(12x)\sqrt{4x^2+1}}\frac{1}{2}.$$ (70) The proof is completed by a straight forward verification of the equivalence between the rational functions (60) and (70). The identities (6), (7) follow now by separating even and odd powers of $`n`$ in (47) and (56). In terms of Kravtchouk polynomials these identities can be written as $$\underset{q=1}{\overset{2m1}{}}\left(\genfrac{}{}{0pt}{}{2m1}{2mq}\right)\left[\frac{2m}{q}P_{2m1,2mq}(q)\right]^2=2^{2m+1}+(1)^m\left(\genfrac{}{}{0pt}{}{2m}{m}\right)2$$ (71) and $$\underset{q=1}{\overset{2m}{}}\left(\genfrac{}{}{0pt}{}{2m}{2m+1q}\right)\left[\frac{2m+1}{q}P_{2m,2m+1q}(q)\right]^2=2^{2m+2}2(1)^m\left(\genfrac{}{}{0pt}{}{2m}{m}\right)2.$$ (72) Finally we will derive the CUE result (50) for the ensemble of graphs defined in the previous subsection starting from the periodic-orbit expansion (51). We find $`K_2(n)`$ $`=`$ $`{\displaystyle _0^{\pi /2}}d\mu (\eta )K_2(n;\eta ).`$ (73) Inserting (51), expanding into a double sum and using $$_0^{\pi /2}d\eta \mathrm{sin}^{2(\nu +\nu ^{})+1}\eta \mathrm{cos}^{2(n\nu \nu ^{})+1}\eta =\frac{1}{2(n+1)}\left(\genfrac{}{}{0pt}{}{n}{\nu +\nu ^{}}\right)^1$$ (74) we get $`K_2(n)`$ $`=`$ $`{\displaystyle \frac{1}{n+1}}+`$ (76) $`+{\displaystyle \frac{n^2}{4(n+1)}}{\displaystyle \underset{q=1}{\overset{n1}{}}}{\displaystyle \underset{\nu ,\nu ^{}}{}}{\displaystyle \frac{(1)^{\nu +\nu ^{}}}{\nu \nu ^{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{\nu +\nu ^{}}}\right)^1\left({\displaystyle \genfrac{}{}{0pt}{}{q1}{\nu 1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{nq1}{\nu 1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{q1}{\nu ^{}1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{nq1}{\nu ^{}1}}\right).`$ Comparing this to the equivalent result (50) we were led to the identity (1) involving a multiple sum over binomial coefficients. In this case, an independent computer-generated proof was found , which is based on the recursion relation $$q^2F_{\nu ,\nu ^{}}(n,q)(nq1)^2F_{\nu ,\nu ^{}}(n,q+1)+(n1)(n2q1)F_{\nu ,\nu ^{}}(n+1,q+1)=0.$$ (77) This recursion relation was obtained with the help of a Mathematica routine , but it can be checked manually in a straight forward calculation. By summing (77) over the indices $`\nu ,\nu ^{}`$, the same recursion relation is shown to be valid for $`S(n,q)`$ and the proof is completed by demonstrating the validity of (1) for a few initial values. Having proven (1) we can use it to perform the summation over $`\nu ,\nu ^{}`$ in (76) and find $`K_2(n)={\displaystyle \frac{1}{n+1}}+{\displaystyle \underset{q=1}{\overset{n1}{}}}{\displaystyle \frac{n}{n^21}}={\displaystyle \frac{1}{n+1}}+{\displaystyle \frac{n}{n+1}}(1\delta _{n,1}),`$ (78) which is now obviously equivalent to the random-matrix form factor (50). ### D Trace identities Let $`S`$ be an arbitrary unitary matrix with a non degenerate spectrum and $`s_n=\mathrm{tr}S^n`$. Then $`\underset{ϵ0}{lim}ϵ{\displaystyle \underset{n=n_0}{\overset{\mathrm{}}{}}}s_n^{}s_{n+\nu }\mathrm{e}^{nϵ}=s_\nu ,`$ (79) for arbitrary integers $`n_0`$ and $`\nu `$ . We shall now apply this identity to the ring graph with $`\eta =\pi /4`$ in order to prove (9). Again, the traces of $`S^n`$ can be expanded in periodic orbits which are grouped into $`n+1`$ families with equal phases $`\mathrm{\Phi }(n,q)=q\varphi _1+(nq)\varphi _2`$ ($`0qn`$). Thus, one can write $$s_n(k)=\underset{q=0}{\overset{n}{}}A(n,q)\mathrm{e}^{\text{i}\mathrm{\Phi }(n,q)},$$ (80) where $`A(n,q)`$ is the coherent sum of all amplitudes of PO’s in the corresponding set. $`A(n,q)`$ can be expressed in terms of Kravtchouk Polynomials as $$A(n,q)=\frac{1}{\sqrt{2^n}}\{\begin{array}{cc}1\hfill & \mathrm{for}q=0\mathrm{or}n\hfill \\ (1)^{n+q}(n/q)\left(\genfrac{}{}{0pt}{}{n1}{nq}\right)^{1/2}P_{n1,nq}(q)\hfill & \mathrm{for}0<q<n\hfill \end{array}$$ (81) (compare Eq. (52)). This is equivalent to (8). Substituting (80) into the trace identity (79), we get for arbitrary integers $`\nu `$ and $`n_0`$ $`\underset{ϵ0}{lim}ϵ`$ $`{\displaystyle \underset{n=n_0}{\overset{\mathrm{}}{}}}`$ $`\mathrm{e}^{nϵ}{\displaystyle \underset{q=0}{\overset{n}{}}}{\displaystyle \underset{p=0}{\overset{n+\nu }{}}}A(n+\nu ,p)A(n,q)\mathrm{e}^{\text{i}\left[(pq)\varphi _1+(\nu (pq))\varphi _2\right]}`$ (82) $`=`$ $`{\displaystyle \underset{\kappa =0}{\overset{\nu }{}}}A(\nu ,\kappa )\mathrm{e}^{\text{i}\left[\kappa \varphi _1+(\nu \kappa )\varphi _2\right]}.`$ (83) This is valid for arbitrary phases $`\varphi _1`$ and $`\varphi _2`$ and therefore the coefficients of the phase factors $`\mathrm{e}^{\text{i}\mathrm{\Phi }(\nu ,\kappa )}`$ on both sides are equal. Eq. (9) follows. ## IV Conclusions We have shown how within periodic-orbit theory the problem of finding the form factor (the spectral two-point correlation function) for a quantum graph can be reduced exactly to a well-defined combinatorial problem. Even for the very simple graph model that we considered in the last section the combinatorial problems involved were highly non-trivial. In fact we encountered previously unknown identities which we could not have obtained if it were not for the independent method of computing the form factor directly from the spectrum. However, since the pioneering work documented in the investigation of sums of the type we encountered in this paper is a rapidly developing subject, and it can be expected that finding identities like (1), (6) and (7) will shortly be a matter of computer power. Numerical simulations in which the form factor was computed for fully connected graphs ($`v_i=V1i`$) and, e. g., Neumann boundary conditions indicate that the spectral correlations of $`S`$ match very well Dyson’s predictions for the circular ensembles COE or CUE, respectively . The agreement is improved as $`V`$ increases. We conjecture that this might be rigorously substantiated by asymptotic combinatorial theory. A first step towards this goal was taken in the present paper. ## V Acknowledgements This research was supported by the Minerva Center for Physics of Nonlinear Systems, and by a grant from the Israel Science Foundation. We thank Uri Gavish for introducing us to the combinatorial literature, and Brendan McKay and Herbert Wilf for their interest and support. We are indebted to Gregory Berkolaiko for his idea concerning the proof of (6) and (7), and to Akalu Tefera for his kind help in obtaining a computer-aided proof of (1).
warning/0003/cond-mat0003364.html	ar5iv	text	# Monte Carlo simulation of nonlinear Couette flow in a dilute gas ## I Introduction One of the most interesting states for analyzing transport processes far from equilibrium is the steady planar Couette flow. The physical situation corresponds to a system enclosed between two infinite, parallel plates in relative motion and, in general, kept at different temperatures. These boundary conditions lead to combined heat and momentum transport. If $`x`$ and $`y`$ denote the coordinates parallel to the flow and orthogonal to the plates, respectively, then the corresponding steady hydrodynamic balance equations are $$\frac{P_{xy}}{y}=\frac{P_{yy}}{y}=0,$$ (1) $$P_{xy}\frac{u_x}{y}+\frac{q_y}{y}=0,$$ (2) where $`𝐮=u_x\widehat{𝐱}`$ is the flow velocity, $`𝖯`$ is the pressure tensor, and $`𝐪=q_x\widehat{𝐱}+q_y\widehat{𝐲}`$ is the heat flux. The presence of $`q_y`$ in Eq. (2) indicates that a thermal gradient $`T/y`$ is induced by the velocity gradient, even if both plates are kept at the same temperature. The balance equations (1) and (2) do not constitute a closed set unless the dependence of the pressure tensor and the heat flux on the hydrodynamic fields is known. If the gradients are small, the fluxes $`𝖯`$ and $`𝐪`$ are described by the Navier-Stokes (NS) constitutive relations, which in this problem yield $$P_{xx}=P_{yy}=P_{zz},P_{xy}=\eta _0\frac{u_x}{y},$$ (3) $$q_x=0,q_y=\kappa _0\frac{T}{y},$$ (4) where $`\eta _0`$ and $`\kappa _0`$ are the NS shear viscosity and thermal conductivity coefficients, respectively. As a consequence of the absence of normal stress differences in the NS description, the hydrostatic pressure $`p=(P_{xx}+P_{yy}+P_{zz})/3`$ is a constant, on account of the balance equation (1). Even in the linear regime described by the NS equations, one still needs to know the spatial dependence of the transport coefficients to obtain the exact solution of the hydrodynamic equations. The problem becomes tractable in the case of a low density gas, where the state of the system is completely specified by the velocity distribution function $`f(𝐫,𝐯;t)`$, which obeys the Boltzmann equation. A relevant dimensionless quantity in a dilute gas is the Knudsen number $`\text{Kn}=\lambda /\mathrm{}_h`$, defined as the ratio of the mean free path $`\lambda `$ to the scale length of the hydrodynamic gradients $`\mathrm{}_h`$. In many laboratory conditions, $`\text{Kn}1`$ and so the Boltzmann equation can be solved by means of the Chapman-Enskog method as an expansion of the distribution function in powers of the Knudsen number. The zeroth order approximation leads to the Euler hydrodynamic equations, while the first order approximation yields the NS equations with explicit expressions for the transport coefficients $`\eta _0`$ and $`\kappa _0`$. The results show that the ratio $`\eta _0/\kappa _0`$ is a constant. Consequently, it then follows from the NS hydrodynamic equations (1)–(4) that the flow velocity profile is quasi-linear, $$\eta _0\frac{u_x}{y}=\text{const},$$ (5) and the temperature is quasi-parabolic, $$\left(\kappa _0\frac{}{y}\right)^2T=\frac{\kappa _0}{\eta _0}\left(\eta _0\frac{u_x}{y}\right)^2=\text{const}.$$ (6) Note that the profile of $`u_x`$ is not strictly linear, due to the space dependence of $`\eta _0`$ through the temperature. Analogously, the temperature profile is not strictly quadratic. In fact, the specific form of both profiles depends on the interaction potential under consideration. On the other hand, from Eqs. (5) and (6) it is easy to derive a nice result, namely that if the temperature $`T`$ is seen as a function of $`u_x`$ rather than as a function of the coordinate space $`y`$, then one has $$\frac{^2T}{u_x^2}=\frac{\eta _0}{\kappa _0}.$$ (7) This is a sort of nonequilibrium “equation of state,” according to which the temperature is a quadratic function of the flow velocity. Moreover, the “curvature” of the profile is practically universal, given the weak influence of the interaction potential on the Prandtl number $`\text{Pr}=5k_B\eta _0/2m\kappa _0\frac{2}{3}`$, where $`k_B`$ is the Boltzmann constant and $`m`$ is the mass of a particle. Since the mean free path is inversely proportional to the density, the Knudsen number, at a given value of the scale length $`\mathrm{}_h`$, increases as the gas becomes more rarefied. In general, when the Knudsen number is not small, the NS relations are not expected to hold and the transport must be described by nonlinear constitutive equations. In the special case of Maxwell molecules (particles interacting via an $`r^4`$ potential), it has been shown that the Boltzmann equation admits a consistent solution in the nonlinear Couette flow characterized by a constant pressure $`p`$ and profiles similar to those obtained in the NS regime, Eqs. (5)–(7), except that $`\eta _0`$ and $`\kappa _0`$ are replaced by a generalized shear viscosity coefficient $`\eta (a)=\eta _0F_\eta (a)`$ and a generalized thermal conductivity coefficient $`\kappa _{yy}(a)=\kappa _0F_\kappa (a)`$, respectively. Here, $`a=(\eta _0/p)u_x/y`$ is a constant (dimensionless) shear rate and $`F_\eta `$ and $`F_\kappa `$ are nonlinear functions of $`a`$. In addition, $`P_{xx}P_{yy}P_{zz}`$ and $`q_x0`$. In this problem, the hydrodynamic scale length can be identified as $`\mathrm{}_h\sqrt{k_BT/m}(u_x/y)^1`$, while the mean free path is $`\lambda \sqrt{k_BT/m}(\eta _0/p)`$. Thus, the reduced shear rate $`a`$ represents the Knudsen number in this problem, i.e. $`a\text{Kn}`$. Henceforth, we will use the reduced shear rate $`a`$ to refer to the Knudsen number Kn. The solution considered in Refs. describes heat and momentum transport for arbitrary velocity and thermal gradients in the bulk domain, where boundary effects are negligible. On the other hand, the full nonlinear dependence of $`F_\eta (a)`$ and $`F_\kappa (a)`$ is not explicitly known, since it involves the infinite hierarchy of moment equations. Their knowledge is limited to super-Burnett order and the result is $`F_\eta (a)=13.111a^2`$ and $`F_\kappa (a)=17.259a^2`$. Consequently, if one wants to get the transport properties for arbitrary values of the shear rate and the thermal gradient, one must resort to approximate schemes or to computer simulations. In the first alternative, explicit expressions for the nonlinear transport coefficients in the Couette flow have been obtained from exact solutions of the Bhatnagar-Gross-Krook (BGK) model and related models for general interactions, as well as from the Grad method. In the simulation side, Risso and Cordero have recently studied the shear-rate dependence of $`F_\eta `$ and $`F_\kappa `$ by means of molecular dynamics simulations of a hard disk gas. Comparison between the different analytical results with those obtained from the simulation shows that the predictions given by kinetic models are in better agreement than those given by the Grad method, especially in the case of the thermal conductivity. Nevertheless, given the difficulties associated with molecular dynamics simulations to achieve large shear rates in a dilute gas, the above comparison is restricted to a range of shear rates for which non-Newtonian effects are hardly significant. For instance, the shear viscosity has only decreased around $`10\%`$ with respect to its Navier-Stokes value for the largest value of the shear rate considered by Risso and Cordero. In order to overcome such difficulties and extend the range of values of $`a`$, one may use the so-called Direct Simulation Monte Carlo (DSMC) method, which is known to qualify as an efficient and accurate method to numerically solve the Boltzmann equation. The aim of this paper is to solve the Boltzmann equation by means of the DSMC method for a gas subjected to the planar Couette flow. The motivation of this work is twofold. On the one hand, we want to test the reliability of the far from equilibrium results obtained from kinetic models and the Grad method by making a comparison with the Boltzmann solution in the case of Maxwell molecules, for which the form of the hydrodynamic profiles in the bulk region (far from the boundaries) is known. We will determine not only the hydrodynamic profiles but also the nonlinear transport coefficients and the velocity distribution function. On the other hand, we want to investigate whether the above results for a system of Maxwell molecules extend to other interaction potentials. This extension holds when the Boltzmann equation is replaced by kinetic models where, in terms of a conveniently scaled space variable, all the results are independent of the interaction law. Thus, we will also solve numerically the Boltzmann equation by the DSMC method for a hard-sphere gas. Since we are interested in describing transport properties in the bulk region, i.e., free from finite-size effects, we need to use appropriate boundary conditions in the simulations. In the conventional boundary conditions, the gas is assumed to be enclosed between two baths at equilibrium in relative motion and, in general, at different temperatures. Under these conditions, a particle leaving the system is formally replaced by a particle coming from the bath, so the incoming velocity is sampled from a local equilibrium distribution. As a consequence, there exists a mismatch between the velocity distribution function of the reemitted particles and that of those particles of the gas located near the wall and moving along the same direction. In this case, in order to inhibit the influence of boundary effects one needs to take very large systems (normal distance between the plates much larger than the mean free path), what is not practical from a computational point of view. To overcome this difficulty, a possibility is to assume that both baths are out of equilibrium in a state close to that of the actual gas. Since such a state is not known “a priori”, in this paper we assume that the state of the baths is described by the BGK solution of the planar Couette flow. Although the boundary effects are still unavoidable, we expect that the above mismatch between reemitted and gas particles will be much smaller. As a matter of fact, the use of these conditions has been shown to be appropriate to analyze bulk transport properties in the special case of planar Fourier flow (both walls at rest). The plan of the paper is as follows. In Sec. II we give a brief description of the planar Couette flow and a summary of the main results obtained from the Boltzmann equation and kinetic models. The boundary conditions used in the simulations and the DSMC method are described in Sec. III. Section IV presents the main results of the paper, where special attention is paid to the nonlinear transport coefficients. A comparison with the analytical results derived from the BGK and the ellipsoidal statistical (ES) models and from the Grad method is also carried out. The comparison shows in general a good agreement of the kinetic models with computer simulations, even for large shear rates. In addition, the velocity distribution function obtained from the simulation in the bulk domain is compared with the ones given by the kinetic models. It is shown again that the agreement is qualitatively good. We close the paper in Sec. V with some concluding remarks. ## II Description of the problem and summary of theoretical results Let us consider a dilute gas. In this case, a kinetic description is sufficient to characterize the state of the system by means of the velocity distribution function $`f(𝐫,𝐯;t)`$. This distribution function obeys the nonlinear Boltzmann equation, which in the absence of external forces is given by $$\frac{f}{t}+𝐯f=J[f,f],$$ (8) where $$J[f,f]=𝑑𝐯_1𝑑\widehat{𝐤}gI(g,\widehat{𝐤})\left[f(𝐯^{})f(𝐯_1^{})f(𝐯)f(𝐯_1)\right]$$ (9) is the collision operator. In this equation, $`I(g,\widehat{𝐤})`$ is the differential cross section, $`g\|𝐯𝐯_1\|`$ being the relative velocity, and $`(𝐯^{},𝐯_1^{})`$ are precollisional velocities yielding postcollisional velocities $`(𝐯,𝐯_1)`$. From the distribution function, one may define the hydrodynamic quantities $$n=𝑑𝐯f,$$ (10) $$𝐮=\frac{1}{n}𝑑𝐯𝐯f,$$ (11) $$\frac{3}{2}nk_BT=\frac{m}{2}𝑑𝐯𝐕^2f,$$ (12) and the momentum and heat fluxes $$𝖯=m𝑑𝐯\mathrm{𝐕𝐕}f,$$ (13) $$𝐪=\frac{m}{2}𝑑𝐯𝐕^2𝐕f.$$ (14) Here, $`n`$ is the number density, $`𝐮`$ is the flow velocity, $`T`$ is the temperature, $`𝖯`$ is the pressure tensor, $`𝐪`$ is the heat flux, and $`𝐕=𝐯𝐮`$ is the peculiar velocity. In addition, the equation of state is that of an ideal gas, i.e., $`p=nk_BT`$. Most of the known solutions to Eq. (8) for spatially inhomogeneous states correspond to the special case of Maxwell molecules, namely, a repulsive potential of the form $`V(r)r^4`$. For this potential, the collision rate $`gI(g,\widehat{𝐤})`$ is independent of the relative velocity and this allows the infinite hierarchy of velocity moments to be recursively solved in some specific situations. Furthermore, the NS transport coefficients $`\eta _0`$ and $`\kappa _0`$ can be exactly obtained from the Chapman-Enskog method. They are given by $$\eta _0=\frac{p}{\nu },\kappa _0=\frac{15}{4}\frac{k_B}{m}\eta _0,$$ (15) where $`\nu =\theta n`$, $`\theta `$ being an eigenvalue of the linearized Boltzmann collision operator. In this paper we are interested in studying the planar Couette flow for a dilute gas. We consider a gas enclosed between two parallel plates in relative motion and maintained at different temperatures. Under these conditions, the system is driven out of equilibrium by the combined action of the velocity and thermal gradients along the direction normal to the plates. After a transient period, the gas is expected to reach a steady state and the corresponding Boltzmann equation reads $$v_y\frac{}{y}f=J[f,f],$$ (16) where we have chosen the axis $`y`$ as the one normal to the plates. In general, this equation must be solved subjected to specific boundary conditions. Nevertheless, in the same spirit as in the Chapman-Enskog method, one may look for “normal” solutions in which all the space dependence of the distribution function occurs through a functional dependence on the hydrodynamic fields. The normal solution describes the state of the gas in the hydrodynamic regime, namely, for times much longer than the mean free time and for distances from the walls much larger than the mean free path. As mentioned in the Introduction, it has been proved that, in the case of Maxwell molecules, Eq. (16) admits an exact solution corresponding to the planar Couette flow problem. This solution belongs to the normal class, so that no explicit boundary conditions appear. In other words, all the space dependence of $`f`$ is given through the local density, the local velocity, the local temperature, and their gradients. The solution is characterized by hydrodynamic profiles that are a simple generalization of those predicted by the NS approximation, Eqs. (5)–(7), namely $$p=\text{const},$$ (17) $$\frac{1}{\nu (y)}\frac{u_x}{y}a=\text{const},$$ (18) $$\left[\frac{1}{\nu (y)}\frac{}{y}\right]^2T=\text{Pr}\frac{2m}{k_B}\gamma (a).$$ (19) The dimensionless parameter $`\gamma (a)`$ is a nonlinear function of the reduced shear rate $`a`$ that, by construction, behaves as $`\gamma a^2/5`$ in the limit $`a0`$. Again, the temperature can be seen as a quadratic function of the flow velocity, but now the coefficient $`\eta _0/\kappa _0`$ appearing in Eq. (7) is replaced by a shear-rate dependent coefficient $`(\eta _0/\kappa _0)5\gamma (a)/a^2`$. Furthermore, in this solution the pressure tensor is independent of the thermal gradient and the heat flux vector is linear in the thermal gradient, but all these fluxes are nonlinear functions of the shear rate. This nonlinear dependence can be characterized through five generalized transport coefficients. First, the shear stress $`P_{xy}`$ defines a generalized shear viscosity $`\eta (a)`$ as $$P_{xy}=\eta (a)\frac{u_x}{y}\eta _0F_\eta (a)\frac{u_x}{y}.$$ (20) Analogously, the component of the heat flux parallel to the thermal gradient defines a generalized thermal conductivity coefficient $`\kappa _{yy}(a)`$: $$q_y=\kappa _{yy}(a)\frac{T}{y}\kappa _0F_\kappa (a)\frac{T}{y}.$$ (21) The dimensionless functions $`F_\eta (a)`$ and $`F_\kappa (a)`$ are the most relevant quantities of the problem. They are related to the function $`\gamma (a)`$ by $`\gamma (a)=a^2F_\eta (a)/5F_\kappa (a)`$. Normal stress differences are different from zero and are measured by two viscometric functions $`\mathrm{\Psi }_{1,2}(a)`$ defined by $$\frac{P_{yy}P_{xx}}{p}=\mathrm{\Psi }_1(a)a^2,$$ (22) $$\frac{P_{zz}P_{yy}}{p}=\mathrm{\Psi }_2(a)a^2.$$ (23) Another interesting quantity related to the pressure tensor is the friction coefficient $`\mu (a)=P_{xy}/P_{yy}`$. To NS order, we simply have $`\mu (a)=a`$. In the non-Newtonian regime, we generalize this coefficient as $`\mu (a)=aF_\mu (a)`$, where the friction function $`F_\mu (a)`$ is $$F_\mu (a)=\frac{F_\eta (a)}{1[\mathrm{\Psi }_2(a)\mathrm{\Psi }_1(a)]a^2/3}.$$ (24) Finally, there exists a nonzero component of the heat flux orthogonal to the thermal gradient given by $$q_x=\kappa _{xy}(a)\frac{}{y}T\kappa _0\mathrm{\Phi }(a)a\frac{}{y}T.$$ (25) The three functions $`\mathrm{\Psi }_{1,2}(a)`$ and $`\mathrm{\Phi }(a)`$ are generalizations of Burnett coefficients. In fact, $`\mathrm{\Psi }_1(0)=14/5`$, $`\mathrm{\Psi }_2(0)=4/5`$, and $`\mathrm{\Phi }(0)=7/2`$ for Maxwell molecules and for hard spheres in the first Sonine approximation. The determination of the nonlinear transport coefficients $`F_\eta `$, $`F_\kappa `$, $`\mathrm{\Psi }_{1,2}`$, and $`\mathrm{\Phi }`$ would imply the solution of an infinite hierarchy that cannot be solved in a recursive way. This hierarchy can only be solved step by step when one performs a perturbation expansion in powers of the shear rate. In fact, Tij and Santos determined the solution up to super-Burnett order: $$F_\eta (a)=13.111a^2+𝒪(a^4),$$ (26) $$F_\kappa (a)=17.259a^2+𝒪(a^4).$$ (27) From Eqs. (26) and (27) it follows that $`\gamma (a)=(a^2/5)[1+4.148a^2+𝒪(a^4)]`$. In addition, $`F_\mu (a)=11.911a^2+𝒪(a^3)`$. Although the above analyses are valuable, they have two main limitations. On the one hand, they are restricted to the special case of Maxwell molecules. For other interaction potentials (e.g., hard spheres), the collisional moments involve all the moments of the distribution function and, as a consequence, the hydrodynamic profiles (17)–(19) are not strictly true. On the other hand, even for Maxwell molecules, the above perturbative solution is not useful for finite shear rates. These two limitations can be overcome, in the context of analytical methods, by introducing additional approximations, such as the Grad method, or by describing the system by means of kinetic models. In these approaches, one looks for a solution having the same hydrodynamic profiles as in the case of the Boltzmann equation, cf. Eqs. (17)–(19). As before, this solution describes the properties of the system in the bulk region, which is insensitive to the details of the boundary conditions. From the Grad method and from the BGK and ES kinetic models, one explicitly gets the full nonlinear shear-rate dependence of the transport coefficients in a non-perturbative way. Moreover, the results are universal, namely the functions $`F_\eta (a)`$, $`F_\kappa (a)`$, …are independent of the interaction potential, provided that the reduced shear rate is defined as in Eq. (18) with $`\nu =p/\eta _0`$. The thirteen-moment Grad method consists of replacing the actual distribution by $$ff_L\left\{1+\frac{m}{n(k_BT)^2}\left[\left(\frac{mV^2}{5k_BT}1\right)𝐕𝐪+\frac{1}{2}\left(P_{ij}p\delta _{ij}\right)V_iV_j\right]\right\},$$ (28) where $$f_L=n\left(\frac{m}{2\pi k_BT}\right)^{3/2}\mathrm{exp}\left(\frac{mV^2}{2k_BT}\right)$$ (29) is the local equilibrium distribution function. When the approximation (28) is inserted into the Boltzmann equation (8) and velocity moments are taken, one gets a closed set of equations for $`n`$, $`𝐮`$, $`𝖯`$, and $`𝐪`$. According to the geometry of the planar Couette flow, there are eight independent moments instead of the original thirteen moments appearing in Eq. (28). Risso and Cordero found that the set of independent moment equations, neglecting nonlinear terms in the fluxes, admits a solution consistent with the profiles (17)–(19). In addition, they obtained explicit expressions for the transport coefficients introduced above as nonlinear functions of the reduced shear rate $`a`$. Those expressions are displayed in the Appendix. Now we consider the kinetic model approach. The basic idea is to replace the detailed Boltzmann collision operator by a much simpler term that otherwise retains the main physical properties of the original operator $`J[f,f]`$. The most familiar choice is the BGK model, $$J[f,f]\nu (ff_L),$$ (30) where $`\nu `$ is an effective collision frequency that depends on the temperature, according to the interaction potential. The NS transport coefficients of the BGK model are $`\eta _0=p/\nu `$ and $`\kappa _0=5pk_B/2m\nu `$. Thus the BGK model has the drawback that it predicts an incorrect value for the Prandtl number, $`\text{Pr}=1`$. This is a consequence of the fact that $`\nu `$ is the only adjustable parameter in the model. This deficiency is corrected by the so-called ellipsoidal statistical (ES) model, in which case $$J[f,f]\zeta (ff_R),$$ (31) where $`\zeta `$ is again an effective collision frequency and the reference function $`f_R`$ is $$f_R=n\pi ^{3/2}(det𝖠)^{1/2}\mathrm{exp}\left(𝖠^1:\mathrm{𝐕𝐕}\right),$$ (32) where $`A_{ij}=(2k_BT/m)\text{Pr}^1\delta _{ij}2(\text{Pr}^11)P_{ij}/mn`$. The NS coefficients are now $`\eta _0=p/(\zeta \text{Pr}^1)`$ and $`\kappa _0=5pk_B/2m\zeta `$. If, as in Eq. (15), we define a collision frequency $`\nu =p/\eta _0`$, then $`\nu =\zeta \text{Pr}^1`$ in the ES model. Note that the ES model reduces to the BGK model if we formally make $`\text{Pr}=1`$. Therefore, the ES model can be seen as an extension of the BGK model to account for the correct Prandtl number, which can be particularly relevant in the Couette flow where heat transport and momentum transport are coupled. The results derived from the BGK and the ES models for the Couette flow problem are also given in the Appendix. Apart from obtaining the transport properties, the velocity distribution function can be explicitly written. In particular, the BGK distribution is given by $`f(𝐫,𝐯)`$ $`=`$ $`n\left({\displaystyle \frac{m}{2\pi k_BT}}\right)^{3/2}{\displaystyle \frac{2\alpha (1+\alpha )^{3/2}}{ϵ\|\xi _y\|}}{\displaystyle _{t_0}^{t_1}}𝑑t\left[2t(1\alpha )t^2\right]^{5/2}`$ (34) $`\times \mathrm{exp}\left\{{\displaystyle \frac{2\alpha }{1+\alpha }}{\displaystyle \frac{1t}{ϵ\xi _y}}{\displaystyle \frac{1+\alpha }{2t(1\alpha )t^2}}\left[\left(\xi _x+{\displaystyle \frac{2a\alpha }{1+\alpha }}{\displaystyle \frac{1t}{ϵ}}\right)^2+\xi _y^2+\xi _z^2\right]\right\}.`$ Here, $`(t_0,t_1)=(0,1)`$ if $`\xi _y>0`$ and $`(t_0,t_1)=[1,2/(1\alpha )]`$ if $`\xi _y<0`$. Besides, $`𝝃(m/2k_BT)^{1/2}𝐕`$, $$\alpha =\frac{ϵ}{\left(ϵ^2+8\gamma \right)^{1/2}},$$ (35) and $$ϵ=\left(\frac{2k_BT}{m}\right)^{1/2}\frac{1}{\nu }\frac{}{y}\mathrm{ln}T$$ (36) is a (local) reduced thermal gradient. Equation (34) shows that the distribution function is a highly nonlinear function of the reduced gradients $`a`$ and $`ϵ`$. In Ref. a comparison between the analytical results obtained from the Grad method and the BGK and ES models with those obtained from molecular dynamics simulations for hard disks was carried out. It was found that the three theories reproduced quite well the simulation data for $`F_\eta `$, but the Grad method failed for $`F_\kappa `$ and $`\mathrm{\Phi }`$. The latter quantity was reproduced better by the ES model than by the BGK model. Notwithstanding this, more definite conclusions could not be reached because the simulation data were restricted to rather small shear rates, namely $`a0.2`$. In this range of shear rates, non-Newtonian effects are not especially significant. In addition, although the molecular dynamics simulations correspond to a very dilute gas (area fraction $`\varphi 1\%`$), the collisional contributions to the transport coefficients (absent in a Boltzmann description) are not strictly zero. Finally, as a minor point, conclusions drawn in the context of two-dimensional systems should not be extrapolated without caution to the more realistic case of three dimensions. As said in the Introduction, the aim of this paper is to solve numerically the Boltzmann equation by means of the DSMC method for the planar Couette flow and compare the results with the theoretical predictions. We will consider three-dimensional systems of Maxwell molecules and hard spheres subjected to shear rates as large as $`a1.2`$. In addition to the nonlinear transport coefficients, the comparison will be extended to the level of the velocity distribution function itself. ## III Boundary conditions and Monte Carlo simulation ### A Boundary conditions The goal now is to solve numerically the Boltzmann equation corresponding to the planar Couette flow by using the successful DSMC method. The gas is enclosed between two parallel plates located at $`y=0`$ and $`y=L`$, which are moving along the $`x`$-direction with velocities $`𝐔_0=U_0\widehat{𝐱}`$ and $`𝐔_L=U_L\widehat{𝐱}`$, respectively. In addition, they are kept at temperatures $`T_0`$ and $`T_L`$, respectively. In order to solve Eq. (16), we need to impose the corresponding boundary conditions. They can be expressed in terms of the kernels $`K_{0,L}(𝐯,𝐯^{})`$ defined as follows. When a particle with velocity $`𝐯^{}`$ hits the wall at $`y=L`$, the probability of being reemitted with a velocity $`𝐯`$ within the range $`d𝐯`$ is $`K_L(𝐯,𝐯^{})d𝐯`$; the kernel $`K_0(𝐯,𝐯^{})`$ represents the same but at $`y=0`$. The boundary conditions are then $`\mathrm{\Theta }(\pm v_y)\|v_y\|f(y=\{0,L\},𝐯)`$ $`=`$ $`\mathrm{\Theta }(\pm v_y){\displaystyle 𝑑𝐯^{}\|v_y^{}\|K_{0,L}(𝐯,𝐯^{})}`$ (38) $`\times \mathrm{\Theta }(v_y^{})f(y=\{0,L\},𝐯^{},t).`$ In this paper we consider boundary conditions of complete accommodation with the walls, so that $`K_{0,L}(𝐯,𝐯^{})=K_{0,L}(𝐯)`$ does not depend on the incoming velocity $`𝐯^{}`$ and can be written as $$K_{0,L}(𝐯)=A_{0,L}^1\mathrm{\Theta }(\pm v_y)\|v_y\|\varphi _{0,L}(𝐯),A_{0,L}=𝑑𝐯\mathrm{\Theta }(\pm v_y)\|v_y\|\varphi _{0,L}(𝐯),$$ (39) where $`\varphi _{0,L}(𝐯)`$ represents the probability distribution of a fictitious gas in contact with the system at $`y=\{0,L\}`$. Equation (39) can then be interpreted as meaning that when a particle hits a wall, it is absorbed and then replaced by a particle leaving the fictitious bath. Of course, any choice of $`\varphi _{0,L}(𝐯)`$ must be consistent with the imposed wall velocities and temperatures, i.e., $$U_{0,L}=𝑑𝐯v_x\varphi _{0,L}(𝐯),$$ (40) $$k_BT_{0,L}=\frac{1}{3}m𝑑𝐯(𝐯𝐔_{0,L})^2\varphi _{0,L}(𝐯).$$ (41) The simplest and most common choice is that of a Maxwell-Boltzmann (MB) distribution: $$\varphi _{0,L}^{\text{MB}}(𝐯)=\left(\frac{m}{2\pi k_BT_{0,L}}\right)^{3/2}\mathrm{exp}\left[\frac{m(𝐯𝐔_{0,L})^2}{2k_BT_{0,L}}\right].$$ (42) Under these conditions, the system is understood to be enclosed between two independent baths at equilibrium. While the conditions (42) are adequate for analyzing boundary effects, they are not very efficient when one is interested in obtaining the transport properties in the bulk region. In order to inhibit the influence of boundary effects, it is more convenient to imagine that the two fictitious baths are in nonequilibrium states resembling the state of the actual gas near the walls. More specifically, we can assume that the fictitious gases are described by the BGK equation, whose exact solution for the steady planar Couette flow is given by Eq. (34). In this case, the probability distributions $`\varphi _{0,L}(𝐯)`$ are $`\varphi _{0,L}^{\text{BGK}}(𝐯)`$ $`=`$ $`\pi ^{3/2}{\displaystyle \frac{m}{k_BT_{0,L}}}{\displaystyle \frac{\alpha _{0,L}(1+\alpha _{0,L})^{3/2}}{ϵ_{0,L}\|v_y\|}}{\displaystyle _{t_0}^{t_1}}𝑑t\left[2t(1\alpha _{0,L})t^2\right]^{5/2}`$ (45) $`\times \mathrm{exp}\{\left({\displaystyle \frac{2k_BT_{0,L}}{m}}\right)^{1/2}{\displaystyle \frac{2\alpha _{0,L}}{1+\alpha _{0,L}}}{\displaystyle \frac{1t}{ϵ_{0,L}v_y}}{\displaystyle \frac{m}{2k_BT_{0,L}}}{\displaystyle \frac{1+\alpha _{0,L}}{2t(1\alpha _{0,L})t^2}}`$ $`\times [(v_xU_{0,L}+{\displaystyle \frac{2a^{}\alpha _{0,L}}{1+\alpha _{0,L}}}{\displaystyle \frac{1t}{ϵ_{0,L}}})^2+v_y^2+v_z^2]\},`$ where $`(t_0,t_1)=(0,1)`$ if $`v_y>0`$ and $`(t_0,t_1)=[1,2/(1\alpha _{0,L})]`$ if $`v_y<0`$, and $`\alpha _{0,L}=ϵ_{0,L}/[ϵ_{0,L}^2+8\gamma _{\text{BGK}}(a^{})]^{1/2}`$. Here, $`a^{}`$ is the estimated value of the reduced shear rate, as predicted by the BGK model for specific values of the boundary parameters $`U_{0,L}`$ and $`T_{0,L}`$, and $`\gamma _{\text{BGK}}(a^{})`$ is obtained from Eq. (68). Given the values of the four independent boundary parameters $`U_{0,L}`$ and $`T_{0,L}`$ (as well as the distance $`L`$), the shear rate $`a^{}`$ and the local thermal gradients $`ϵ_{0,L}`$ are fixed by the conditions (17)–(19). Therefore, $$a^{}=\frac{U_LU_0}{},$$ (46) $$ϵ_{0,L}=\frac{1}{}\left(\frac{2k_BT_{0,L}}{m}\right)^{1/2}\left[\frac{T_LT_0}{T_{0,L}}\pm \frac{m\gamma _{\text{BGK}}(a^{})\text{Pr}}{k_BT_{0,L}}^2\right].$$ (47) In these equations, $``$ is related to the actual separation $`L`$ between the plates through the nonlinear equation $$L=_0^{}\frac{ds}{\nu (s)},$$ (48) where $`s`$ is a variable in terms of which the temperature is a quadratic function, namely $`T(s)=T_0[1+ϵ_0(m/2k_BT_0)^{1/2}sm\gamma _{\text{BGK}}(a^{})\text{Pr}s^2/k_BT_0]`$, and the $`s`$-dependence of the collision frequency $`\nu (s)`$ appears only through the temperature (taking into account that $`p=\text{const}`$). The solution of the nonlinear set of equations (46)–(48) gives $`a^{}`$, $`ϵ_{0,L}`$, and $``$ for any choice of $`U_{0,L}`$, $`T_{0,L}`$, and $`L`$. Nevertheless, from a practical point of view, it is more convenient to fix $`U_0`$, $`T_{0,L}`$, $`a^{}`$, and $`ϵ_0`$ as independent parameters. Without loss of generality we take $`U_0=0`$ and $`T_L=1`$. In addition, we will choose $`ϵ_0=0`$. This implies that, if boundary effects were absent, the temperature would have a maximum at the lower plate $`y=0`$. Thus, $$U_L=a^{},=\left[\frac{\mathrm{\Delta }}{2\text{Pr}\gamma _{\text{BGK}}(a^{})}\right]^{1/2},$$ (49) $$ϵ_L=2\frac{\mathrm{\Delta }}{}.$$ (50) In the above equations $`\mathrm{\Delta }T_01`$ and we have taken $`m=1`$ and $`k_B=\frac{1}{2}`$. Finally, the actual distance $`L`$ is given by Eq. (48). However, from the simulation point of view, it is more useful to express $`L`$ in terms of the collision frequency $`\overline{\nu }`$ corresponding to the temperature $`T_L`$ and the average density $`\overline{n}`$. In other words, instead of Eq. (48) we use $$\overline{n}=\frac{1}{L}_0^L𝑑yn(y)=\frac{1}{L}_0^{}𝑑s\frac{n(s)}{\nu (s)}$$ (51) to determine $`L`$. For the sake of concreteness, let us consider repulsive potentials, for which $`\nu =\overline{\nu }(n/\overline{n})(T/T_L)^\omega `$, where $`\omega `$ ranges from 0 (Maxwell molecules) to $`\frac{1}{2}`$ (hard spheres). In that case, $$L=\frac{1}{\overline{\nu }}_0^{}𝑑s[T(s)]^\omega .$$ (52) For Maxwell molecules, this simply reduces to $`L=/\overline{\nu }`$, while for hard spheres one has $`L=(/\overline{\nu })\mathrm{tan}^1(\sqrt{\mathrm{\Delta }})/\sqrt{\mathrm{\Delta }}`$. In summary, given the values of $`a^{}`$ and $`\mathrm{\Delta }`$, the separation $`L`$ between the plates and the velocity of the upper plate $`U_L`$ are uniquely determined. In addition, the thermal gradient at the upper plate $`ϵ_L`$ is also determined, while the value at the lower plate is fixed as $`ϵ_0=0`$. The knowledge of these boundary parameters allows one to obtain the distributions $`\varphi _{0,L}^{\text{BGK}}(𝐯)`$, according to Eq. (45). ### B The DSMC method Now, we briefly describe the numerical algorithm we have employed to solve the Boltzmann equation by means of the so-called Direct Simulation Monte Carlo (DSMC) method. In this method, the velocity distribution function is represented by the velocities $`\{𝐯_i\}`$ and positions $`\{y_i\}`$ of a sufficiently large number of particles $`N`$. Given the geometry of the problem, the physical system is split into layers of width $`\delta y`$, sufficiently smaller than the mean free path. The velocities and coordinates are updated from time $`t`$ to time $`t+\delta t`$, where the time step $`\delta t`$ is much smaller than the mean free time, by applying a streaming step followed by a collision step. In the streaming step, the particles are moved ballistically, $`y_iy_i+v_{iy}\delta t`$. In addition, those particles crossing the boundaries are reentered with velocities sampled from the corresponding probability distribution $`K_{0,L}(𝐯)`$. Suppose that a particle $`i`$ crosses the lower plate between times $`t`$ and $`t+\delta t`$, i.e., $`y_i+v_{iy}\delta t<0`$. Then, regardless of the incoming velocity $`𝐯_i`$, a new velocity $`\stackrel{~}{𝐯}_i`$ is assigned according to the following rules. First, a velocity $`\stackrel{~}{𝐯}_i`$ (with $`\stackrel{~}{v}_{iy}>0`$) is sampled with a probability proportional to $`\|v_y\|\varphi _0^{\text{MB}}(𝐯)`$. If one is considering the MB boundary conditions, this velocity is accepted directly. Otherwise, the above acts as a “filter” to optimize the acceptance-rejection procedure and the velocity $`\stackrel{~}{𝐯}_i`$ is accepted with a probability proportional to the ratio $`\varphi _0^{\text{BGK}}(𝐯)/\varphi _0^{\text{MB}}(𝐯)`$. If the velocity is rejected, a new velocity $`\stackrel{~}{𝐯}_i`$ is sampled and the process is repeated until acceptance. The new position is assigned as $`\stackrel{~}{v}_{iy}(\delta t+y_i/v_{iy})`$. The process is analogous in the case of the upper plate. The collision step proceeds as follows For each layer $`\alpha `$, a pair of potential collision partners $`i`$ and $`j`$ are chosen at random with equiprobability. The collision between those particles is then accepted with a probability equal to the corresponding collision rate times $`\delta t`$. For hard spheres, the collision rate is proportional to the relative velocity $`\|𝐯_i𝐯_j\|`$, while it is independent of the relative velocity for Maxwell molecules (an angle cut-off is needed in the latter case). If the collision is accepted, the scattering direction is randomly chosen according to the interaction law and post-collisional velocities are assigned to both particles, according to the conservation of momentum and energy. After the collision is processed or if the pair is rejected, the routine moves again to the choice of a new pair until the required number of candidate pairs has been taken. In the course of the simulations, the following “coarse-grained” local quantities are computed. The number density in layer $`\alpha `$ is $$n_\alpha =\overline{n}\frac{N_\alpha }{(N/L)\delta y}=\frac{\overline{n}L}{N\delta y}\underset{i=1}{\overset{N}{}}\mathrm{\Theta }_\alpha (y_i),$$ (53) where $`\mathrm{\Theta }_\alpha (y)`$ is the characteristic function of layer $`\alpha `$, i.e., $`\mathrm{\Theta }_\alpha (y)=1`$ if $`y`$ belongs to layer $`\alpha `$ and is zero otherwise. Similarly, the flow velocity, the temperature, the pressure tensor, and the heat flux of layer $`\alpha `$ are $$𝐮_\alpha =\frac{1}{N_\alpha }\underset{i=1}{\overset{N}{}}\mathrm{\Theta }_\alpha (y_i)𝐯_i,$$ (54) $$k_BT_\alpha =\frac{p_\alpha }{n_\alpha }=\frac{m}{3N_\alpha }\underset{i=1}{\overset{N}{}}\mathrm{\Theta }_\alpha (y_i)(𝐯_i𝐮_\alpha )^2,$$ (55) $$𝖯_\alpha =\frac{L}{N\delta y}m\underset{i=1}{\overset{N}{}}\mathrm{\Theta }_\alpha (y_i)(𝐯_i𝐮_\alpha )(𝐯_i𝐮_\alpha ),$$ (56) $$𝐪_\alpha =\frac{L}{N\delta y}\frac{m}{2}\underset{i=1}{\overset{N}{}}\mathrm{\Theta }_\alpha (y_i)(𝐯_i𝐮_\alpha )^2(𝐯_i𝐮_\alpha ).$$ (57) From these quantities one can get local values of the gradients and of the transport coefficients. For instance, the reduced shear rate is $$a_\alpha =\frac{\overline{n}}{\overline{\nu }n_\alpha }\left(\frac{T_L}{T_\alpha }\right)^\omega \frac{u_{\alpha +1,x}u_{\alpha ,x}}{\delta y}$$ (58) and the viscosity function is $$F_{\eta ,\alpha }=\frac{P_{\alpha ,xy}}{a_\alpha p_\alpha }.$$ (59) As said before, in the simulations we take units such that $`m=1`$, $`T_L=1`$, and $`k_B=\frac{1}{2}`$. It remains to fix the time unit or, equivalently, the length unit. The standard definition of mean free path in the case of hard spheres is $$\lambda =\frac{1}{\sqrt{2}n\pi \sigma ^2},$$ (60) where $`\sigma `$ is the diameter of the spheres. The Navier-Stokes shear viscosity is (in the first Sonine approximation) $`\eta _0=5(mk_BT/\pi )^{1/2}/16\sigma ^2`$. Consequently, the effective collision frequency $`\nu =p/\eta _0`$ and the mean free path $`\lambda `$ are related as $`\nu =(8/5\sqrt{\pi })(2k_BT/m)^{1/2}/\lambda `$. As usual, we choose the mean free path corresponding to the average density $`\overline{n}`$ as the length unit. This in turn implies that $`\overline{\nu }=8/5\sqrt{\pi }0.903`$. For convenience, we take the latter value for Maxwell molecules as well. The typical values of the simulation parameters are $`N=2\times 10^5`$ particles, a layer width $`\delta y=0.02`$, and a time step $`\delta t=0.003`$. The procedure to measure the relevant quantities of the problem is as follows. First the values of the imposed shear rate $`a^{}`$ and the temperature difference $`\mathrm{\Delta }`$ are chosen. This choice fixes the system size $`L`$, as well as the upper velocity $`U_L`$ and the upper thermal gradient $`ϵ_L`$, according to Eqs. (49), (50), and (52). Starting from an equilibrium initial state with $`T=T_0`$, the system evolves driven by the boundary conditions described before. After a transient period (typically up to $`t=25`$), the system reaches a steady state in which the quantities fluctuate around constant values. In this state, the balance equations predict that the quantities $`u_y`$, $`P_{xy}`$, $`P_{yy}`$, and $`u_xP_{xy}+q_y`$ are spatially uniform and this is used in the simulations as a test to make sure that the steady state has been achieved. Once the steady state is reached, the local quantities (53)–(59) are averaged over typically 100 snapshots equally spaced between $`t=25`$ and $`t=55`$, which corresponds to about 60 collisions per particle in the case of hard spheres. ### C Test of the numerical algorithm Before closing this Section, it is worthwhile carrying out a test of the reliability of the numerical method. To that end, we have simulated the BGK equation by a DSMC-like method similar to the one described in Ref. . If the boundary conditions are implemented correctly and the simulation parameters are well chosen, then the simulation results should agree with the theoretical BGK predictions. We have checked that this indeed the case. As an example, consider the hard-sphere situation with $`a^{}=1`$ and $`\mathrm{\Delta }=5`$. This corresponds to $`\gamma _{\text{BGK}}=0.248`$, $`L=1.81`$, $`U_L=3.17`$, and $`ϵ_L=3.15`$, where in this case $`\text{Pr}=1`$. Figure 1 shows the marginal velocity distributions of particles reemitted by the walls, $$𝒦_{0,L}(\xi _y)=\left(\frac{2k_BT_{0,L}}{m}\right)^{1/2}_{\mathrm{}}^{\mathrm{}}𝑑v_x_{\mathrm{}}^{\mathrm{}}𝑑v_zK_{0,L}(𝐯),$$ (61) as functions of $`\xi _y=(m/2k_BT_{0,L})^{1/2}v_y`$. The case $`\xi _y>0`$ ($`\xi _y<0`$) corresponds to particles that are reemitted from the lower (upper) plate. The agreement with the imposed distribution is excellent. Note the strong asymmetry between the distributions corresponding to both plates, in contrast to the symmetry of the MB distributions obtained from (42). The temperature and velocity profiles are shown in Fig. 2. The simulation values overlap, within statistical fluctuations, with the theoretical predictions. Apart from the profiles, we have verified that the generalized transport coefficients obtained from the fluxes also agree with the theory. A more stringent test is provided in Fig. 3, where the marginal distribution function $$\phi (\xi _y)=\frac{1}{n}\left(\frac{2k_BT}{m}\right)^{1/2}_{\mathrm{}}^{\mathrm{}}𝑑v_x_{\mathrm{}}^{\mathrm{}}𝑑v_zf(𝐯),\xi _y=\left(\frac{m}{2k_BT}\right)^{1/2}v_y$$ (62) is plotted at the point $`y/L=0.5`$, which corresponds to a local thermal gradient $`ϵ=0.60`$. It is apparent again that an excellent agreement exists between simulation and theory. ## IV Results This Section is devoted to a comparison between the simulation results for the Boltzmann equation obtained by the simulation method described in the previous Section and the theoretical predictions provided by the Grad method and the BGK and ES kinetic models. The comparison will be carried out at the levels of the transport coefficients and the velocity distribution, both for Maxwell molecules and hard spheres. Before that, the hydrodynamic profiles obtained from the simulations by using the two types of boundary conditions considered in Sec. III are presented. ### A Hydrodynamic profiles As said in Secs. I and II, the Boltzmann equation for Maxwell molecules admits an exact solution characterized by Eqs. (17)–(19). This solution applies to the bulk region, i.e., the region where boundary effects are negligible. Obviously, in a simulation with a finite size of the system, it is not possible to avoid boundary effects completely. On the other hand, one can expect that the “nonequilibrium” boundary conditions based on the BGK distribution, Eq. (45), inhibit the boundary effects, as compared with the conventional “equilibrium” boundary conditions (42). We have confirmed that this is indeed the case. As an illustrative example, let us consider $`\mathrm{\Delta }=5`$ and $`a^{}=0.92`$ for Maxwell molecules. This corresponds to $`\gamma _{\text{BGK}}=0.21`$, $`L=4.68`$, $`U_L=3.90`$, and $`ϵ_L=2.37`$. Figure 4 shows the temperature and velocity profiles as obtained by using the MB and BGK boundary conditions. It is apparent that the velocity slips and the temperature jumps at the walls are much larger in the former case than in the latter. Note that the maximum temperature is not exactly located in the layer adjacent to the lower plate but it is slightly shifted. It is interesting to remark that when plotting the temperature $`T`$ as a function of the flow velocity $`u_x`$, a parabolic curve is observed with both boundary conditions, as expected from Eqs. (18) and (19). A more evident proof of the advantage of the BGK conditions is shown in Fig. 5. Since the pressure is a constant in the exact solution valid in the bulk, any deviation from $`p=\text{const}`$ can be attributed to boundary effects. The pressure obtained with the BGK boundary conditions is practically constant, except near the upper plate, while the one obtained with the MB conditions is only nearly constant in a small region around $`y/L0.75`$. Finally, Fig. 6 shows the ratio $`a/a^{}`$ between the actual (local) shear rate $`a`$ measured in the simulations, cf. Eq. (58), and the imposed shear rate $`a^{}`$. The ratio is closer to 1 in the case of the BGK boundary conditions than in that of the MB conditions. In addition, in the former case the region where $`a`$ is practically constant extends to layers closer to the lower plate. In summary, the above example illustrates that the BGK boundary conditions are much more efficient than the MB ones to measure transport properties in the bulk. Therefore, in what follows we will only consider the BGK conditions. In each case, we identify a bulk domain comprised between the layers $`y=y_0`$ and $`y=y_1`$ where $`a\text{const}`$, $`p\text{const}`$, and $`\gamma \text{const}`$, and take averages of $`a`$, $`F_\eta `$, $`F_\kappa `$, $`\mathrm{\Psi }_{1,2}`$, and $`\mathrm{\Phi }`$ over those layers. Typical values are $`y_0/L0.2`$ and $`y_1/L0.8`$. ### B Nonlinear transport coefficients In this subsection we compare the simulation results for Maxwell molecules and hard spheres obtained from the Monte Carlo simulations with the (universal) predictions of the Grad method and the BGK and ES kinetic models. As said in the Introduction, we are interested in situations where the Knudsen number is not small, so that nonlinear effects are important. An interesting quantity in the nonlinear Couette flow is the parameter $`\gamma `$, which is related to the curvature of the temperature profile. Its shear-rate dependence is shown in Fig. 7. A remarkable feature is that the simulation data for both interactions seem to lie on a common curve. This indicates that, as predicted by the models, the transport properties are hardly sensitive to the interaction potential, provided that the quantities are conveniently scaled. While the kinetic models exhibit a good quantitative (in the case of the ES model) or qualitative (BGK model) agreement, the Grad method fails, except for small shear rates. Since $`\gamma (a)=a^2F_\eta (a)/5F_\kappa (a)`$, one can interpret $`5\gamma (a)/a^2`$ as an effective, shear-rate dependent Prandtl number (relative to the usual Pr). This quantity is bounded between 1 for small shear rates and $`5/3`$ (in the BGK model) or $`5/2`$ (in the ES model) in the limit of large shear rates. This is consistent with the fact that the BGK model underestimates the value of $`\gamma `$. The most important transport coefficient is the nonlinear viscosity represented by the function $`F_\eta (a)`$. According to Fig. 8, the three theories retain the qualitative trends of the simulation data, namely the decrease of $`F_\eta `$ with increasing $`a`$ (shear thinning effect). In general, however, the kinetic models (especially the BGK model) have a better agreement than the Grad method. It is also interesting to remark that a slight influence of the interaction potential seems to exist, the shear thinning effect being a little bit more significant for hard spheres than for Maxwell molecules. The nonlinear thermal conductivity $`F_\kappa (a)`$ is plotted in Fig. 9. It is quite apparent that Grad’s solution does not capture even the qualitative shape of $`F_\kappa `$, as was already noted in the case of hard disks. Again, the kinetic models present a good agreement, especially in the case of the ES model. Normal stress differences are characterized by the viscometric functions $`\mathrm{\Psi }_{1,2}(a)`$. These quantities are well described by the kinetic models, as shown in Figs. 10 and 11. At a quantitative level, the agreement is better in the case of Maxwell molecules. In fact, the viscometric functions, especially the second one, exhibit a certain influence of the potential. Although the functions $`\mathrm{\Psi }_{1,2}(a)`$ were not evaluated from the Grad method in Refs. and , the friction function $`F_\mu (a)`$ was calculated. This quantity is plotted in Fig. 12, showing an agreement between the theories and the simulation data similar to the one found in Fig. 8 for the viscosity function. The last transport coefficient is the cross-coefficient $`\mathrm{\Phi }`$ defined by Eq. (25). The comparison with simulation results for this quantity is a stringent test of the theories, since it is a generalization of a Burnett coefficient that measures the component of the heat flux orthogonal to the thermal gradient. Figure 13 shows that, as happened with the thermal conductivity function $`F_\kappa (a)`$, the Grad method gives a wrong prediction for the shear-rate dependence of $`\mathrm{\Phi }(a)`$. On the other hand, the kinetic models describe fairly well the nonlinear behavior of this function. In the case of the ES model, the agreement with the simulation data is practically perfect. ### C Velocity distribution function Apart from the transport coefficients, the kinetic models provide the velocity distribution function. In the case of the BGK model, the solution is given by Eq. (34), while the ES distribution function can be found in Refs. and . In order to assess their reliability, we have computed in the simulations the marginal distribution (62) at the layer $`y/L=0.5`$. The ratio $`R(\xi _y)\phi (\xi _y)/[\pi ^{1/2}\mathrm{exp}(\xi _y^2)]`$, where $`\phi (\xi _y)`$ is defined by Eq. (62), is a measure of the departure of the distribution function from the local equilibrium. This quantity is plotted in Fig. 14 for Maxwell molecules in the case of a reduced shear rate $`a=0.636`$ and a reduced (local) thermal gradient $`ϵ=0.272`$. Both theories capture the main features of the actual distribution. While the BGK distribution exhibits a better agreement near the maximum (around $`\xi _y0.5`$), the ES distribution seems to describe better the region $`\xi _y1`$. The case of hard spheres is illustrated in Fig. 15, which corresponds to $`a=0.419`$ and $`ϵ=0.195`$. In this case, the ES model shows a superiority over the BGK model both near the maximum and for large positive velocities. ## V Concluding remarks This paper has dealt with the steady planar Couette flow in a dilute gas beyond the scope of the Navier-Stokes description. This nonlinear problem had been previously studied by means of kinetic theory tools, such as the Grad method, and the BGK and ES kinetic models. These theories predict momentum and heat fluxes characterized by five shear-rate dependent generalized transport coefficients: a viscosity function $`F_\eta (a)`$, Eq. (20), a thermal conductivity function $`F_\kappa (a)`$, Eq. (21), two viscometric functions $`\mathrm{\Psi }_{1,2}(a)`$, Eqs. (22) and (23), and a cross coefficient $`\mathrm{\Phi }(a)`$, Eq. (25). The main motivation of our study has been to perform DSMC simulations for Maxwell molecules and hard spheres in order to assess the reliability of the above theories in the non-Newtonian regime. Since we have been interested in the bulk properties, we have used “nonequilibrium” boundary conditions to inhibit the influence of finite-size effects. An important outcome of the simulation results is that, as predicted by the kinetic theories considered here, the shear-rate dependence of the transport coefficients is practically insensitive to the interaction potential when the quantities are properly nondimensionalized. In particular, the actual shear rate has been reduced with respect to an effective collision frequency defined from the Navier-Stokes shear viscosity coefficient. The simulation results show, however, that the second viscometric function presents a non negligible influence of the interaction model, so that the normal stress difference $`P_{zz}P_{yy}`$ is smaller for Maxwell molecules than for hard spheres. The comparison with the theoretical predictions shows that the kinetic models give a fairly good description of the five transport coefficients. On the other hand, the Grad method yields a shear viscosity in qualitative agreement with the simulations but it dramatically fails for the coefficients measuring the heat flux. This is basically due to the truncation scheme of the Grad method at the level of the heat flux. The physical idea behind a kinetic model is quite different, since it consists of replacing the true Boltzmann collision operator by a simple relaxation term but otherwise all the velocity moments are taken into account. As a consequence, while in the Grad method one has to solve a closed set of coupled differential equations for the moments, in the case of the kinetic model one gets the velocity distribution function and determines the fluxes from it. In the ES kinetic model the reference distribution function appearing in the collision operator is an anisotropic Gaussian parameterized by the pressure tensor. This allows the model to give the correct Prandtl number $`\text{Pr}=\frac{2}{3}`$ but at the expense of complicating its mathematical structure. In the case of the BGK model, however, the reference distribution is that of local equilibrium but the model leads to $`\text{Pr}=1`$. The agreement with simulation of the ES model is generally better than that of the BGK model, especially in the case of $`F_\kappa `$ and $`\mathrm{\Phi }`$. In spite of this, it is fair to say that the performance of the BGK model is quite good, given its simplicity relative to that of the ES model. Finally, the results reported in this paper clearly shows the usefulness of kinetic models to analyze nonlinear transport phenomena in the Couette flow problem. This complements previous conclusions drawn from other nonlinear problems, such as the uniform shear flow and the Fourier flow. ###### Acknowledgements. The authors acknowledge partial support from the DGES (Spain) through grant No. PB97-1501 and from the Junta de Extremadura (Fondo Social Europeo) through grant No. IPR99C031. ## Theoretical expressions for the transport coefficients In this Appendix we list the explicit shear-rate dependence of the dimensionless transport coefficients defined in Sec. II, according to the Grad method, the BGK model, and the ES model. ### 1 The Grad method From the Appendix of Ref. , corrected in Ref. , one has $$F_\eta (a)=\frac{2}{1+\frac{72}{25}a^2+\mathrm{\Delta }(a)},$$ (63) $$F_\kappa (a)=\frac{4}{1\frac{216}{25}a^2+3\mathrm{\Delta }(a)},$$ (64) $$\mathrm{\Phi }(a)=7\frac{1\frac{36}{125}a^2}{1+\frac{6}{5}a^2+(1\frac{63}{25}a^2)\mathrm{\Delta }(a)},$$ (65) where $`\mathrm{\Delta }(a)\sqrt{1+\frac{116}{25}a^2\frac{864}{625}a^4}`$. The viscometric functions are not evaluated in Refs. and , although the friction function is provided. It is given by $$F_\mu (a)=\frac{2}{1+\frac{12}{25}a^2+\mathrm{\Delta }(a)}.$$ (66) Note that $`F_\eta (a)`$ and $`F_\mu (a)`$ become meaningless for $`a^225(\sqrt{1057}+29)/4323.56`$, while $`F_\kappa (a)`$ and $`\mathrm{\Phi }(a)`$ are unphysical for $`a^250/630.79`$. For small shear rates, the above transport coefficients behave as $`F_\eta 1\frac{13}{5}a^2`$, $`F_\kappa 1+\frac{21}{50}a^2`$, $`\mathrm{\Phi }\frac{7}{2}(1\frac{197}{250}a^2)`$, and $`F_\mu 1\frac{7}{5}a^2`$. ### 2 The BGK model The derivation of the transport coefficients from the BGK model implies the resummation of asymptotic series by means of the Borel method. As a consequence, the results are expressed in terms of the functions $`F_r(x)`$ defined by the recurrence relation $`F_r(x)=[(d/dx)x]^rF_0(x)`$, where $$F_0(x)=\frac{2}{x}_0^{\mathrm{}}𝑑tte^{t^2/2}K_0(2t^{1/2}/x^{1/4}),$$ (67) $`K_0`$ being the zeroth-order modified Bessel function. The curvature parameter $`\gamma (a)`$ is given by the solution to the following implicit equation $$a^2=\gamma \frac{3F_1+2F_2}{F_1},$$ (68) where the functions $`F_r`$ are evaluated at $`x=\gamma `$. Next, the transport coefficients are expressed in terms of $`a`$ and $`F_r(\gamma )`$ as $$F_\eta (a)=F_0,$$ (69) $$F_\kappa (a)=\frac{F_0}{5}\frac{3F_1+2F_2}{F_1},$$ (70) $$\mathrm{\Psi }_1(a)=2F_1\frac{3F_1+4F_2}{3F_1+2F_2},$$ (71) $$\mathrm{\Psi }_2(a)=4\frac{F_1F_2}{3F_1+2F_2},$$ (72) $$\mathrm{\Phi }(a)=\left[5F_2+2F_3+a^2(F_2+5F_3+8F_4+4F_5)\right].$$ (73) For small shear rates, one gets $`F_\eta 1\frac{18}{5}a^2`$, $`F_\kappa 1\frac{162}{25}a^2`$, $`\mathrm{\Psi }_1\frac{14}{5}(1\frac{1476}{175}a^2)`$, $`\mathrm{\Psi }_2\frac{4}{5}(1\frac{288}{25}a^2)`$, $`\mathrm{\Phi }\frac{14}{5}`$, and $`F_\mu 1\frac{12}{5}a^2`$. ### 3 The ES model In Refs. and the solution of the ES model is worked out keeping the Prandtl number Pr as a free parameter. Here we particularize the results to the correct value $`\text{Pr}=\frac{2}{3}`$. It is convenient to express the transport coefficients in terms of an auxiliary parameter $`\beta `$, defined as the solution of the implicit equation $$a^2=\frac{4\beta }{9}\frac{[2\beta (F_1+2F_2)3]^2[3F_1+2F_22\beta F_1(F_1+2F_2)]}{F_0^2[2\beta (F_1+F_2)3]+F_1[2\beta (F_1+2F_2)3]^2},$$ (74) where now the functions $`F_r`$ are evaluated at $`x=\beta `$. The relationship between the curvature parameter $`\gamma `$ and $`\beta `$ is $$\gamma (a)=\frac{2}{9}\beta [32\beta (F_1+2F_2)].$$ (75) The transport coefficients are $$F_\eta (a)=\frac{9F_0}{[2\beta (F_1+2F_2)3]^2},$$ (76) $$F_\kappa (a)=\frac{a^2}{5\gamma (a)}F_\eta (a),$$ (77) $$\mathrm{\Psi }_1(a)=\frac{12\beta }{a^2}\frac{3F_1+4F_22\beta F_1(F_1+2F_2)}{(32\beta F_1)[32\beta F_1(F_1+2F_2]},$$ (78) $$\mathrm{\Psi }_2(a)=\frac{24\beta }{a^2}\frac{F_2}{(32\beta F_1)[32\beta F_1(F_1+2F_2]},$$ (79) $`\mathrm{\Phi }(a)`$ $`=`$ $`{\displaystyle \frac{1}{40}}C_1\{a^2[C_1^4F_0^3(F_1+2F_2)9C_1^3F_0^2(F_2+2F_3)+54C_1^2F_0(F_2+4(F_3+F_4))`$ (82) $`108C_1(F_2+5F_3+4(2F_4+F_5))]C_3(C_1F_0F_16F_2)+4C_1^2F_0(F_1+2F_2)`$ $`+4C_1[C_2F_0F_16(F_2+2F_3)]24C_2F_2\}.`$ In Eq. (82), $$C_1\frac{3}{32\beta (F_1+2F_2)},$$ (83) $$C_2\frac{3}{32\beta F_1},$$ (84) $$C_33a^2C_1^3F_0^2+4C_2\{2\beta [C_1(F_1+4F_2)+3F_1]3C_1\}.$$ (85) For small shear rates, one has $`F_\eta 1\frac{21}{5}a^2`$, $`F_\kappa 1\frac{197}{25}a^2`$, $`\mathrm{\Psi }_1\frac{14}{5}(1\frac{2126}{175}a^2)`$, $`\mathrm{\Psi }_2\frac{4}{5}(1\frac{413}{25}a^2)`$, $`\mathrm{\Phi }\frac{7}{2}`$, and $`F_\mu 13a^2`$.
warning/0003/cs0003006.html	ar5iv	text	# Materialized View Selection and Maintenance Using Multi-Query OptimizationWork partly supported by a Govt. of India, Department of Science and Technology Grant, and by an IBM University Partnership Program Grant. The work of Prasan Roy was supported by an IBM Research Fellowship. Ramamritham was also supported in part by NSF grant IRI-9619588. ## 1 Introduction Materialization of views can help speed up query and update processing. Views are especially attractive in data warehousing environments because of the query intensive nature of data warehouses. However, when a warehouse is updated, the materialized views must also be updated. Typically, updates are accumulated and then applied to a data warehouse. Loading of updates and view maintenance in warehouses has traditionally been done at night. While the need to provide up-to-date responses to an increasing query load is growing and the amount of data that gets added to data warehouses has been increasing, the time window available for making the warehouse up-to-date has been shrinking. These trends call for efficient techniques for maintaining the materialized views as and when the warehouse is updated. Given multiple views, the view maintenance problem can be seen as computing the expressions corresponding to the “delta” of the views, given the “delta”s of the base relations that are used to define the views. The contributions of this paper lie in the exploitation of the Multi-Query Optimization (MQO) framework along with our recently developed efficient algorithms for MQO, to compute the delta expressions corresponding to the multiple views defined in a data warehouse. It is not difficult to motivate that query optimization techniques are important for choosing an efficient plan for maintaining a view. For example, consider the expression $`(AB)C`$, where $`A`$, $`B`$ and $`C`$ are multisets (i.e., relations with duplicates). Given that the multiset of tuples $`\delta _C^+`$ is inserted into $`C`$, the change to the view is given by $`(AB)\delta _C^+`$. This expression can equivalently be computed as $`(A\delta _C^+)B`$ and by $`(B\delta _C^+)A`$, one of which may be substantially cheaper to compute. Further, in some cases the view may be best maintained by recomputing it, rather than by finding the differentials as above. Vista \[Vis98\] describes how to extend the Volcano query optimizer \[GM91\] to choose the best plan for computing the differential of the result of an expression. Given a set of queries, multiquery optimization \[Sel88\] provides the possibility of reducing costs by computing shared subexpressions once, materializing them temporarily, and reusing them where required in the given set of queries. Although multiquery optimization was earlier viewed as expensive, our recent work \[RSSB00\] has provided efficient algorithms for multiquery optimization, making it practical. In this paper, we provide practical solutions to the problem of optimizing the update of a set of materialized views, by exploiting these algorithms. Specifically, our contributions are as follows. 1. We extend the multiquery optimization algorithms to find the best plan for computing the differential of a set of expressions, by exploiting shared subexpressions. Sharing of subexpressions occurs when multiple views are being maintained, since related views may share subexpressions, and as a result the maintenance expressions may also be shared. Furthermore, sharing can occur even within the plan for maintaining a single view, as we illustrate later in the paper. Our algorithms choose shared expressions to be temporarily materialized during view maintenance, and choose view maintenance plans that utilize these temporarily materialized results. 2. Just as the presence of views allows queries to be evaluated more efficiently, the maintenance of these views can be made more efficient by the presence of additional views/indices \[RSS96\]. That is, given a set of materialized views to be maintained, we need to choose what additional indices and views should be materialized to minimize overall view maintenance costs. The choice of additional views must be done in conjunction with choosing plans for maintaining the views. For instance, a plan that seems quite inefficient could become the best plan if some intermediate result of the plan is chosen to be materialized and maintained. Our contribution here is to show how to extend the multiquery optimization algorithms of \[RSSB00\] to tackle the problems of selecting permanent materialized views, in conjunction with choosing the best plans for updating the views. We show how to cleanly integrate the choice of expressions/indices to be permanently materialized, with the choice of expressions/indices to be temporarily materialized. It is worth pointing out that although our focus in this paper is to speed up view maintenance, our algorithms can also be used to choose extra temporary and permanent views in order to speed up a workload containing queries and updates (that trigger view maintenance). There has been much earlier work on choosing a set of views to be materialized and maintained to optimize given workloads of queries and updates. The major differences between our work and earlier work can be summarized as follows (we outline the differences in detail in Section 2): 1. Given a set of related materialized views, temporarily materializing common subexpressions could have significant benefit. However, earlier work did not consider how to exploit common subexpressions by temporarily materializing them because of their focus on permanent materialization and common subexpressions involving differential relations cannot be permanently materialized. 2. The earlier work does not cover efficient techniques for the implementation of materialized view selection algorithms, in particular, their integration with query optimizers. In the context of materialized view maintenance, this is an important problem since the cost of view maintenance can be reduced by the presence of (additional) indices on relations, and of appropriate extra materialized views. In contrast, we show how to efficiently choose views/indices to be (permanently) materialized by extending the multiquery optimization algorithms of \[RSSB00\]. The rest of the paper is organized as follows. We outline related work in Section 2, and provide the reader with some background in view maintenance in Section 3. We describe the DAG structure used to represent queries in Section 4, and algorithms to find optimal update plans (without materializing additional views) in Section 5. Section 6 presents an optimized greedy algorithm for selecting extra views for materialization. Section 7 outlines results of a performance study, and Section 8 concludes the paper. ## 2 Related Work There has been a large volume of research on incremental view maintenance in the past decade. Amongst the early work on computing the differential results of operations/expressions was Blakeley et al. \[BCL86\]. More recent work in this area includes \[GL95, CGL<sup>+</sup>96, MQM97\]. Gupta and Mumick \[GM95\] provide a survey of view maintenance techniques. Blakeley et al. \[BCL86\] and Ross et al. \[RSS96\] noted that the computation of the expression differentials has the potential for benefiting from multiquery optimization. In the past, multiquery optimization was viewed as too expensive for practical use. As a result they did not go beyond stating that multiquery optimization could be useful for view maintenance. Our recent work in \[RSSB00\] provides efficient heuristic algorithms for multiquery optimization, and demonstrates that multiquery optimization is feasible and effective. There has been much work on selection of views to be materialized. One notable early work in this area was by Roussopolous \[Rou82\]. Ross et al. \[RSS96\] considered the selection of extra materialized views to optimize maintenance of other materialized views/assertions, and mention some heuristics. Labio et al. \[LQA97\] provide further heuristics. The problem of materialized view selection for data cubes has seen much work, such as \[HRU96\], who propose a greedy heuristic for the problem. Gupta \[Gup97\] and Gupta and Mumick \[GM99\] extend some of these ideas to a wider class of queries. However, (a) none of the above papers consider implementation details that are important for efficient selection of views, and (b) none of these consider how to optimize view maintenance expressions. Vista \[Vis98\] describes how to extend the Volcano query optimizer to optimize view maintenance. However, she does not consider the materialization of expressions, whether temporary or permanent. Optimizations that exploit knowledge of foreign key dependencies can be used to detect that certain join results involving differentials will be empty \[QGMW96, Vis98\]. Roy et al. \[RSSB00\] consider how to perform multiquery optimization by selecting subexpressions/indices for temporary materialization. They present important optimizations of a greedy heuristic for materialized view selection that makes the heuristic practical. However, they do not consider updates or view maintenance, which is the focus of this paper. We utilize the optimizations proposed in \[RSSB00\], but the extensions required to to take update costs into account, and optimize view maintenance expressions, are non-trivial. There has been earlier work on multiquery optimization, including \[Sel88, SSN94, SG90\] and more recently \[SV98\], but none of these consider updates. ## 3 Background and Motivation We assume that updates (inserts/deletes) to relations are logged in corresponding delta relations, which are made available to the view refresh mechanism. We assume for each relation $`r`$, there are two relations $`\delta _r^+`$ and $`\delta _r^{}`$ denoting, respectively, the (multiset of) tuples inserted into and deleted from the relation $`r`$. The view refresh mechanism is invoked as part of an update transaction for immediate update, or periodically for deferred updates. In the second case, updates performed by many transactions may be collected together in the delta relations $`\delta _r^+`$ and $`\delta _r^{}`$ for a relation $`r`$. Our techniques work regardless of whether the updates are immediate or deferred. ### 3.1 Computing the Differential of an Operation There is a considerable amount of literature on computing differentials of operations, as outlined in Section 2. For completeness, we briefly review techniques for computing the differential of an operation in the multiset relational algebra. #### 3.1.1 Differentials of Joins Consider a multiset join $`AB`$, and suppose $`A`$ and $`B`$ are updated by inserting the multisets of tuples $`\delta _A^+`$ and $`\delta _B^+`$ respectively. Let $`A^{old}`$ and $`B^{old}`$ refer to the old values of $`A`$ and $`B`$, that is their contents before the update. The multiset of tuples that get added to the view $`V`$ are denoted by $`\delta _V^+`$, and can be computed as: $$\delta _V^+=(\delta _A^+B^{old})(A^{old}\delta _B^+)(\delta _A^+\delta _B^+)$$ View $`V`$ is then updated as follows: $`VV\delta _V^+`$ Similarly, if tuples $`\delta _A^{}`$ and $`\delta _B^{}`$ are deleted from $`A`$ and $`B`$ respectively, the multiset of tuples $$\delta _V^{}=(\delta _A^{}B^{old})(A^{old}\delta _B^{})(\delta _A^{}\delta _B^{})$$ get deleted from $`V`$, which is then updated by: $`VV\delta _V^{}`$. Updates can be modeled as deletes followed by inserts. If both inserts and deletes are present on a relation, we get a more complex expression for updating the relation. $$VV(A^{old}\delta _B^+)(\delta _A^+B^{old})(\delta _A^+\delta _B^+)(A^{old}\delta _B^{})(\delta _A^{}B^{old})(\delta _A^{}\delta _B^{})$$ In contrast if only one input, say $`A`$, is updated by only insertion the change in the view is much easier to compute: $$\delta _V^+=\delta _A^+B^{old}$$ and similarly for deletions on $`A`$, $$\delta _V^{}=\delta _A^{}B^{old}$$ To keep expressions simple (and for another reason which we describe later in Section 3.2.3) we assume that updates are propagated one relation at a time, and only one type of update at a time. This simply means we compute the effect of all inserts on $`A`$, then update $`A`$ with $`\delta _A^+`$, then compute the effects of all deletes on $`A`$, update $`A`$ with $`\delta _A^{}`$. We then proceed with inserts to $`B`$, and then with deletes from $`B`$. The net result is the same as if the more complicated expressions are used, but the expressions we need to deal with are much simpler. Note that an operation may have two complex expressions as inputs, and if both use a particular relation, even with this restriction there may be differential results on both its inputs, in which case the more complex expression allowing differentials on both inputs must be used. #### 3.1.2 Differentials of Other Operations If the result of a groupby/aggregate operation, such as $`{}_{A}{}^{}𝒢_{count(B)}^{}(E)`$, has been materialized (and the aggregate function is distributive), the change in the aggregate result can be computed using only the changes ($`\delta _E^+`$ and $`\delta _E^{}`$) to the input $`E`$, and the old result of the aggregation.<sup>1</sup><sup>1</sup>1For some operations like $`average`$, the count of tuples in each group must also be materialized. Even for the $`sum`$ operation, the count of tuples is needed to deal with deletions. The group-by/aggregation operation is executed on the tuples from the delta relations, and the results are merged into the existing materialized view using a merge operation. For more details, see, e.g., \[GM95\], and for extensions to operations such as median, see \[RSSS94\]. To compute the differential result of an aggregate/grouping operation whose result has not been materialized, we would have to recompute the aggregate values for all groups which are affected by the update. This may involve significant extra work.<sup>2</sup><sup>2</sup>2There are techniques, such as \[CGL<sup>+</sup>96, MQM97\], that use differentials of more complex forms, such as changes in the value of an aggregate result, to avoid recomputing aggregate values in some special cases. Our techniques can be extended to deal with such differentials, but for simplicity we do not consider such techniques here. Standard techniques are available for computing the differentials of other operations, such as duplicate elimination (and projection), and outer joins \[GL95, GJM97\]. We omit details but note that we can use these techniques without any changes to our optimization algorithms. ### 3.2 Computing the Differential of an Expression Views are defined by potentially complex expressions, hence we need to find the differential of an entire expressions. #### 3.2.1 Generating a Differential Expression Techniques, such as that of \[GL95\] can be used to generate an expression that computes the differential of a given expression. However, the resultant expression can be very large – exponential in the size of the query. For instance consider the view $`V=ABC`$, with inserts on all three relations. The differential in the result of $`V`$ can be computed as > $`(\delta _A^+BC)(A\delta _B^+C)(AB\delta _c^+)(A\delta _B^+\delta _C^+)`$ > $`(\delta _A^+B\delta _C^+)(\delta _A^+\delta _B^+C)(\delta _A^+\delta _B^+\delta _C^+)`$ The size of this expression is exponential in the number of relations. Optimizing such large expressions can be quite expensive, since query optimization is exponential in the size of the expression. There are many common subexpressions in the above expression, and the above expression could be simplified by factoring, to get: $`(\delta _A^+BC)((A\delta _A^+)\delta _B^+C)((A\delta _A^+)(B\delta _B^+)\delta _C^+)`$ But creating simplified forms of differential expressions is difficult with more complex expressions containing operations other than join. Therefore our algorithms use an alternative technique, which we outline in the next section. #### 3.2.2 Propagating Updates Up An Expression An alternative to generating a differential expression is to propagate differentials up an expression \[Rou82, RSS96\]. Propagation is best understood by visualizing an expression as a tree. The differential of a node in the tree is computed using the differential (and if necessary, the old value) of its inputs, as described earlier. We start at the leaves of the tree, and proceed upwards, computing the differential expressions at each node. For example, consider an expression $`(AB)C)`$, and suppose we wish to propagate inserts to $`A`$. We can do so by first computing the differential of the node $`AB`$ as $`\delta _A^+B`$. We then join this differential with $`C`$, which is the other input of its parent node, to get the differential of the parent. As mentioned earlier, if there are updates to multiple relations, we propagate one type of update to one relation at a time. Doing so simplifies the expressions for computing the differentials, as outlined in Section 3.1, and permits a different evaluation plan to be chosen for each expression; this is essential for efficient view maintenance, as we will see next, in Section 3.2.3. The process of computing the differential of an expression can be expressed purely in terms of how to compute the differentials for each operation in the expression. There is no need to rewrite the entire expression. Note also that the procedure for computing differentials of an expression can be easily extended to handle expressions using new types of operations, so long as we have a way of incrementally computing the differential of the operation. #### 3.2.3 The Role of Query Optimization Consider an expression $`A(BC)`$, and suppose tuples are inserted into $`A`$. We can compute the differential of the result as $`\delta _A^+(BC)`$. If we compute this expression as shown, we would need to compute $`BC`$, which does not involve any $`\delta `$ relation, and hence may be large and expensive. A better way of evaluating the differential may be $`(\delta _A^+B)C`$. Note that the two variants are logically equivalent. Thus, for efficient differential computation, query optimization must be applied to the update expressions to choose the cheapest variant, as proposed in \[Vis98\]. Furthermore, note that if we wish to compute the differential when tuples are inserted into $`C`$, the plan $`(\delta _C^+B)A`$ or $`(\delta _C^+A)B`$ may be preferable to $`(AB)\delta _C^+`$. Thus, using a single expression, such as $`(AB)C`$ to propagate differentials to $`A`$, $`B`$ and $`C`$ is likely to perform badly for at least one of the differentials. For this reason, we propagate differentials of only one relation at a time, and choose a separate plan for each differential propagation. We use a query optimizer for choosing best plans for such propagation, and our optimizer uses a DAG representation that compactly represents all expressions equivalent to a given expression. Since all alternative expressions above are available, the best one can be chosen for the propagation of each differential. We present details in Sections 4 and 5. Note also that recomputation of a materialized view is always an alternative to computing the differential in its result and updating it. Thus, the query optimizer must choose recomputation over incremental view maintenance, if recomputation is cheaper. ### 3.3 The Role of Multi-Query Optimization in View Update Multi-query optimization attempts at exploiting common sub-expressions within a query, or across queries in a batch of queries submitted together, to reduce the query evaluation cost. In the context of view update, sharing can occur across the tasks of computing differentials of different views, or even within the task of computing the differential of a single view, as we show below. It is easy to see that related queries may share subexpressions, and if so, it may be best to compute the shared subexpression once, materialize it, and reuse it. However, this decision must be done in a cost based manner, as the following example from \[RSSB00\] illustrates. ###### Example 3.1 Let $`Q_1`$ and $`Q_2`$ be two queries whose locally optimal plans (i.e., individual best plans) are $`(RS)P`$ and $`(RT)S`$ respectively. The best plans for $`Q_1`$ and $`Q_2`$ do not have any common sub-expressions. However, if we choose the alternative plan $`(RS)T`$ (which may not be locally optimal) for $`Q_2`$, then, it is clear that $`RS`$ is a common sub-expression and can be computed once and used in both queries. This alternative with sharing of $`RS`$ may be the globally optimal choice. On the other hand, blindly using a common sub-expression may not always lead to a globally optimal strategy. For example, there may be cases where the cost of joining the expression $`RS`$ with $`T`$ is very large compared to the cost of the plan $`(RT)S`$; in such cases it may make no sense to reuse $`RS`$ even if it were available. $`\mathrm{}`$ In the context of view maintenance, if two materialized views have common subexpressions, as in the example above, the expressions for computing the differential of the common subexpression would also be shared. To illustrate subexpression sharing possibilities within a single view maintenance query, consider a view $`V`$ defined as in the example below. ###### Example 3.2 Let view $`V=ABCD`$. Suppose there are inserts on all four relations. The differential of $`V`$ can then be computed using > $`(\delta _A^+BCD)((A\delta _A^+)\delta _B^+CD)`$ > $`((A\delta _A^+)(B\delta _B^+)\delta _C^+D)((A\delta _A^+)(B\delta _B^+)(C\delta _C^+)\delta _D^+)`$ The above expression represents algebraically the effect of propagating differentials one at a time, as described in Section 3.2.2. Note that there are several potential common subexpressions in the above expression. For instance, if the plans chosen for the first two terms of the above union are $`(\delta _A^+(B(CD)))`$ and $`((A\delta _A^+)(\delta _B^+(CD)))`$, then $`CD`$ is a common subexpression of the two. If the above plans are chosen, the subexpression can be computed once and shared. Similarly, $`(A\delta _A^+)(B\delta _B^+)`$ is potentially a common subexpression. The alternative plans of $`(((\delta _A^+B)C)D)`$ and $`((((A\delta _A^+)\delta _B^+)C)D)`$ offer no sharing possibilities, but may still be cheaper. $`\mathrm{}`$ Which the above plans should be used, and whether the common subexpressions should materialized and shared is a decision for the multiquery optimizer to make in a cost-based manner. Thus, it is the job of the multiquery optimizer to find the best overall plan taking sharing possibilities into account. ## 4 DAG Representation of Queries Our algorithms use an extended form of the DAG representation of queries used, for instance, in Volcano \[GM91\]. In this section we summarize the DAG representation and terminology from \[RSSB00\]. An AND–OR DAG is a directed acyclic graph whose nodes can be divided into AND-nodes and OR-nodes; the AND-nodes have only OR-nodes as children and OR-nodes have only AND-nodes as children. An AND-node in the AND-OR DAG corresponds to an algebraic operation, such as the join operation ($``$) or a select operation ($`\sigma `$). It represents the expression defined by the operation and its inputs. Hereafter, we refer to the AND-nodes as operation nodes. An OR-node in the AND-OR DAG represents a set of logical expressions that generate the same result set; the set of such expressions is defined by the children AND nodes of the OR node, and their inputs. We shall refer to the OR-nodes as equivalence nodes henceforth. ### 4.1 Representing a Single Query A given query is initially represented directly in the AND-OR DAG formulation. For example, the query tree of Figure 1(a) is initially represented in the AND-OR DAG formulation, as shown in Figure 1(b). Equivalence nodes (OR-nodes) are shown as boxes, while operation nodes (AND-nodes) are shown as circles. The initial AND-OR DAG is then expanded by applying all possible transformations on every node of the initial query DAG representing the given set of queries. Suppose the only transformations possible are join associativity and commutativity. Then the plans $`A(BC)`$ and $`(AC)B`$, as well as several plans equivalent to these modulo commutativity can be obtained by transformations on the initial AND-OR-DAG of Figure 1(b). These are represented in the DAG shown in Figure 1(c). We shall refer to the DAG after all transformations have been applied as the expanded DAG. Note that the expanded DAG has exactly one equivalence node for every subset of $`\{A,B,C\}`$; the node represents all ways of computing the joins of the relations in that subset. ### 4.2 Representing Sets of Queries in a DAG Queries are inserted into the DAG structure one at a time. When a query is inserted, equivalence nodes and operation nodes are created for each of the operations in its initial query tree. Some of the subexpressions of the initial query tree may be equivalent to expressions already in the DAG. Further, subexpressions of a query may be equivalent to each other, even if syntactically different. For example, query may contain a subexpression that is logically equivalent to, but syntactically different from another subexpression of the query (e.g., $`(AB)C`$, and $`A(BC)`$). Before the second subexpression is expanded, the DAG would contain two different equivalence nodes representing the two subexpressions. \[RSSB00\] modifies the Volcano DAG generation algorithm such that whenever it finds nodes to be equivalent (in the above example, after applying join associativity) it unifies the nodes, replacing them by a single equivalence node. The Volcano optimizer \[GM91\] already has a hashing-based scheme to efficiently detect repeated expressions, thereby avoiding creation of new nodes that are equivalent to existing nodes. The extension of \[RSSB00\] additionally unifies existing logically equivalent nodes. Another extension is to detect and handle subsumption derivations. For example, $`\sigma _{A<5}(E)`$ can be computed from $`\sigma _{A<10}(E)`$ if they both appear in a set of queries. Similarly, if we have aggregations $`{}_{dno}{}^{}𝒢_{sum(Sal)}^{}(E)`$ and $`{}_{age}{}^{}𝒢_{sum(Sal)}^{}(E)`$, we can introduce a new equivalence node $`{}_{dno,age}{}^{}𝒢_{sum(Sal)}^{}(E)`$ and add derivations of the other two from this one. For more details of unification and subsumption derivations involving selections and aggregation, see \[RSSB00\]. ### 4.3 Physical Properties It is straightforward to refine the above AND-OR DAG representation to represent physical properties \[GM91\], such as sort order, that do not form part of the logical data model, and obtain a physical AND-OR DAG <sup>3</sup><sup>3</sup>3For example, an equivalence node is refined to multiple physical equivalence nodes, one per required physical property, in the physical AND-OR DAG.. The presence of an index on a result is also modeled as a physical property of the result by \[RSSB00\], making the code that handles physical properties also perform index selection. Physical properties of intermediate results are important; for example, if an intermediate result is sorted on a join attribute, the join cost can potentially be reduced by using a merge join. This also holds true of intermediate results that are materialized and shared. Our implementation indeed handles physical properties, including sort orders and indices, but to keep the description simple we do not explicitly consider physical properties. ## 5 Finding Optimal Plans We first outline how to find optimal plans for queries, following \[RSSB00\], and then outline extensions to find optimal plans for view maintenance. In both cases, we assume that the set of views chosen for materialization is fixed. In Section 6 we outline how to integrate the choice of views to materialize with the choice of optimal plans for view maintenance. ### 5.1 Finding Optimal Plans for Queries The Volcano optimization algorithm finds the best plan for each node of the expanded DAG by performing a depth first traversal of the DAG. Costs are defined for operation and equivalence nodes. The computation cost of an operation node is $`o`$ is defined as follows: $`compcost(o)=`$ cost of executing $`(o)`$ \+ $`\mathrm{\Sigma }_{e_ichildren(o)}compcost(e_i)`$ The children of $`o`$ (if any) are equivalence nodes. The computation cost of an equivalence node $`e`$ is given as $`compcost(e)=min\{compcost(o_i)\|o_ichildren(e)\}`$ and is $`0`$ if the node has no children (i.e., it represents a relation).<sup>4</sup><sup>4</sup>4Relation scans are explicitly represented as an operation and assigned a cost. Note that the cost of executing an operation $`o`$ also takes into account the cost of reading the inputs, if they are not pipelined. Volcano also caches the best plan it finds for each equivalence node, in case the node is re-visited during the depth first search of the DAG. A simple extension of the Volcano algorithm to find best plans given a set of materialized views is described in \[RSSB00\]. We outline this extension below. Let $`reusecost(m)`$ denote the cost of reusing the materialized result of $`m`$, and let $`M`$ denote the set of materialized nodes. To find the cost of a node given a set of nodes $`M`$ have been materialized, we simply use the Volcano cost formulae above for the query, with the following change. When computing the cost of an operation node $`o`$, if an input equivalence node $`e`$ is materialized (i.e., in $`M`$), the minimum of $`reusecost(e)`$ and $`compcost(e)`$ is used for computing $`compcost(o)`$. Thus, we use the following expression instead: $`compcost(o,M)=`$ cost of executing $`(o)`$ \+ $`\mathrm{\Sigma }_{e_ichildren(o)}C(e_i,M)`$ \| where $`C(e_i,M)=compcost(e_i)`$ if $`e_iM`$ \| \| --- \| \| \| $`=min(compcost(e_i,M),reusecost(e_i))`$ if $`e_iM`$. \| and $`compcost`$ for equivalence nodes is defined as before. Thus, the extended optimizer computes best plans for the query in the presence of materialized results. The extra optimization overhead is quite small. ### 5.2 Extending the DAG Structure for Computing Differentials We now outline how to extend the DAG structure to represent the differentials of a set of expressions. We first construct the expanded DAG for the given expression (or set of expressions). As in Volcano, each equivalence node in the DAG has a set of logical properties such as the schema of the expression result, and estimated statistics about the result such as number of tuples. Once the best plan is computed for a node, it is cached in case it is needed later during optimization. If there are $`n`$ relations, $`R_1,\mathrm{},R_n`$, we need to store information about the differentials of the node with respect to $`\delta _{R_1}^+`$, $`\delta _{R_1}^{}`$, $`\delta _{R_2}^+`$, $`\delta _{R_2}^{}`$, and so on until $`\delta _{R_n}^+`$, $`\delta _{R_n}^{}`$. We number these updates as $`1,\mathrm{},2n`$, and use these numbers to identify the update. To optimize differential plans, each equivalence node stores information for the differentials of the expression with respect to each update type, in addition to information about the full result. Each equivalence node $`e`$ therefore stores an array of $`2n`$ records, as below. Each odd numbered entry $`2i1,i=1..n`$, of this array contains: 1. logical properties (such as schema and estimated statistics) of the differential of $`e`$ with respect to inserts on $`R_i`$ 2. the best plan for computing the differential of $`e`$ with respect to inserts on $`R_i`$ 3. the logical properties of the full result of the equivalence node after inserts and deletes to relations $`R_1,\mathrm{},R_{i1}`$ have been propagated Similarly, each even numbered entry $`2i,i=1..n`$, of this array contains similar information on differentials and best plans with respect to the deletes on $`R_i`$, and the logical properties for the full result of the equivalence node after inserts and deletes to relations $`R_1,\mathrm{},R_{i1}`$, and inserts to $`R_i`$ have been propagated. In addition, as in the original representation, each node stores the best plan for (and cost of) recomputing the entire result of the node after all updates have been made on the base relations. The logical properties of the differentials are computed by a bottom-up traversal of the DAG. We describe later how the best plans for computing differentials are computed and stored. If an equivalence node does not depend on relation $`R_i`$, we flag this during the above bottom-up traversal, and set the plans in entry $`2i`$ and $`2i+1`$ to be null. The traversal also computes and stores an estimate of the execution cost of the differential version of each operation in the DAG (such as a join or an aggregation). The properties of the differentials of its inputs, as well as the full version of the input, where required, are used to compute this estimate. ### 5.3 Finding Optimal Plans for Updates We now outline how to find optimal plans for updates, using the above mentioned DAG representation. Recall the example from Section 3.2.3, with the expression being $`A(BC)`$, and an insert on $`A`$, the plan $`(\delta _A^+B)C`$, is likely to be more efficient than $`\delta _A^+(BC)`$, and should be considered. Luckily for us, the DAG representation of the query represents $`(AB)C`$ in addition to $`A(BC)`$ (see Figure 1). We now extend the technique for finding optimal plans for queries described in Section 5.1, to find the optimal way of propagating the differential $`\delta _A^+`$. Some equivalence nodes do not depend on some relations, and their differential with respect to the relation will be empty. Let $`\text{diffChildren}(o,i)`$ denote all equivalence node children of $`o`$ whose differential is non-empty on update $`i`$, and $`\text{fullChildren}(o,i)`$ denotes all children of $`o`$ whose full results are required to compute the differential of $`o`$, in conjunction with $`\text{diffChildren}(o,i)`$. For instance, diffChildren for an operation that joins $`A`$ with $`(BC)`$, with respect to an insert on $`B`$, is the node $`BC`$, and correspondingly fullChildren of the node is $`A`$. Given an operation node $`o`$ in the DAG, let $`\delta (o,i)`$ denote the differential of operation $`o`$ with respect to update $`i`$. Also let $`\text{localDiffCost}(o,i)`$ denote the cost of executing the operations in $`\delta (o,i)`$, without counting the cost of generating its inputs. Similarly, for an equivalence node $`e`$, let $`\delta (e,i)`$ denote the differential result of $`e`$ with respect to update $`i`$. Then, the total cost of generating the differential result of an operation node $`o`$ with respect to update $`i`$, $`\text{diffCost}(o,i)`$ can be computed by: $$\text{localDiffCost}(o,i)+\mathrm{\Sigma }_{e_j\text{diffChildren}(o,i)}\text{diffCost}(e_j,i)+\mathrm{\Sigma }_{e_j\text{fullChildren}(o,i)}compcost(e_j)$$ The cost of computing the differential of an equivalence node $`e`$ with respect to update $`i`$ is given as $`\text{diffCost}(e,i)=min\{\text{diffCost}(o_j,i)\|o_jchildren(e)\}`$ and is $`0`$ if the node has no children (i.e., it represents a relation or a relation differential). The definition of $`compcost`$ is as defined earlier in Section 5.1, and represents the cost of recomputation of the node after updates have been performed on the database relations. The above formula is extended for the case where some nodes are materialized, as follows. Note that the full result of a node may be materialized, and independently, any of its differential results may also be (temporarily) materialized. Let the set of materialized results be $`M`$; For an operation node $`o`$, we compute $`\text{diffCost}(o,M,i)`$ as: $$\text{localDiffCost}(o,i)+\mathrm{\Sigma }_{e_j\text{diffChildren}(o,i)}C(e_j,M,i)+\mathrm{\Sigma }_{e_j\text{fullChildren}(o,i)}C(e_j,M)$$ where $`C`$ is defined as follows: if $`\delta (e,i)`$ is not materialized (i.e., not in $`M`$), $`C(e,M,i)=min\{\text{diffCost}(o_j,M,i)\|o_jchildren(e)\}`$ and if $`\delta (e,i)`$ is materialized (i.e., in $`M`$), $`C(e,M,i)=min(reusecost(e,i),min\{\text{diffCost}(o_j,M,i)\|o_jchildren(e)\})`$ and $`reusecost(e,i)`$ denotes the cost of reusing the materialized result of $`\delta (e,i)`$. Also, $`C(e_j,M)`$ plays the same role for the full result of node $`e_j`$, as defined in Section 5.1. For an equivalence node $`e`$, $`\text{diffCost}(e,M,i)=min\{\text{diffCost}(o_j,M,i)\|o_jchildren(e)\}`$ and is $`0`$ if the node has no children (i.e., it represents a relation or a relation differential). That is, diffCost represents the cost of computing the differential, even if the differential is materialized. For each equivalence node $`e`$, the operation node corresponding to the minimum cost in the above formula defines the (top node of the) best plan for $`\delta (e,i)`$, given that results in $`M`$ are materialized. Further, we can compute the total cost of computing the differential of a node as $`\text{totalDiffCost}(e,M)=\mathrm{\Sigma }_{i=1\mathrm{}2n}\text{diffCost}(e,M)`$. We perform a single traversal of the DAG to compute the costs for each equivalence/operation node, based on the above equations. During the traversal we also cache the best (minimum cost) plan computed for each differential, just as we cache the best plans for each full result. Note that if both inputs to a join $`E_1E_2`$ are expressions using a common relation $`R`$, an update to $`R`$ results in changes to both inputs, and as a result, the update expression for the join is $`(\delta _{E_1}^+E_2)((E_1\delta _{E_1}^+)\delta _{E_2}^+)`$. In this case, a join in the original expression has been converted into a union of two joins. The best plan for each join is found, giving the best plan for the entire expression, and this combined best plan must be stored (and used to compute the cost of finding the differential). Optimizations that exploit knowledge of foreign key dependencies can be used to detect that certain join results involving differentials will be empty \[QGMW96, Vis98\]. For instance, if $`r.B`$ is a foreign key into $`s.A`$, then the join of $`\delta _s^+`$ and $`r`$ will be empty. Based on this, parts of the differential expression can be detected to be empty, and eliminated during optimization. ## 6 The Greedy Algorithm for Selecting Materialized Views Till now we assumed that the set of materialized nodes is fixed. We now describe how to integrate the choice of extra materialized views/indices with the choice of best plans for view maintenance. Our algorithm is based on a greedy heuristic. We first present the basic algorithm, then some optimizations, and extensions, below. ### 6.1 The Basic Greedy Algorithm As outlined earlier, we first take the given set of materialized views $`𝒱`$, and build a DAG structure on the expressions defining the views. The nodes of the DAG corresponding to views in $`𝒱`$ are marked as already chosen for materialization. We consider both full and differential results for materialization. A result is identified by a node and an update number (in our implementation a full result is identified by the update number $`0`$, and differential results by numbers $`1\mathrm{}2n`$). If a result is chosen for temporary materialization, we must take into account the cost of computing it. And if it is chosen for permanent materialization, we must take into account the cost of maintaining it (we need not consider the cost of initial materialization since it is a one time cost). The cost of maintaining a node incrementally is the sum of the costs of its differentials: $`maintcost(n,M)=\text{totalDiffCost}(n,M)+mergeCost(n)`$ where $`mergeCost(n)`$ denotes the cost of updating the materialized result of $`n`$ using the differentials. For a full result $`n`$, we define $`cost(n,M)=min(compcost(n,M)+matcost(n),maintcost(n,M))`$ where $`matcost(n)`$ denotes the cost of writing out the computed result of $`n`$. That is, when finding the cost of the full result of a materialized node, we take the minimum of the cost via recomputation and the cost via computing the differentials. For a differential result $`x=\delta (n,i)`$, we define $`cost(x,M)=\text{diffCost}(n,M,i)+matcost(x)`$. Given a set $`S`$ of results (full/differential), let $`cost(S,M)`$ be defined as $`cost(S,M)=\mathrm{\Sigma }_{qS}cost(q,M)`$ Given a set of results $`M`$ already chosen to be materialized, and a result $`x`$, the benefit of additionally materializing $`x`$, $`benefit(x,M)`$, is defined as: $`benefit(x,M)=cost(M,M)(cost(M,\{x\}M)+cost(x,M))`$ Note that $`(cost(M,\{x\}M)+cost(x,M))`$ is equivalent to $`cost(M\{x\},M\{x\})`$. Figure 2 outlines a greedy algorithm that iteratively picks nodes to be materialized. The procedure takes as input the set of candidate results (equivalence nodes, and their differentials) for materialization. At each iteration, the node $`x`$ that gives the maximum reduction in the cost if it is materialized, is chosen to be added to $`X`$. The procedure not only selects results for maintenance, but also decides on how they should be maintained. Specifically, for full results, it chooses the cheaper of recomputation (including the cost of storing the result), and differential computation (including the cost of performing the computed differential updates). If recomputation is cheaper for a result, and the result was not part of the given set of materialized view, the result can be materialized temporarily during view maintenance, and discarded later. Differential results that are chosen to be materialized are materialized only temporarily since they are only used during view maintenance. ### 6.2 Optimizations The greedy algorithm as described above can be expensive due to the large number of times the function $`benefit`$ is called, (which in turn calls the expensive function $`cost()`$). Some important optimizations to the greedy algorithm for multi-query optimization are presented in \[RSSB00\]. We use two of the optimizations, with some extensions for handling differentials: 1. There are many calls to benefit (and thereby to $`cost()`$) at line L1 of Figure 2, with different parameters. A simple option is to process each call to $`cost`$ independent of other calls. However, observe that the set of materialized nodes which is the second argument of $`cost`$ changes minimally in successive calls — successive calls take parameters of the form $`cost(R,\{x\}X)`$, where only $`x`$ varies. That is, instead of considering $`x_1X`$ for materialization, we are now considering storing $`x_2X`$ for materialization. The best plans computed earlier does not change for nodes that are not ancestors of either $`x_1`$ or $`x_2`$. It makes sense for a call to leverage the work done by a previous call by recomputing best plans only for ancestors of $`x_1`$ and $`x_2`$. A novel incremental cost update algorithm is presented in \[RSSB00\]. This algorithm maintains the state of the DAG (which includes previously computed best plans for the equivalence nodes) across calls to $`cost`$, and may even avoid visiting many of the ancestors of $`x_1`$ (which is unmaterialized) and $`x_2`$ (which is materialized). In our context of finding update plans, we have to modify the incremental cost update algorithm slightly. 1. If the full result of a node is materialized, we update not only the cost of computing the full result of each ancestor node, but also the costs for the $`2n`$ differentials of each ancestor node since the full result may be used in any of the $`2n`$ differentials. Propagation up from an ancestor node can be stopped if there is no change in cost to computing the full result or any of the differentials. 2. If the differential of a node with respect to update $`i`$ is materialized, we update only the differentials of its ancestors with respect to update $`i`$. Propagation can stop on ancestors whose differentials with respect to $`i`$ do not change in cost. 2. The monotonicity optimization works as follows. With the greedy algorithm as presented above, in each iteration the benefit of every candidate node that is not yet materialized is recomputed since it may have changed. The monotonicity optimization is based on the assumption that the benefit of a node cannot increase as other nodes are chosen to be materialized – while this is not always true, it is often true in practice. The monotonicity optimization makes the above assumption, and does not recompute the benefit of a node $`x`$ if the new benefit of some node $`y`$ is higher than the previously computed benefit of $`x`$. It is clearly preferable to materialize $`y`$ at this stage, rather than $`x`$ — assuming monotonicity holds, the benefit of $`x`$ could not have increased since it was last computed, and it cannot be the node with highest benefit now, hence its benefit need not be recomputed now. Thus, recomputations of benefit are greatly reduced. \[RSSB00\] presents a third optimization based on potential sharability of nodes between queries. This optimization is not relevant here, since even a node that is not sharable may be worth materializing permanently. For the purpose of this paper, the details of the above optimizations are not critical, but the interested reader may refer to \[RSSB00\] for details. The Greedy procedure can be extended in a straightforward manner to consider a workload of queries, along with periodic updates, and to choose the best set of results (and indices) to materialize, to minimize the cost of the queries and view update. Some optimizations are needed to handle large workloads \[RSSB00\]. We can also introduce limits on space for storing permanently materialized results and temporarily materialized results. Results can then be materialized in the order of benefit per unit space, instead of just benefit. ## 7 Performance Study We implemented the algorithm described earlier for finding optimal plans for view maintenance. Like the existing multiquery optimization code, the new code implements index selection along with selection of results to materialize. Our current implementation has a restriction in that it only considers full results for materialization, although a version which also considers differential results for materialization should be ready shortly. Thus our estimated benefits are actually conservative, and we may be able to get even better results once the full implementation is ready. However, the benefits are already very significant. ### 7.1 Performance Model We used a benchmark consisting of TPC-D queries (and some variants based on the same TPC-D schema). The performance measure is estimated execution cost, called plan cost in the performance graphs. Our cost model extends the cost model used in the multiquery optimizer, by taking differential computation into account. The cost model used takes into account number of seeks, amount of data read, amount of data written, and CPU time for in-memory processing. Since we do not currently have a query execution engine which we can extend to perform differential view maintenance, we are unable get actual numbers. However, the cost model is fairly sophisticated, and further, benefits for multiquery optimization predicted by the basic cost model have been verified by running rewritten queries on commercial database systems (\[RSSB00\], and results in a companion paper on query result caching), giving support to the accuracy of estimated benefits. We provide performance numbers for different percentages of updates to the database relations; we assume that all relations are updated by the same percentage. To model a growing database, we have twice as many inserts as deletes. In our notation, a 10 percent update to a relation consists of inserting 10% as many tuples are currently in the relation, and deleting 5% of the current tuples. We compare the performance of our greedy algorithm (referred to as Greedy in our discussion and figures) with plain Volcano query optimization extended to choose between recomputation and incremental maintenance of views (referred to as NoGreedy). (The algorithm of \[Vis98\] falls in the same class as NoGreedy, although the optimization method is somewhat different.) For each query, we present results at different update percentages, ranging from 1% to 80%. The cost of view maintenance is affected by the presence of indices. Normally, databases have indices on the primary key attributes of each relation, to check for uniqueness. Hence we assume that for each of the TPC-D relations, an index is present on the primary key attributes. However, we also ran our benchmark assuming that no indices are initially present, and found that all required indices got chosen for permanent materialization by our algorithm. Thus, the cost of the plans we generate were not significantly affected by the presence of indices, although the cost of plans without our optimizations rose if indices were not already present. We assume a TPC-D database at scale factor of 0.1, that is the relations occupy a total of 100 MB. The buffer size is set at 8000 blocks, each of size 4KB, although we also ran some tests at a much smaller buffer size of 1000 blocks. The tests were run on an Ultrasparc 10, with 256 MB of memory. ### 7.2 Performance Results Maintaining Individual Views. Figure 3 shows our results on two queries, the first consisting of the join of 4 relations, without aggregation, and the second consisting of aggregation on the same join. As can be seen from the figures significant benefits are to be had, especially at low update percentages, but there are benefits even at relatively high update percentages. Maintaining a Set of Views. Figure 4 shows our results on two sets of queries, the first containing five queries without aggregation and the second containing five queries with aggregation. The benefit ratio due to Greedy is again excellent at lower update percentages. There is a jump in cost at one point, which is because of the use of an algorithm that depends on an input fitting in memory, and when the input does not fit in memory its cost increases sharply. Figure 5 shows the results on a set of 10 queries, with indices already present on all primary key attributes, and without any indices present initially; all required indices got chosen for materialization. Cost of Optimization. For a set of 10 materialized views, each a join of 3 to 4 TPC-D relations, (whose results are shown in Figure 5), the time for Greedy optimization was 31 seconds. Note however that 31 seconds is low compared to the savings of up to 1000 seconds obtained for one run of view maintenance, and besides it is a one-time cost whereas view maintenance is typically at least a daily task in a data warehouse. Thus, the extra cost for our algorithms is worthwhile. The number of candidates for materialization grows exponentially with the number of relations in a query. We are currently working on techniques to prune the set of candidates, in order to keep optimization time in tight control even with a higher number of relations. Temporary vs. Permanent Materialization. Out of a total of 1600 results that were materialized (totalling across many different queries and query sets that we considered, and across update percentages ranging from 1 percent to 90 percent), we found that for about 1000 the recomputation cost was less, meaning they are materialized temporarily, and for 600 the maintenance cost was less, meaning they were materialized permanently. At 1 to 5 % update rates, the ratio was 281 to 306, while at 50 to 90 % update rates, the ratio changed to 360 to 88, in favor of recomputation. Effect of Buffer Size. With a buffer size of 1000 blocks (instead of 8000 blocks), we found that the costs of plans with and without Greedy optimization went up, but the increase was more for recomputation plans and the benefit ratio for small update percentages was actually more strongly in favor of our algorithms. ## 8 Conclusions and Future Work The problem of finding the best way to maintain a given set of materialized views is an important practical problem, especially in data warehouses/data marts, where the maintenance windows are shrinking. We have presented solutions that exploit commonality between different tasks in view maintenance, to minimize the cost of maintenance. Our techniques are easy to implement on an existing multiquery optimizer. As shown by the results in section 7, our techniques can generate significant speedup in view maintenance cost, and the increase in cost of optimization is acceptable. We therefore believe that our results are a timely solution to an important practical problem. Future work includes further heuristics to decrease optimization cost, and implementing extensions to efficiently handle workloads containing queries. We also plan to port the system to a dynamic query result caching environment; in a companion paper, we study the issue of selecting results to cache dynamically, in the absence of updates.
warning/0003/hep-ph0003235.html	ar5iv	text	# An Analysis of the Asymptotic Limit of Gluon Shadowing. ## Abstract We examine the gluon distributions in nuclei in the asymptotic region defined by $`Q^2\mathrm{}`$, $`x0`$. An analysis using the Double Asymptotic Scaling variables of Ball and Forte is proposed. New scaling relations are predicted which can help disentangling the different mechanisms of low $`x`$ perturbative QCD evolution in nuclei. Submitted to Physics Letters B 1. Nuclear shadowing or the depletion at low Bjorken $`x`$ of the nuclear Deep Inelastic (DI) structure function, $`F_2^A`$, with respect to the nucleon one, $`F_2^N`$, has been observed in a number of experiments ( and references therein). Assuming the universality of parton distributions in nuclei, one expects nuclear shadowing to be present in other high energy processes as well, such as Drell-Yan pair, $`J/\psi `$ and $`\mathrm{{\rm Y}}`$ production in lepton-nucleus, hadron-nucleus and nucleus-nucleus collisions ( and references therein). In particular nuclear shadowing might be concurring with the other mechanisms, among which is the quark-gluon plasma formation, that result in a depletion of the observed cross sections for these processes. Moreover, as the very low $`x`$ regime has become accessible at HERA, new phenomenological studies of high density QCD at the saturation scale predicted in are now possible . A quantitative understanding of both the $`x`$ and $`Q^2`$ dependences of the nuclear parton distributions at low $`x`$ therefore constitutes a practical and necessary step both for interpreting the outcome of future experiments at RHIC and at the LHC and for investigating the onset of parton saturation. Recent calculations rely on non-perturbative models for the nuclear parton distributions at a given (low) scale, $`Q_o^2`$, combined with DGLAP perturbative evolution. They are all therefore affected by the uncertainty in the initial parton distributions and, in particular, in the gluon distribution which governs evolution at low $`x`$, and which is poorly known experimentally. A strong effect is seen by changing the value of $`Q_o^2`$ itself which can plausibly vary within the range, $`Q_o^2=0.8few\mathrm{GeV}^2`$, leading to sensibly different values for the shadowing of both the structure function and the gluon distribution at large $`Q^2`$. Now, as $`Q^2`$ increases the (low $`x`$) gluon distribution in a proton should tend to a universal asymptotic value, corresponding to the Double Leading Logarithmic Approximation (DLLA) result of Ref.. Numerically, this value is attained at $`Q^210^3\mathrm{GeV}^2`$ provided the initial distributions grow much slower than $`x^{1/2}`$. Because of the coupling between the singlet quark and gluon distributions’ evolution, a similar behavior is predicted for the structure function, $`F_2`$ . It is therefore natural to address the question of whether the differences in the initial non perturbative nuclear shadowing will decrease with growing $`Q^2`$ and give rise to a universal asymptotic curve which is within the range of planned experiments. Our first result is a negative one: we will show that because of the form that DGLAP evolution takes in a nucleus the influence of initial conditions is carried on to the largest attainable $`Q^2`$ values. We then examine carefully the asymptotic behavior of the shadowing ratios $`R_G=G_A/G_N`$ and $`R_F=F_2^A/F_2^N`$, in the DLLA (notations are: $`G_{N(A)}`$ and $`F_2^{N(A)}`$ for the gluon distributions and the structure function in a nucleon, $`(N)`$, and in an isoscalar nucleus, $`(A)`$, respectively). Our goal is to ascertain whether it is possible to distinguish among the different approximations the perturbation series takes in a nucleus and at very low $`x`$. The proton structure function data analyzed recently at HERA have been shown to lie in the asymptotic region ($`Q^2Q_o^2=1\mathrm{GeV}^2`$ and $`xx_o=0.1`$, and to evolve according to DLLA . The key test is to prove that the data obey Double Asymtptotic Scaling (DAS) in the variables $`\rho =\gamma ((YY_o)/\xi )^{1/2}`$ and $`\sigma =\gamma ^1((YY_o)\xi )^{1/2}`$, $`\gamma =6/(332N_f)^{1/2}`$, $`Y=\mathrm{ln}1/x`$, $`\xi =\gamma ^2\mathrm{ln}(\mathrm{ln}Q^2/\mathrm{\Lambda }_{QCD}^2/\mathrm{ln}Q_o^2/\mathrm{\Lambda }_{QCD}^2`$). Violations from DAS (other than due to the fact that the data lie in a pre-asymptotic region ) would signal either the onset of contributions beyond standard pQCD evolution , including the beginning of parton saturation . The approach to asymptotia in a nucleus can be first analyzed by assuming that evolution proceeds through DLLA equations as well. We have found that in this case the ratios $`R_F`$ and $`R_G`$ display exact scaling in $`\sigma `$, thus becoming a function of $`\rho `$ only. As in the proton case, this is a model independent result in that we obtain a scaling form, independent from the initial conditions. We then use this result as a basis for addressing the next question i.e. the detection of violations of DAS scaling in nuclei, which could possibly originate at different values of $`xx_o^A`$ and $`Q^2Q_{o,A}^2`$, than in the proton. In particular Unitarity Shadowing Corrections (USC) are expected to affect evolution at $`x_o^A>x_o`$ because of the increase of the tranverse gluon density in a nucleus due to the overlapping of nucleons in the longitudinal direction. On a more speculative basis one might also expect the transition to the $`\mathrm{ln}(1/x)`$ resummation to appear in a different regime or, in the most “exotic” scenario, that medium modifications of the anomalous dimensions could be observed. In our approach such questions can be addressed systematically as they introduce specific scaling violations from the DLLA result, appearing as different $`\sigma `$ dependences in the ratios $`R_G`$ and $`R_F`$. Our main observation is therefore that although it is technically predictable that for a proton target DGLAP evolution and DLLA should break down at very low values of $`x`$ and sufficiently large $`Q^2`$ and give way to $`\mathrm{ln}(1/x)`$ summation and to USC, it is still a major task to be able to pinpoint where and if the transition from the different regimes is going to take place in the kinematical regimes currently under exploration. Our goal is to obtain some new insight by using nuclear targets where the asymptotic regime can in principle be reached at larger $`x`$. As a by-product we obtain quantitative predictions in the asymptotic kinematic regime which should be attainable at RHIC and at the LHC. 2. We first summarize results for ordinary DGLAP evolution applied to the nuclear ratios at low $`x`$, assuming that the proton and the nuclear distributions evolve similarly. As it is well known evolution is driven by the gluon distribution which dominates over the sea quarks one and one can predict the behavior of the shadowing ratios, $`R_G`$ and $`R_F`$ with $`Q^2`$: $`{\displaystyle \frac{R_G}{\mathrm{ln}Q^2}}`$ $``$ $`{\displaystyle _x^1}P_{GG}({\displaystyle \frac{x}{y}},\alpha _S(Q^2)){\displaystyle \frac{G_N(y,Q^2)}{G_N(x,Q^2)}}\left[R_G(y,Q^2)R_G(x,Q^2)\right]{\displaystyle \frac{dy}{y}}`$ (1) $``$ $`{\displaystyle \frac{G_N/Q^2}{G_N}}\left({\displaystyle \frac{G_A(x,Q^2)/\mathrm{ln}Q^2}{G_N(x,Q^2)/\mathrm{ln}Q^2}}R_G(x,Q^2)\right),`$ (2) $$\frac{R_F}{\mathrm{ln}Q^2}_x^1P_{qG}(\frac{x}{y},\alpha _S(Q^2))\frac{G_N(y,Q^2)}{\mathrm{\Sigma }_N(x,Q^2)}\left[R_G(y,Q^2)R_F(x,Q^2)\right]\frac{dy}{y},$$ (3) where we have disregarded the sea quarks distribution on the r.h.s. of the coupled DGLAP evolution equations; $`P_{qG}`$ and $`P_{GG}`$ are the splitting functions evaluated at NLO; and we used the approximation $`F_2^{N(A)}5/18\mathrm{\Sigma }_{N(A)}`$, with $`\mathrm{\Sigma }=_iq_i(x,Q^2)+\overline{q}_i(x,Q^2)`$. For ease of presentation we will use the following notation: $`G_{N(A)}^{}=G_{N(A)}(x,Q^2)/\mathrm{ln}Q^2`$. Eqs.(1) and (2) show that the $`Q^2`$ dependence of the ratios $`R_G`$ and $`R_F`$ is determined by a subtle balance involving both the parton distributions and the ratios themselves . Based on Eqs.(1) and (2), and defining $`R_G(x,Q_o^2)R_G^o`$ and $`R_F(x,Q_o^2)R_F^o`$ for the initial distributions, one can make the following predictions for the $`Q^2`$ dependence of $`R_G`$ and $`R_F`$: i) $`R_G`$ grows with $`Q^2`$. In fact, if as predicted by non-perturbative shadowing models $`R_G`$ is a growing function of $`x`$, then it also grows with $`Q^2`$, the r.h.s. of Eq.(1) being positive ($`yx`$); ii) if $`R_F^o<R_G^o`$, then $`R_F`$ grows with $`Q^2`$; if $`R_F^o>R_G^o`$, then $`R_F`$ initially decreases with $`Q^2`$ until it reaches the value of $`R_G^o`$ and it subsequently starts increasing along with $`R_G`$. Note that from the behavior of $`R_F`$ and $`R_G`$ obtained from the straighforward application of DGLAP equations, one cannot predict the approach of e.g. $`R_G`$ to a universal limiting curve at large $`Q^2`$. As a matter of fact, although the form of Eq.(1) might seem suggestive of a fixed point behavior , this is a priori not the case, since the quantity $`G_A/G_N`$ depends on $`Q^2`$. The rate of change with $`Q^2`$ is instead governed both by: (a) by the ratio, $`G_N^{}/G_N`$; (b) by the difference $`\mathrm{\Delta }_G=G_A^{}/G_N^{}R_G(x,Q^2)`$ at $`Q^2=Q_o^2`$ . Current parametrizations feature an ultra-soft behavior of the gluon distribution at $`Q_o^21\mathrm{GeV}^2`$, i.e. $`G_N0`$ as $`x0`$, thus causing (because of (a)) a rapid evolution which strongly reduces the shadowing in both $`R_G`$ and $`R_F`$, between $`Q_o^2`$ and $`Q^212\mathrm{GeV}^2`$. If on the contrary one assumes $`Q_o^2`$ ranging from $`2`$ to $`5\mathrm{GeV}^2`$, where harder low $`x`$ initial gluons are expected, then the evolution is slower ($`G_N`$ is larger) and the nuclear ratio is basically unchanged at $`Q^210100\mathrm{GeV}^2`$. A further model dependence follows from the usage of different non-perturbative shadowing mechanisms. We examine two in particular: the Aligned Jet Model (AJM) (see and references therein), and Initial State Recombination (ISR) . Both models explain qualitatively the initial onset of shadowing. Accurate quantitative calculations have been performed using the AJM in . Results are summarized in Fig.1, where we show the $`Q^2`$ dependence of the ratios $`R_F(x,Q^2)`$ and $`R_G(x,Q^2)`$ in $`{}_{}{}^{40}Ca`$ at fixed $`x=10^4`$ (Fig. 1(a)) and $`x=10^2`$ (Fig. 1(b)), for both the AJM (full lines) and the ISR model (short dashes). In order to show the dependence on the initial scale $`Q_0^2`$, results are presented for both models by taking, $`Q_0^2=0.8`$ GeV<sup>2</sup>, and $`Q_0^2=5`$ GeV<sup>2</sup> (the latter can be easily distinguished in the graph by observing the shift in the starting point of the curves). Moreover, as the main purpose of the figure is to illustrate the main features of both the nuclear models and the initial parton distributions that will lead to our description on the asymptotic behavior, we have not sought for $`R_F`$ the best agreement with the data. Details on this part are going to be given elsewhere . A remarkable feature is the sensitivity of the ISR model to this initial scale $`Q_0^2`$. It is only for $`Q_0^21`$ GeV<sup>2</sup> that a sizeable gluon shadowing is obtained due to the fact that the amount of initial shadowing is proportional to the square of the initial gluon distributions and to $`\alpha _s(Q_0^2)/Q_0^2`$ . Results obtaind with the AJM vary less dramatically with the initial scale, $`Q_0^2`$, which in this case enters just the $`(q\overline{q})nucleon`$ (or $`(gg)nucleon`$) cross sections . Nonetheless, the initial difference is carried on to $`Q^2`$ as large as 10<sup>3</sup> GeV<sup>2</sup>. Moreover, in both cases we can observe the $`Q^2`$ behaviour outlined before: a rapid suppression of shadowing in the range of $`Q^2`$ up to 2 GeV<sup>2</sup> and a subsequent softer evolution. The comparison between the $`Q^2`$ behaviour for different fixed $`x`$ values (Fig. 1(a) and Fig. 1(b)) shows that evolution is slower for smaller $`x`$. This can be technically understood by noting that, in the DLLA limit the logarithmic derivative of $`R_G`$ in Eq. (2) vanishes. This is why initial discrepancies between models are more likely to persist at smaller $`x`$. In summary, the asymptotic behavior of $`R_G`$ depends on the initial conditions up to the largest values of $`Q^2`$ attainable at low $`x`$. This feature is in common with the proton gluon distributions themselves as shown e.g. in . As it is well known, neither the data nor theoretical arguments can help us defining the optimal values of $`R_G(x,Q_o^2)`$ and $`R_F(x,Q_o^2)`$. The situation that we have described calls for some redefinition of the approach to asymptotia in deep inelastic scattering from nuclei. In the next Section we illustrate how different behaviors of the data could be revealed by extending the double asymptotic scaling analysis of Ref. to nuclei. 3. We now examine the nuclear DI structure function and gluon distribution in the asymptotic regime defined by $`x0`$, $`Q^2\mathrm{}`$. Our goal is to explore scaling relations in nuclei in order to be able to compare theory with data in a model independent way. The derivation of the equations of the DLLA in a nucleus parallels the one for the proton, namely one first writes the DGLAP evolution equation for the gluon distribution in the limit $`n1`$, $`n`$ being the variable in moments space (we have omitted the subscripts $`N(A)`$ unless necessary): $$\frac{g(n,Q^2)}{\mathrm{ln}Q^2}=\frac{\alpha _s(Q^2)}{2\pi }\gamma _{GG}^0(n)g(n,Q^2),$$ (4) $`\gamma _{gg}^0(n)2C_A/(n1)+\kappa `$ being the anomalous dimension in the limit $`n1`$ ($`\kappa =11/6n_f/3C_A`$ is the next-order or subleading contribution in this limit). Solutions in $`(x,Q^2)`$ are found by evaluating the anti-Mellin transform, $$G(x,Q^2)=\frac{1}{2\pi i}_C𝑑ng(n,Q^2)\mathrm{exp}\left[Y(n1)\right],$$ (5) with the saddle point method . In Eq.(5), $`g(n,Q^2)=g(n,Q_0^2)\mathrm{exp}\left[\xi /(n1)+\xi \kappa /2\right]`$; $`g(n,Q_0^2)=_0^1𝑑xx^{(n1)}g(x,Q_0^2)`$, $`g(x,Q_o^2)`$ being the initial gluon distribution, and $`G(x,Q^2)=xg(x,Q^2)`$; $`Y=\mathrm{ln}(1/x)`$, and $`\xi =\gamma ^2\mathrm{ln}(\mathrm{ln}(Q^2/\mathrm{\Lambda }^2)/(\mathrm{ln}Q_o^2/\mathrm{\Lambda }^2))`$, $`\gamma ^2=C_A/(\pi b)`$, $`b=(332N_f)/12\pi `$. We rewrite the integrand in Eq.(5) as: $$\stackrel{~}{g}(n,Q^2)=g(n,Q_o^2)\mathrm{exp}\xi \kappa /2\mathrm{exp}[Yf_1(n)+\xi f_2(n)],$$ (6) where $`f_1(n)=n1`$, $`f_2(n)=(n1)^1`$ and $`Y`$ and $`\xi `$ are both similarly large. If one takes “soft” initial conditions such as, $`G(x,Q_0^2)A_Nx^\lambda `$, $`\lambda 0`$, then $`g(n,Q_0^2)A_N/(n(\lambda +1))`$ has a pole to the left of the saddle point which is found by setting $`[Yf_1(n)+\xi f_2(n)]/n=0`$. The integral in Eq.(5) is then approximated by: $`G(x,Q^2)`$ $`=`$ $`\sqrt{2\pi }\left({\displaystyle \frac{\stackrel{~}{g}(n_o,Q_o^2)}{\stackrel{~}{g}^{\prime \prime }(n_o,Q_o^2)}}\right)^{1/2}\stackrel{~}{g}(n_o,Q^2)=`$ (7) $`=`$ $`\sqrt{2\pi }g(n_o,Q_o^2)\mathrm{exp}(\xi \kappa /2){\displaystyle \frac{\mathrm{exp}[Yf_1(n_o)+\xi f_2(n_o)]}{\left\|Yf_1^{\prime \prime }(n_o)+\xi f_2^{\prime \prime }(n_o)\right\|^{1/2}}}`$ (8) where $`n_0=1+\sqrt{\xi /\mathrm{\Delta }Y}`$, $`\mathrm{\Delta }Y=YY_o`$, $`Y_o=\mathrm{ln}(1/x_o)`$, $`x_o0.1`$, is the saddle point and $`f_{1(2)}^{\prime \prime }(n_o)=^2f_{1(2)}(n)/n^2\|_{n=n_o}`$. By introducing the variables, $`\rho =\gamma ((YY_o)/\xi )^{1/2}`$ and $`\sigma =\gamma ^1((YY_o)\xi )^{1/2}`$, one has: $$n_on_o(\rho )=1+\gamma /\rho ,$$ (9) and, $$GG^{DAS}(\rho ,\sigma )=\sqrt{\pi }f_G(\rho /\gamma )\left(\frac{\gamma }{\rho }\right)\frac{\mathrm{exp}\left[2\gamma \sigma \gamma ^2\kappa /2\right]}{\sqrt{\sigma \gamma }},$$ (10) where $`f_G(\rho /\gamma )`$ is a smooth function describing the initial conditions. DAS is the prediction that, in the hypothesis of soft initial conditions, and in the asymptotic limit defined by $`\sigma \mathrm{}`$ and $`\rho O(1)`$, $`\mathrm{ln}(G^{DAS}/f_G(\rho ,\sigma ))`$ becomes a linear function of $`\sigma `$, independent of the value of $`\rho `$, with slope fixed by the known constant, $`\gamma `$. A similar behavior is found for $`F_2^p`$ and can therefore be compared to the available data . The structure function’s asymptotic behavior is in fact obtained by solving the equation <sup></sup><sup></sup>*We refer here to the singlet part of the structure function: $$\frac{F_2^p(x,Q^2)}{\mathrm{ln}Q^2}=\frac{5}{18}\frac{\alpha _S}{\pi }G(x,Q^2),$$ (11) yielding: $$F_2^p(\rho ,\sigma )=f_\mathrm{\Sigma }(\rho ,\sigma )\mathrm{exp}(2\gamma \sigma ),$$ (12) where $`f_\mathrm{\Sigma }`$ was derived by using the LO expression for $`\alpha _s`$ and, as for the gluons, it depends on the initial conditions. The scaling of $`F_2^p`$ can be seen from Fig.2 (top-right). From Fig.2 one can also see that the data from NMC (triangles) , corresponding to larger $`x`$ with respect to the 1995 HERA ones (open squares) , as well as some of the more recent HERA data with very low $`x`$ and $`Q^2`$ (open dots) , violate scaling (for a better reading compare with top-left). These scaling violations have been interpreted as due to the fact that the kinematics is not yet asymptotic, as it can be easily seen from the fact that the data lie well below $`\sigma 1`$. Deviations from DAS can be predicted also in the case of hard initial conditions i.e. when $`g(n_0,Q_0^2)A_N/(n(\lambda +1))`$ with $`\lambda 0.2`$. Intuitively this corresponds to taking the limit $`Y\xi `$ in Eq.(5), thus defining the new saddle point: $`n_o=(1+\lambda )+1/\mathrm{\Delta }Y`$. The corresponding gluon distribution is then $`G(x,Q^2)`$ $`=`$ $`f_G^h\mathrm{exp}\left[\lambda (YY_o)+\xi /\lambda +\xi \kappa /2\right]`$ (13) $`=`$ $`f_G^h\mathrm{exp}\left[\lambda \sigma \rho +(\gamma ^2/\lambda +\kappa /2)\xi \right]`$ (14) with $`f_G^h(\rho ,\sigma )=\sqrt{2\pi \rho \sigma }A_N`$ (see also ). This behavior supports the presence of USC appearing as a non-linear term in the evolution equation , which has the effect of damping the steep rise of the gluon distribution at small $`x`$. An alternative possible explanation of DAS violations in the recent HERA data which accounts also for the peculiar stooping of the logarithmic slope of the proton structure function, $`F_2/\mathrm{ln}Q^2`$, at low $`x`$ and $`Q^2`$, is that indeed USC need to be taken into account. We now study these two different scenarios for the asymptotic behavior in nuclei, where it is well known that some aspects of perturbative evolution such as USC, are expected (simply based on geometrical arguments) to arise at larger values of $`x`$. In Fig.2 we present the world low $`x`$ data on the nuclear ratios, Eq.(2) as a function of the DAS variable, $`\sigma `$ (bottom-right), and we compare both data and their kinematics (bottom-left) with the proton ones (top). We use in both cases the following values of $`Q_o^2`$, $`\mathrm{\Lambda }_{QCD}`$, $`\rho `$ and $`\sigma `$: $`Q_o^2=1\mathrm{GeV}^2`$, $`\mathrm{\Lambda }_{QCD}=185`$ MeV (LO). In the NLO analysis performed in it was found that: $`Q_o^2=1.8\mathrm{GeV}^2`$, $`\mathrm{\Lambda }_{QCD}=200`$ MeV (NLO), $`\rho 1`$ and $`\sigma 1.2`$. From the figure it appears that the existing nuclear data are scarse and do not presently support a DAS type analysis and, moreover, they seem to lie mainly in a pre-asymptotic region. Experiments at RHIC and LHC are however expected to be able to cover the asymptotic region. Although experimental extractions of the logarithmic slopes in nuclei have been performed in , very little can be concluded from these data as well . We evaluate the ratios $`R_G`$ and $`R_F`$, Eq.(2) in DAS, i.e. within the hypotheses: i) the quark and gluon distributions are initially shadowed due to some non-perturbative mechanism; ii) the pQCD evolution mechanism is not affected by the nuclear medium; iii) the initial distributions are soft. The DLLA predictions are: $$R_G(x,Q^2)=\left[\frac{\stackrel{~}{g}_A(n_o,Q^2)}{\stackrel{~}{g}_N(n_o,Q^2)}\right]^{3/2}[\frac{\stackrel{~}{g}_N^{\prime \prime }(n_o,Q^2)}{\stackrel{~}{g}_A^{\prime \prime }(n_o,Q^2)}]^{1/2},$$ (15) where $`n_o`$ is the saddle point defined by Eq.(9) for both the proton and the nucleus We rewrite Eq.(15) as a function of the DAS variables by using Eq.(10): $$R_G^{DAS}=\frac{f_G^A(\rho /\gamma )}{f_G^N(\rho /\gamma )}$$ (16) i.e. the exponential terms appearing in $`G^{DAS}(\rho ,\sigma )`$, Eq.(10), are the same in a nucleus and in a single nucleon respectively, thus canceling the $`\sigma `$ dependence in the ratio $`R_G`$: $`R_G`$ is predicted to scale exactly in $`\sigma `$ to a smooth function of $`\rho `$ that is determined entirely by the (soft) initial conditions. The onset of a different evolution mechanism in the nucleus will appear as a $`\sigma `$-scaling violation modifying the exponential behavior of Eqs.(10)-(15) with respect to the single nucleon case. Since the low $`x`$ behavior of recent HERA data seem to show evidence for rather large screening corrections, and at the same time they do not rule out the hard pomeron contribution, we consider the effect of: (A) USC combined with soft initial conditions; (B) USC combined with hard initial conditions; (C) Hard initial conditions in both nucleon and nucleus, no USC. The effect of USC is taken into account through a “damping factor”, $`D_G(x,Q^2)1`$, (Ref. and references therein), evaluated using Mueller-Glauber’s eikonal approximation . We extended the calculation to nuclei by assuming that two gluons inside a nucleus are correlated by a larger confinement radius, $`R_Ar_0A^{1/3}`$, than in a nucleon, corresponding to a smeared impact parameter space two-gluon form factor (details of this calculation will be included in ). As a result the effect USC is enhanced in a nucleus with respect to a nucleon target, due to the larger transverse gluon density at similar values of $`x`$. The asymptotic gluon distribution function is written in terms of the damping factor as: $$G_{N(A)}^{SC}(\rho ,\sigma )=D_G^{N(A)}\times G^{DAS}(\rho ,\sigma ).$$ (17) By using Eq.(17) We can now calculate the ratio $`R_G`$, for the cases listed above. We obtain: $`R_G^{(A)}(\rho ,\sigma )`$ $`=`$ $`R_G^{DAS}(\rho )\times {\displaystyle \frac{\mathrm{exp}\left[2\gamma (\sigma \sigma _A)\frac{\gamma ^2\kappa }{2}\frac{\sigma }{\rho }\right]}{\mathrm{exp}\left[2\gamma (\sigma \sigma _N)\frac{\gamma ^2\kappa }{2}\frac{\sigma }{\rho }\right]}}`$ (19) $``$ $`R_G^{DAS}(\rho )\times \mathrm{exp}\left[(\sigma _A\sigma _N)\right].`$ (20) $`R_G^{(B)}(\rho ,\sigma )`$ $`=`$ $`{\displaystyle \frac{f_G^{h,A}}{f_G^{h,N}}}\times {\displaystyle \frac{\mathrm{exp}\left[\lambda _A\sigma \rho +\frac{\gamma ^2}{\lambda _A}\frac{\sigma }{\rho }\sigma _A\right]}{\mathrm{exp}\left[\lambda _N\sigma \rho +\frac{\gamma ^2}{\lambda _N}\frac{\sigma }{\rho }\sigma _N\right]}}`$ (21) $``$ $`C\mathrm{exp}\left[(\lambda _A\lambda _N)\sigma \rho +\gamma ^2({\displaystyle \frac{1}{\lambda _A}}{\displaystyle \frac{1}{\lambda _N}}){\displaystyle \frac{\sigma }{\rho }}\right]\times \mathrm{exp}\left[(\sigma _A\sigma _N)\right].`$ (22) $`R_G^{(C)}(\rho ,\sigma )`$ $`=`$ $`C\mathrm{exp}\left[(\lambda _A\lambda _N)\sigma \rho +\gamma ^2({\displaystyle \frac{1}{\lambda _A}}{\displaystyle \frac{1}{\lambda _N}}){\displaystyle \frac{\sigma }{\rho }}\right].`$ (23) Here $`R_G^{DAS}`$ is the same as in Eq.(16); $`\sigma _{N(A)}(\rho ,\sigma )=1/2\gamma \mathrm{ln}(1/D_G)`$; and the shadowing for the initial hard distributions has been parametrized as $`R_G^{(o)}=f_G^{h,A}/f_G^{h,N}=Cx^\alpha `$, with $`C1.3`$ and $`\alpha 0.080.1`$. Moreover, we have chosen $`Q_o^2=1\mathrm{GeV}^2`$ for both soft and hard initial conditions and $`\lambda _N=0.35`$. Our scaling result, Eq.(16), and the scaling violating ones, Eqs.(15a)-(15c), are shown in Fig.3 for two different values of $`\rho `$: $`\rho =1.8`$ well inside the asymptotic region shown in Fig.2, and $`\rho =3.4`$ corresponding to very low $`x`$ and almost fixed $`Q^2`$ where we expect standard DGLAP to break down. A few comments are in order. Starting from soft initial conditions one obtains either the $`\sigma `$-scaling curves (full lines) or the $`\sigma `$-scaling violation, Eq.(15a), induced by USC (dashed lines). These corrections are driven by the damping factor given for each value of $`\rho `$ by the dashed curves below. The decreasing trend with growing $`\sigma `$ is larger at large $`\rho `$ because of the correspondingly decreasing values of $`x`$ at similar values of $`Q^2`$. On the other side, hard initial conditions, Eq.(14) and dot-dashed curves in Fig.3, show sensible deviations from $`\sigma `$-scaling at large $`\rho `$ where the $`\mathrm{ln}(1/x)`$ term dominates over the $`\mathrm{ln}Q^2`$ one. This effect is enhanced by USC (dotted curves). An interesting observation is the change in slope of the scaling violations both when passing from soft to hard initial conditions at low $`\rho `$, and when passing from low $`\rho `$ to large $`\rho `$. These and similar other regularities could be studied systematically both for the gluon and the structure function ratios once a much larger set of data will be available. Most importantly, with the approach proposed here we eliminate the ambiguities in the determination of the value of gluons and quarks nuclear shadowing illustrated in Fig.1, in that we identify scaling relations that must be verified independently from the initial non-perturbative nuclear shadowing. In conclusion, by applying DAS to nuclei we have shown model independent predictions, i.e. scaling in the variable $`\rho `$ for the ratios of the nuclear gluon distributions to the free nucleon ones in the asymptotic region. Similar relations hold for the nuclear structure functions. We have considered a few possible sources of scaling violations due to the onset of USC and to the domination of a hard pomeron. As a result, by studying the $`A`$-dependence of shadowing we find some new constraints on perturbative evolution at low $`x`$. Our calculations are relevant for the regime accessible at future experiments at RHIC, LHC and at the $`eA`$ project at DESY .
warning/0003/astro-ph0003278.html	ar5iv	text	# A new approach to the evolution of cosmological perturbations on large scales ## I Introduction Structure in the Universe is generally supposed to originate from the quantum fluctuation of the inflaton field. As each scale leaves the horizon during inflation, the fluctuation freezes in, to become a perturbation of the classical field. The resulting cosmological inhomogeneity is commonly characterized by the intrinsic curvature of spatial hypersurfaces defined with respect to the matter. This metric perturbation is a crucial quantity, because at approach of horizon re-entry after inflation it determines the adiabatic perturbations of the various components of the cosmic fluid, which seem to give a good account of large-scale structure . To compare the inflationary prediction for the curvature perturbation with observation, we need to know its evolution outside the horizon, through the end of inflation, until re-entry on each cosmologically relevant scale. The standard assumption is that the curvature perturbation is practically constant. This has recently been called into question in the context of preheating models at the end of inflation where non-inflaton perturbations can be resonantly amplified . The purpose of the present paper is to investigate the circumstances under which the curvature perturbation may vary. Using only the local conservation of energy–momentum, we show that the rate of change of the curvature perturbation on uniform-density hypersurfaces<sup></sup><sup></sup>*The “conserved quantity” $`\zeta `$ was originally defined in Bardeen, Steinhardt and Turner , but constructed from perturbations defined in the uniform Hubble-constant gauge., $`\zeta `$, on large scales is due to the non-adiabatic part of the pressure perturbation. This result is independent of the form of the gravitational field equations, demonstrating for the first time that the curvature perturbation remains constant on large scales for purely adiabatic perturbations in any relativistic theory of gravity where the energy–momentum tensor is covariantly conserved, $`T_{\nu ;\mu }^\mu =0`$. We also show that for adiabatic perturbations produced during single field inflation the curvature perturbation on uniform-density hypersurfaces, $`\zeta `$ , can be identified with the comoving curvature perturbation, $``$ . The pressure perturbation must be adiabatic if there is a definite equation of state for the pressure as a function of density, which is the case during both radiation domination and matter domination. On the other hand, a change in $`\zeta `$ on super-horizon scales will occur during the transition from matter to radiation domination if there is an isocurvature matter density perturbation . We give a simple derivation of this effect in terms of the curvature perturbations on uniform-radiation and uniform-matter hypersurfaces which remain constant throughout. A simple intuitive understanding of how the curvature perturbation on large scales changes, due to the different integrated expansion in locally homogeneous but causally-disconnected regions of the universe, can be obtained within the ‘separate universes’ picture which we describe in section IV. This enables one to model the evolution of the large-scale curvature perturbation using the equations of motion for an unperturbed Robertson–Walker universe. In section V we use this approach to discuss the evolution of the curvature perturbation in single- and multi-field inflation models. ## II Linear scalar perturbations In this section we summarize the essential results from cosmological perturbation theory, applied to the scalar metric perturbations and the associated perturbations in the pressure and energy density. In contrast with the usual approach to cosmological perturbation theory, we shall not invoke any gravitational field equations. We define energy-momentum in the usual way, $$T_{\mu \nu }2\frac{_{\mathrm{mat}}}{g^{\mu \nu }}+g_{\mu \nu }_{\mathrm{mat}},$$ (1) where $`_{\mathrm{mat}}`$ is any contribution to the Lagrange density from matter fields with no external interactions. General coordinate invariance implies the energy-momentum conservation law $`T_{\nu ;\mu }^\mu =0`$, without invoking the Einstein field equations. There are many different ways of characterizing cosmological perturbations, reflecting the arbitrariness in the choice of coordinates (gauge), which in turn determines the slicing of spacetime into spatial hypersurfaces, and its threading into timelike worldlines. The line element allowing arbitrary linear scalar perturbations of a Friedmann–Robertson–Walker (FRW) background can be written $`ds^2`$ $`=`$ $`(1+2A)dt^2+2a^2(t)_iBdx^idt`$ (3) $`+a^2(t)\left[(12\psi )\gamma _{ij}+2_i_jE\right]dx^idx^j.`$ The unperturbed spatial metric for a space of constant curvature $`\kappa `$ is given by $`\gamma _{ij}`$ and covariant derivatives with respect to this metric are denoted by $`_i`$. For comparison with the notation of Bardeen note that $`AA_BQ^{(0)},`$ $`\psi \left(H_L+{\displaystyle \frac{1}{3}}H_T\right)Q^{(0)},`$ (4) $`B{\displaystyle \frac{B_BQ^{(0)}}{ka}},`$ $`E{\displaystyle \frac{H_TQ^{(0)}}{k^2}},`$ (5) where Bardeen explicitly included $`Q^{(0)}(x^i)`$, the eigenmodes of the spatial Laplacian, $`^2`$, with eigenvalue $`k^2`$. The intrinsic curvature of a spatial hypersurface, $`{}_{}{}^{(3)}R`$, is usually described by the dimensionless curvature perturbationThis quantity is denoted $``$ in Refs. . $`\psi `$, where $${}_{}{}^{(3)}R=\frac{6\kappa }{a^2}+\frac{12\kappa }{a^2}\psi +\frac{4}{a^2}^2\psi .$$ (6) The curvature perturbation on fixed-$`t`$ hypersurfaces is a gauge-dependent quantity and under an arbitrary linear coordinate transformation, $`tt+\delta t`$, it transforms as $$\psi \psi +H\delta t.$$ (7) For a scalar quantity $`x`$, such as the energy density or the pressure, the corresponding transformation is $$\delta \rho \delta \rho \dot{\rho }\delta t,$$ (8) where a dot denotes differentiation with respect to coordinate time $`t`$. The curvature perturbation on uniform-density hypersurfaces, can be written as<sup>§</sup><sup>§</sup>§The sign of $`\zeta `$ is chosen here to coincide with Refs. . $$\zeta =H\xi ,$$ (9) where the displacement between the uniform-density ($`\delta \rho =0`$) hypersurface and the uniform-curvature ($`\psi =0`$) hypersurface has the gauge-invariant definition: $$\xi \frac{\psi }{H}+\frac{\delta \rho }{\dot{\rho }}.$$ (10) Alternatively one can work in terms of the density perturbation on uniform-curvature hypersurfaces $$\delta \rho _\psi =\dot{\rho }\xi ,$$ (11) where the subscript $`\psi `$ indicates the uniform-curvature hypersurface. The curvature perturbation on uniform-density hypersurfaces, $`\zeta `$, is often chosen as a convenient gauge-invariant definition of the scalar metric perturbation on large scales. These hypersurfaces become ill-defined if the density is not strictly decreasing, as can occur in a scalar field dominated universe when the kinetic energy of the scalar field vanishes. In this case one can instead work in terms of the density perturbation on uniform-curvature hypersurfaces, $`\delta \rho _\psi `$, which remains finite. The pressure perturbation (in any gauge) can be split into adiabatic and entropic (non-adiabatic) parts, by writing $$\delta p=c_\mathrm{s}^2\delta \rho +\dot{p}\mathrm{\Gamma },$$ (12) where $`c_s^2\dot{p}/\dot{\rho }`$. The non-adiabatic part is $`\delta p_{\mathrm{nad}}\dot{p}\mathrm{\Gamma }`$, and $$\mathrm{\Gamma }\frac{\delta p}{\dot{p}}\frac{\delta \rho }{\dot{\rho }}.$$ (13) The entropy perturbation $`\mathrm{\Gamma }`$, defined in this way, is gauge-invariant, and represents the displacement between hypersurfaces of uniform pressure and uniform density. ## III Evolution of the curvature perturbation ### A Rate of change of the curvature perturbation on large scales Of primary interest to us, and much of modern cosmology, is the evolution of the curvature perturbation, $`\psi `$, on the constant-time hypersurfaces defined in Eq. (3). These constant-time hypersurfaces are orthogonal to the unit time-like vector field $$n^\mu =(1A,^iB).$$ (14) The expansion of the spatial hypersurfaces with respect to the proper time, $`d\tau (1+A)dt`$, of observers with 4-velocity $`n^\mu `$, is given by $$\theta n_{;\mu }^\mu =3H\left(1A\right)3\dot{\psi }+^2\sigma ,$$ (15) where the scalar describing the shear is $$\sigma =\dot{E}B.$$ (16) However it is useful to define the expansion rate with respect to the coordinate time $$\stackrel{~}{\theta }=(1+A)\theta =3H3\dot{\psi }+^2\sigma .$$ (17) We can write this as an equation for the time evolution of $`\psi `$ in terms of the perturbed expansion, $`\delta \stackrel{~}{\theta }\stackrel{~}{\theta }3H`$, and the shear: $$\dot{\psi }=\frac{1}{3}\delta \stackrel{~}{\theta }+\frac{1}{3}^2\sigma .$$ (18) Note that this is independent of the field equations and follows simply from the geometry. Irrespective of the gravitational field equations we can derive important results from the local conservation of the energy–momentum tensor $`T_{\nu ;\mu }^\mu =0`$. The energy conservation equation $`n^\nu T_{\nu ;\mu }^\mu =0`$ for first-order density perturbations gives $$\dot{\delta \rho }=3H(\delta \rho +\delta p)+(\rho +p)\left[3\dot{\psi }^2\left(\sigma +v+B\right)\right],$$ (19) where $`^iv`$ is the perturbed 3-velocity of the fluid. In the uniform-density gauge, where $`\delta \rho =0`$ and $`\psi =\zeta `$, the energy conservation equation (19) immediately gives $$\dot{\zeta }=\frac{H}{\rho +p}\delta p_{\mathrm{nad}}\frac{1}{3}^2\left(\sigma +v+B\right).$$ (20) We emphasize that we have derived this result without invoking any gravitational field equations, although related results have been obtained in particular non-Einstein gravity theories . We see that $`\zeta `$ is constant if (i) there is no non-adiabatic pressure perturbation, and (ii) the divergence of the 3-momentum on zero-shear hypersurfaces, $`^2(v+B+\sigma )`$, is negligible. On sufficiently large scales, gradient terms can be neglected and $$\dot{\zeta }=\frac{H}{\rho +p}\delta p_{\mathrm{nad}},$$ (21) which implies that $`\zeta `$ is constant if the pressure perturbation is adiabatic. It has been argued that the divergence is likely to be negligible on all super-horizon scales, and in the following we shall make that assumption. Although there have been many previous discussions of conserved quantities in perturbed FRW cosmologies (which coincide with $`\zeta `$ on large scales), we believe that this is the first time that the constancy of $`\zeta `$ has been derived without reference to any equations of motion for the gravitational field. It holds for linear perturbations about an FRW metric for any relativistic theory of gravity, as a consequence of local energy conservation $`n^\nu T_{\nu ;\mu }^\mu =0`$. ### B Non-Einstein gravity theories The most intensively studied example of non-Einstein gravity is provided by scalar–tensor theories, which include a scalar field, $`\varphi `$, non-minimally coupled to the spacetime curvature. One approach to studying the evolution of the metric perturbation previously applied is to perform a conformal transformation to the Einstein frame in which the scalar field is minimally coupled to the metric, and hence the usual Einstein gravitational field equations hold, but non-minimally coupled to other matter fields (whose energy–momentum tensor has non-vanishing trace). The conservation of the total energy–momentum tensor, including the scalar field, in the Einstein frame ensures that the curvature perturbation in this frame, $`\stackrel{~}{\zeta }`$, will remain constant on large scales, but only so long as $`\delta \varphi /\dot{\varphi }=\delta \rho /\dot{\rho }`$, i.e., only for perturbations obeying the generalized adiabatic condition $`\mathrm{\Gamma }_{\varphi \rho }=0`$ \[see Eq. (31)\], in addition to the adiabatic condition for the fluid, $`\mathrm{\Gamma }=0`$ in Eq. (13). However, Eq. (20) shows that $`\zeta `$ must always be conserved on uniform density hypersurfaces in the original frame where ordinary matter is minimally coupled, for adiabatic fluid perturbations ($`\mathrm{\Gamma }=0`$) independently of the perturbations in $`\varphi `$. The two alternative definitions of the curvature perturbation are equal, $`\zeta =\stackrel{~}{\zeta }`$, only in the special case when $`\delta \varphi /\dot{\varphi }=\delta \rho /\dot{\rho }`$ and it then follows that the curvature perturbation is constant in both frames because the generalized adiabatic condition holds. Non-Einstein gravity (in our four spacetime dimensions) may also emerge from theories involving a large extra dimension . In particular, our proof of Eq. (21) validates a recent discussion of chaotic inflation in these theories, which relied on that equation. ### C Matter plus radiation In a multi-fluid system we can define uniform-density hypersurfaces for each fluid and a corresponding curvature perturbation on these hypersurfaces, $`\zeta _{(i)}\psi \delta \rho _{(i)}/\dot{\rho }_{(i)}`$. Equation (20) then shows that $`\zeta _{(i)}`$ remains constant for adiabatic perturbations in any fluid whose energy–momentum is locally conserved: $`n^\nu T_{(i)\nu ;\mu }^\mu =0`$. Thus, for example, in a universe containing non-interacting cold dark matter plus radiation, which both have well-defined equations of state ($`p_\mathrm{m}=0`$ and $`p_\gamma =\rho _\gamma /3`$), the curvatures of uniform-matter-density hypersurfaces, $`\zeta _m`$, and of uniform-radiation-density hypersurfaces, $`\zeta _\gamma `$, remain constant on super-horizon scales. The curvature perturbation on the uniform-total-density hypersurfaces is given by $$\zeta =\frac{(4/3)\rho _\gamma \zeta _\gamma +\rho _m\zeta _\mathrm{m}}{(4/3)\rho _\gamma +\rho _m}.$$ (22) At early times in the radiation dominated era ($`\rho _\gamma \rho _\mathrm{m}`$) we have $`\zeta _{\mathrm{init}}\zeta _\gamma `$, while at late times ($`\rho _\mathrm{m}\rho _\gamma `$) we have $`\zeta _{\mathrm{fin}}\zeta _\mathrm{m}`$. $`\zeta `$ remains constant throughout only for adiabatic perturbations where the uniform-matter-density and uniform-radiation-density hypersurfaces coincide, ensuring $`\zeta _\gamma =\zeta _\mathrm{m}`$. The isocurvature (or entropy) perturbation is conventionally denoted by the perturbation in the ratio of the photon and matter number densities $$S=\frac{\delta n_\gamma }{n_\gamma }\frac{\delta n_\mathrm{m}}{n_\mathrm{m}}=3\left(\zeta _\gamma \zeta _\mathrm{m}\right).$$ (23) Hence the entropy perturbation for any two non-interacting fluids always remains constant on large scales independent of the gravitational field equations. Hence we recover the standard result for the final curvature perturbation in terms of the initial curvature and entropy perturbationThis result was derived first by solving a differential equation , and then by integrating Eq. (21) using Eq. (22). We have here demonstrated that even the integration is unnecessary. $$\zeta _{\mathrm{fin}}=\zeta _{\mathrm{ini}}\frac{1}{3}S.$$ (24) ## IV The separate universe approach One can proceed to use the perturbed field equations, to follow the evolution of linear perturbations in the metric and matter fields in whatever gauge one chooses. This allows one to calculate the corresponding perturbations in the density and pressure and the non-adiabatic pressure perturbation if there is one, and see whether it causes a significant change in $`\zeta `$. However, there is a particularly simple alternative approach to studying the evolution of perturbations on large scales, which has been employed in some multi-component inflation models . This considers each super-horizon sized region of the Universe to be evolving like a separate Robertson–Walker universe where density and pressure may take different values, but are locally homogeneous. After patching together the different regions, this can be used to follow the evolution of the curvature perturbation with time. Figure 1 shows the general idea of the separate universe picture, though really every point is viewed as having its own Robertson–Walker region surrounding it. Consider two such locally homogeneous regions $`(a)`$ and $`(b)`$ at fixed spatial coordinates, separated by a coordinate distance $`\lambda `$, on an initial hypersurface (e.g., uniform-density hypersurface) specified by a fixed coordinate time, $`t=t_1`$, in the appropriate gauge (e.g., uniform-density gauge). The initial large-scale curvature perturbation on the scale $`\lambda `$ can then be defined (independently of the background) as $$\delta \psi _1\psi _{a1}\psi _{b1}.$$ (25) On a subsequent hypersurface defined by $`t=t_2`$ the curvature perturbation at $`(a)`$ or $`(b)`$ can be evaluated using Eq. (18) \[but neglecting $`^2\sigma `$\] to give $$\psi _{a2}=\psi _{a1}\delta N_a,$$ (26) where the integrated expansion between the two hypersurfaces along the world-line followed by region $`(a)`$ is given by $`N_a=N+\delta N_a`$, with $`N\mathrm{ln}a`$ the expansion in the unperturbed background and $$\delta N_a=_1^2\frac{1}{3}\delta \stackrel{~}{\theta }_a𝑑t.$$ (27) The curvature perturbation when $`t=t_2`$ on the comoving scale $`\lambda `$ is thus given by $$\delta \psi _2\psi _{a2}\psi _{b2}=\delta \psi _1\left(N_aN_b\right).$$ (28) In order to calculate the change in the curvature perturbation in any gauge on very large scales it is thus sufficient to evaluate the difference in the integrated expansion between the initial and final hypersurface along different world-lines. In particular, using Eq. (28), one can evolve the curvature perturbation, $`\zeta `$, on super-horizon scales, knowing only the evolution of the family of Robertson–Walker universes, which according to the separate Universe assumption describe the evolution of the Universe on super-horizon scales: $$\mathrm{\Delta }\zeta =\mathrm{\Delta }N,$$ (29) where $`\mathrm{\Delta }\zeta =\psi _a+\psi _b`$ on uniform-density hypersurfaces and $`\mathrm{\Delta }N=N_aN_b`$ in Eq. (28). As we shall discuss in the next section, this evolution is in turn specified by the values of the relevant fields during inflation, and as a result one can calculate $`\zeta `$ at horizon re-entry from the vacuum fluctuations of these fields. While it is a non-trivial assumption to suppose that every comoving region well outside the horizon evolves like an unperturbed universe, there has to be some scale $`\lambda _\mathrm{s}`$ for which that assumption is true to useful accuracy. If there were not, the concept of an unperturbed (Robertson–Walker) background would make no sense. We use the phrase ‘background’ to describe the evolution on a much larger scale $`\lambda _0`$, which should be much bigger even than our present horizon size, with respect to which the perturbations in section II were defined. It is important to distinguish this from regions of size $`\lambda _\mathrm{s}`$ large enough to be treated as locally homogeneous, but which when pieced together over a larger scale, $`\lambda `$, represent the long-wavelength perturbations under consideration. Thus we require a hierarchy of scales: $$\lambda _0\lambda \lambda _\mathrm{s}cH^1.$$ (30) Ideally $`\lambda _0`$ would be taken to be infinite. However it may be that the Universe becomes highly inhomogeneous on some very much larger scale, $`\lambda _\mathrm{e}\lambda _0`$, where effects such as stochastic or eternal inflation determine the dynamical evolution. Nevertheless, this will not prevent us from defining an effectively homogeneous background in our observable Universe, which is governed by the local Einstein equations and hence impervious to anything happening on vast scales. Specifically we will assume that it is possible to foliate spacetime on this large scale $`\lambda _0`$ with spatial hypersurfaces. When we use homogeneous equations to describe separate regions on length scales greater than $`\lambda _\mathrm{s}`$, we are implicitly assuming that the evolution on these scales is independent of shorter wavelength perturbations. This is true within linear perturbation theory in which the evolution of each Fourier mode can be considered independently, but any non-linear interaction introduces mode-mode coupling which undermines the separate universes picture. The separate universe model may still be used for the evolution of linear metric perturbation if the perturbations in the total density and pressure remain small, but a suitable model (possibly a thermodynamic description) of the effect of the non-linear evolution of matter fields on smaller scales may be necessary in some cases. An application to the study of preheating at the end of inflation is discussed in Section V C. Adiabatic perturbations in the density and pressure correspond to shifts forwards or backwards in time along the background solution, $`\delta p/\delta \rho =\dot{p}/\dot{\rho }c_s^2`$, and hence $`\mathrm{\Gamma }=0`$ in Eq. (13). For example, in a universe containing only baryonic matter plus radiation, the density of baryons or photons may vary locally, but the perturbations are adiabatic if the ratio of photons to baryons remains unperturbed. Different regions are compelled to undergo the same evolution along a unique trajectory in field space, separated only by a shift in the expansion. The pressure $`p`$ thus remains a unique function of the density $`\rho `$ and the energy conservation equation, $`d\rho /dN=3(\rho +p)`$, determines $`\rho `$ as a function of the integrated expansion, $`N`$. Under these conditions, uniform-density hypersurfaces are separated by a uniform expansion and hence the curvature perturbation, $`\zeta `$, remains constant. For $`\mathrm{\Gamma }0`$ it is no longer possible to define a simple shift to describe both the density and pressure perturbation. The existence of a non-zero pressure perturbation on uniform-density hypersurfaces changes the equation of state in different regions of the Universe and hence leads to perturbations in the expansion along different worldlines between uniform-density hypersurfaces. This is consistent with Eq. (20) which quantifies how the non-adiabatic pressure perturbation determines the variation of $`\zeta `$ on large scales . The entropy perturbation between any two quantities (which are spatially homogeneous in the background) has a naturally gauge-invariant definition \[which follows from the obvious extension of Eq. (13)\] $$\mathrm{\Gamma }_{xy}\frac{\delta x}{\dot{x}}\frac{\delta y}{\dot{y}}.$$ (31) We define a generalized adiabatic condition which requires $`\mathrm{\Gamma }_{xy}=0`$ for any physical scalars $`x`$ and $`y`$. In the separate universes picture this condition ensures that if all field perturbations are adiabatic at any one time (i.e. on any spatial hypersurface), then they must remain so at any subsequent time. Purely adiabatic perturbations can never give rise to entropy perturbations on large scales as all fields share the same time shift, $`\delta t=\delta x/\dot{x}`$, along a single phase-space trajectory. ## V Inflation ### A Single-component inflaton field In Section III we showed that the curvature perturbation $`\zeta `$ on the uniform-density gauge is constant on large scales for adiabatic perturbations. A common application of this is to perturbations produced by a single scalar field during inflation. Even this apparently simple case is somewhat subtle since a scalar field obeys a second-order equation of motion and cannot in general be described by an equation of state $`p(\rho )`$, since the total energy can be split between potential and kinetic energy. However, the existence of an attractor solution for a strongly-damped inflaton field allows one to drop the decaying mode as inflation progresses and ensures a unique relation between the field value and its first derivative. The specific relations between the inflaton field and curvature perturbations depends on the choice of gauge. In practice the inflaton field perturbation spectrum can be calculated on uniform-curvature ($`\psi =0`$) slices, where the field perturbations have the gauge-invariant definition $$\delta \varphi _\psi \delta \varphi +\frac{\dot{\varphi }}{H}\psi .$$ (32) In the slow-roll limit the amplitude of field fluctuations at horizon crossing ($`\lambda =H^1`$) is given by $`H/2\pi `$. Note that this is the amplitude of the asymptotic solution on large scales. This result is independent of the geometry and holds for a massless scalar field in de Sitter spacetime independently of the gravitational field equations. The field fluctuation is then related to the curvature perturbation on comoving hypersurfaces (on which the scalar field is uniform, $`\delta \varphi _c=0`$) using Eq. (7), by $$\psi _c=\frac{H}{\dot{\varphi }}\delta \varphi _\psi .$$ (33) We will now demonstrate that for adiabatic perturbations we can identify the curvature perturbation on comoving hypersurfaces, $``$, with the curvature perturbation on uniform-density hypersurfaces, $`\zeta `$. In an arbitrary gauge the density and pressure perturbations of a scalar field are given by $`\delta \rho `$ $`=`$ $`\dot{\varphi }\dot{\delta \varphi }A\dot{\varphi }^2+V^{}\delta \varphi ,`$ (34) $`\delta p`$ $`=`$ $`\dot{\varphi }\dot{\delta \varphi }A\dot{\varphi }^2V^{}\delta \varphi ,`$ (35) where $`V^{}dV/d\varphi `$. Thus we find $`\delta \rho \delta p=2V^{}\delta \varphi `$. For adiabatic perturbations on uniform-density hypersurfaces both the density and pressure perturbation must vanish and thus so does the field perturbation $`\delta \varphi _\rho =0`$ for $`V^{}0`$. Hence the uniform-density and comoving hypersurfaces coincide, and $``$ and $`\zeta `$ are identical, for adiabatic perturbations. The asymptotic solution/growing mode for the scalar field vacuum fluctuation corresponds to a perturbation about the background attractor solution and hence generates a purely adiabatic perturbation on super-horizon scales. Thus the density perturbation when a mode re-enters the horizon during the radiation or matter dominated eras can be directly related to the growing mode of the inflaton field perturbation when that mode left the horizon during inflation due to the constancy of $`\zeta `$ once the decaying mode becomes negligible after horizon crossing . We have shown that this does not depend on any slow-roll type approximation for the inflaton field, nor does it depend on the form of the gravitational field equations. The result holds for any metric theory of gravity that respects local conservation of energy–momentum. As an example, the large-scale curvature perturbation spectrum produced during a period of “brane inflation” has recently been calculated in the four-dimensional effective theory of gravity induced on the world-volume of a 3-brane in five-dimensional Einstein gravity , even though the full theory of cosmological perturbations has yet to be determined in this model. ### B Multi-component inflaton field During a period of inflation it is important to distinguish between “light” fields, whose effective mass is less than the Hubble parameter, and “heavy” fields whose mass is greater than the Hubble parameter. Long-wavelength (super-Hubble scale) perturbations of heavy fields are under-damped and oscillate with rapidly decaying amplitude ($`\varphi ^2a^3`$) about their vacuum expectation value as the universe expands. Light fields, on the other hand, are over-damped and may decay only slowly towards the minimum of their effective potential. It is the slow-rolling of these light fields that controls the cosmological dynamics during inflation. The inflaton, defined as the direction of the classical evolution, is one of the light fields, while the other light fields (if any) will be taken to be orthogonal to it in field space. In a multi-component inflation model there is a family of inflaton trajectories, and the effect of the orthogonal perturbations is to shift the inflaton from one trajectory to another. If all the fields orthogonal to the inflaton are heavy then there is a unique inflaton trajectory in field space. In this case even a curved path in field space, after canonically normalizing the inflaton trajectory, is indistinguishable from the case of a straight trajectory, and leads to no variation in $`\zeta `$. When there are multiple light fields evolving during inflation, uncorrelated perturbations in more than one field will lead to different regions that are not simply time translations of each other. In order to specify the evolution of each locally homogeneous universe one needs as initial data the value of every cosmologically significant field. In general, therefore, there will be non-adiabatic perturbations, $`\mathrm{\Gamma }_{xy}0`$. If the local integrated expansion, $`N`$, is sensitive to the value of more than one of the light fields then $`\zeta `$ is able to evolve on super-horizon scales, as has been shown by several authors . Note also that the comoving and uniform-density hypersurfaces need no longer coincide in the presence of non-adiabatic pressure perturbations. In practice it is necessary to follow the evolution of the perturbations on super-horizon scales in order to calculate the curvature perturbation at later times. In most models studied so far, the trajectories converge to a unique one before the end of inflation, but that need not be the case in general. The separate universe approach described in section IV gives a rather straightforward procedure for calculating the evolution of the curvature perturbation, $`\psi `$, on large scales based on the change in the integrated expansion, $`N`$, in different locally homogeneous regions of the universe. This approach was developed in Refs. for general relativistic models where scalar fields dominate the energy density and pressure, though it has not been applied to many specific models. In the case of a single-component inflaton, this means that on each comoving scale, $`\lambda `$, the curvature perturbation, $`\zeta `$, on uniform-density (or comoving) hypersurfaces must stop changing when gradient terms can be neglected ($`\lambda >\lambda _s`$). More generally, with a multi-component inflaton, the perturbations generated in the fields during inflation will still determine the curvature perturbation, $`\zeta `$, on large scales, but one needs to follow the time evolution during the entire period a scale remains outside the horizon in order to evaluate $`\zeta `$ at later times. This will certainly require knowledge of the gravitational field equations and may also involve the use of approximations such as the slow-roll approximation to obtain analytic results. ### C Preheating During inflation, every field is supposed to be in the vacuum state well before horizon exit, corresponding to the absence of particles. The vacuum fluctuation cannot play a role in cosmology unless it is converted into a classical perturbation, defined as a quantity which can have a well-defined value on a sufficiently long time-scale . For every light field this conversion occurs at horizon exit ($`\lambda H^1`$). In contrast, heavy fields become classical, if at all, only when their quantum fluctuation is amplified by some other mechanism. There has recently been great interest in models where vacuum fluctuations become classical (i.e., particle production occurs) due to the rapid change in the effective mass (and hence the vacuum state) of one or more fields. This usually (though not always ) occurs at the end of inflation when the inflaton oscillates about its vacuum expectation value which can lead to parametric amplification of the perturbations — a process which has become known as preheating . The rate of amplification tends to be greatest for long-wavelength modes and this has lead to the claim that rapid amplification of non-adiabatic perturbations could change the curvature perturbation, $`\zeta `$, even on very large scales . Within the separate universes picture this is certainly possible if preheating leads to different integrated expansion in different regions of the universe. In particular $`\zeta `$ can evolve if a significant non-adiabatic pressure perturbation is produced on large scales. However it is also apparent in the separate universes picture that no non-adiabatic perturbation can subsequently be introduced on large scales if the original perturbations were purely adiabatic. This is of course also apparent in the field equations where preheating can only amplify pre-existing field fluctuations. Efficient preheating requires strong coupling between the inflaton and preheating fields which typically leads to the preheating field being heavy during inflation (when the inflaton field is large). The strong suppression of super-horizon scale fluctuations in heavy fields during inflation means that in this case no significant change in $`\zeta `$ is produced on super-horizon scales before back-reaction due to particle production on much smaller scales damps the oscillation of the inflaton and brings preheating to an end . Because the first-order effect is so strongly suppressed in such models, the dominant effect actually comes from second-order perturbations in the fields . The expansion on large scales is no longer independent of shorter wavelength field perturbations when we consider higher-order terms in the equations of motion. Nonetheless in many cases it is still possible to use linear perturbation theory for the metric perturbations while including second-order perturbations in the matter fields.Formally one considers the matter field perturbations to be of order $`ϵ`$, but the metric perturbations to be of order $`ϵ^2`$. In Ref. this was done to show that even allowing for second-order field perturbations, there is no significant non-adiabatic pressure perturbation, and hence no change in $`\zeta `$, on large scales in the original model of preheating in chaotic inflation. More recently a modified version of preheating has been proposed (requiring a different model of inflation) where the preheating field is light during inflation, and the coupling to the inflaton only becomes strong at the end of inflation. In such a multi-component inflation model non-adiabatic perturbations are no longer suppressed on super-horizon scales and it is possible for the curvature perturbation $`\zeta `$ to evolve both during inflation and preheating, as described in Section V-B. ## VI Conclusions In this paper, we have identified the general condition under which the super-horizon curvature perturbation on spatial hypersurfaces can vary as being due to differences in the integrated expansion along different worldlines between hypersurfaces. As long as linear perturbation theory is valid, then, when spatial gradients of the perturbations are negligible, such a situation can be described using the separate universes picture, where regions are evolved according to the homogeneous equations of motion. In particular, the curvature perturbation on uniform-density hypersurfaces, $`\zeta `$, can vary only in the presence of a significant non-adiabatic pressure perturbation. The result follows directly from the local conservation of energy–momentum and is independent of the gravitational field equations. Thus $`\zeta `$ is conserved for adiabatic perturbations on sufficiently large scales in any metric theory of gravity, including scalar–tensor theories of gravity or induced four-dimensional gravity in the brane-world scenario. Multi-component inflaton models are an example where non-adiabatic perturbations may cause the curvature perturbation to evolve on super-horizon scales. ## Acknowledgments We thank Marco Bruni and Roy Maartens for useful discussions. DW is supported by the Royal Society.
warning/0003/hep-th0003276.html	ar5iv	text	# Difference equations for correlation functions of Belavin’s ℤ_𝑛-symmetric model with boundary reflection ## 1 Introduction Integrable models with a boundary have been studied in massive quantum theories and half infinite lattice models . The boundary interaction is specified by the boundary $`S`$-matrix for massive quantum theories , and by the reflection matrix $`K`$ for lattice models . The integrability in the presence of reflecting boundary is ensured by the reflection equation (boundary Yang–Baxter equation) , in addition to the Yang–Baxter equation for bulk (i.e., without boundary) theory . It was shown in that the boundary vacuum of boundary integrable theories can be expressed in terms of the vacuum and the creation operators in the bulk theory. In the explicit bosonic formulae of the boundary vacuum of the boundary XXZ model were obtained by using the bosonization of the vertex operators associated with the bulk XXZ model . The quantum Knizhnik-Zamolodchikov equations are satisfied by both correlation functions and form factors for bulk field theories and for bulk lattice models with the affine quantum group symmetry. It is shown in that correlation functions and form factors in semi-infinite XXZ/XYZ spin chains with integrable boundary conditions satisfy the boundary analogue of the quantum Knizhnik-Zamolodchikov equation . In this paper we study Belavin’s $`_n`$-symmetric vertex model with integrable boundary condition, the boundary Belavin model. The $`R`$-matrix of Belavin’s model is expressed in terms of elliptic functions of the spectral parameter $`z`$ so that the $`R`$-matrix has doubly quasi periodicity. Thus we expect that the $`K`$-matrix of the boundary Belavin model also possesses appropriate transformation properties with respect to $`z`$ compatible to those of the $`R`$-matrix. We shall show that under such assumption the $`K`$-matrix of the boundary Belavin model is inevitably non-diagonal for $`n>2`$. Our solution is diagonal for $`n=2`$ but different from the one used in . On the basis of boundary CTM bootstrap we find that the correlation functions for the boundary Belavin model satisfy a set of difference equations, the boundary analogue of the quantum Knizhnik–Zamolodchikov equation. Furthermore, by solving the simplest difference equations, we obtain the boundary spontaneous polarization which turns out to be the square of that for the bulk $`_n`$-symmetric model . The rest of this paper is organized as follows. In section 2 we review Belavin’s $`_n`$-symmetric model, thereby fixing our notations. In section 3 we give two non-diagonal solutions to the reflection equation, one is a constant $`K`$-matrix, and the other is an elliptic $`K`$-matrix. Furthermore, we consider the boundary analogue of the vertex-face correspondence to discuss the connection between our $`K`$-matrix and the boundary weights of the $`A_{n1}^{(1)}`$ model . In section 4 we construct lattice realization of the boundary vacuum states and vertex operators from the boundary CTM bootstrap approach. In section 5 we derive difference equations for $`N`$-point functions of the boundary Belavin model. We solve the simplest difference equations with $`N=1`$ for free boundary condition to obtain the explicit expression of the boundary spontaneous polarization. The result gives the higher rank generalization of that for the boundary eight vertex model . In section 6 we summarize the results obtained in this paper, and give some concluding remarks. ## 2 Belavin’s vertex model and the reflection equation ### 2.1 Elliptic theta functions For a complex number $`\tau `$ in the upper half-plane, let $`\mathrm{\Lambda }_\tau :=+\tau `$ be the lattice generated by $`1`$ and $`\tau `$, and $`E_\tau :=/\mathrm{\Lambda }_\tau `$ the complex torus which can be identified with an elliptic curve. For $`a,b`$, introduce the Jacobi theta function $$\vartheta \left[\begin{array}{c}a\\ b\end{array}\right](z,\tau ):=\underset{m}{}\mathrm{exp}\left\{\pi \sqrt{1}(m+a)\left[(m+a)\tau +2(z+b)\right]\right\}.$$ (2.1) Hereafter a positive integer $`n2`$ is fixed and we will use the following compact symbols $$\sigma _𝜶^{(n)}(z)=\vartheta \left[\begin{array}{c}\alpha _2/n+1/2\\ \alpha _1/n+1/2\end{array}\right](z,\tau ),\theta _n^{(j)}(z)=\vartheta \left[\begin{array}{c}1/2j/n\\ 1/2\end{array}\right](z,n\tau ),$$ (2.2) for $`𝜶=(\alpha _1,\alpha _2)`$ and for $`j_n`$; and $$h(z):=\underset{j=0}{\overset{n1}{}}\theta ^{(j)}(z)/\underset{j=1}{\overset{n1}{}}\theta ^{(j)}(0).$$ The superscript $`(n)`$ and the subscript $`n`$ will be often suppressed when we have no fear of confusion. The elliptic theta functions are expressed in terms of the product series $$\begin{array}{c}\theta ^{(j)}(z)=\sqrt{1}\omega ^{j/2}t^{n(1/2j/n)^2}u^{1+2j/n}(t^{2n};t^{2n})_{\mathrm{}}(t^{2j}u^2;t^{2n})_{\mathrm{}}(t^{2(nj)}u^2;t^{2n})_{\mathrm{}},\\ h(z)=t^{(n1)/4}\frac{(t^{2n};t^{2n})_{\mathrm{}}^3}{(t^2;t^2)_{\mathrm{}}^3}\sigma _\mathrm{𝟎}(z,\tau )=\sqrt{1}t^{n/4}\frac{(t^{2n};t^{2n})_{\mathrm{}}^3}{(t^2;t^2)_{\mathrm{}}^2}u^1(u^2;t^2)_{\mathrm{}}(t^2u^2;t^2)_{\mathrm{}},\end{array}$$ (2.3) where $$(a;q_1,\mathrm{},q_k)_{\mathrm{}}:=\underset{m_1=0}{\overset{\mathrm{}}{}}\mathrm{}\underset{m_k=0}{\overset{\mathrm{}}{}}(1aq_1^{m_1}\mathrm{}q_k^{m_k}).$$ ### 2.2 Belavin’s vertex model Let $`V=^n`$ and $`\{v_i\}_{i_n}`$ be the standard orthonormal basis of $`V`$ with the inner product $`(v_j,v_k)=\delta _{jk}`$. Let $`V_z`$ be a copy of $`V`$ with a spectral parameter $`z`$. The $`_n`$-Baxter model is a vertex model on a two-dimensional square lattice $``$ such that the state variables take on values of $`_n`$-spin. Each oriented line of $``$ carries a spectral parameter varying from line to line. We assign a $`_n`$-valued local state on each edge. Let be a local Boltzmann weight for a single vertex with bond states $`i,j,k,l_n`$. Arrows denotes orientations of lines. We now define the linear map on $`V_{z_1}V_{z_2}`$ called the $`R`$-matrix as follows: $`R^{V_{z_1},V_{z_2}}(v_jv_l)={\displaystyle \underset{i,k_n}{}}(v_iv_k)R(z_1z_2)_{jl}^{ik}.`$ Belavin considered the $`_n`$-symmetric model satisfying $$\begin{array}{cc}\text{(i)}& R(z)_{jl}^{ik}=0,\text{ unless }i+k=j+l\text{, mod }n,\hfill \\ \text{(ii)}& R(z)_{j+pl+p}^{i+pk+p}=R(z)_{jl}^{ik},\text{ for every }i,j,k,l\text{ and }p_n.\hfill \end{array}$$ (2.4) In terms of two linear map in $`V`$ $$gv_i=\omega ^iv_i,hv_i=v_{i1},$$ (2.5) where $`\omega =\mathrm{exp}(2\pi \sqrt{1}/n)`$, the conditions (2.4) can be rephrased as follows: $$\begin{array}{ccc}\hfill R(z)(gg)& =& (gg)R(z),\hfill \\ \hfill R(z)(hh)& =& (hh)R(z).\hfill \end{array}$$ (2.6) Thus the $`R`$-matrix of Belavin’s $`_n`$-symmetric model is of the form $$R(z)=\frac{1}{\kappa (z)}\overline{R}(z),\overline{R}(z)=\underset{𝜶G_n}{}u_𝜶(z)I_𝜶I_𝜶^1.$$ (2.7) Here $`G_n=_n_n`$, and $`I_𝜶=g^{\alpha _1}h^{\alpha _2}`$ for $`𝜶=(\alpha _1,\alpha _2)`$. The normalization factor $`\kappa (z)`$ will be given lator. The coefficient function $`u_𝜶(z)`$ is determined by imposing the $`R`$-matrix satisfies the Yang-Baxter equation $`R_{12}(z_1z_2)R_{13}(z_1z_3)R_{23}(z_2z_3)=R_{23}(z_2z_3)R_{13}(z_1z_3)R_{12}(z_1z_2),`$ (2.8) where $`R_{ij}(z)`$ denotes the matrix on $`V^3`$, which acts as $`R(z)`$ on the $`i`$-th and $`j`$-th components and as identity on the other one. Belavin’s solution to (2.8) is given as follows: $$u_𝜶(z)=u_𝜶^{(n)}(z,w):=\frac{1}{n}\frac{\sigma _𝜶(z+w/n)}{\sigma _𝜶(w/n)},$$ (2.9) where $`w(0\text{mod }\mathrm{\Lambda }_\tau )`$ is a constant. It is obvious that the following initial condition holds: $$\overline{R}(0)=P,P(xy)=yx.$$ (2.10) In order to facilitate the derivation of the similar results for the $`K`$-matrix of the boundary $`_n`$-symmetric model, we give brief sketches of proofs of several well known properties for Belavin’s $`R`$-matrix. ###### Proposition 2.1 The Boltzmann weights or the elements of $`R`$-matrix are given as follows : $$\overline{R}(z)_{jl}^{ik}=\{\begin{array}{cc}\frac{h(z)\theta ^{(ik)}(z+w)}{\theta ^{(jk)}(z)\theta ^{(ij)}(w)}\hfill & \text{if }i+k=j+l\text{, mod }n\text{,}\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$ (2.11) \[Proof\] Because of the $`_n`$-symmetry, $`R_{0lj}^{ijkj}(z)=R_{jl}^{ik}(z)={\displaystyle \underset{𝜶G_n}{}}u_𝜶(z)(I_𝜶)_j^i(I_𝜶^1)_l^k=\delta _{j+l}^{i+k}{\displaystyle \underset{\alpha _1_n}{}}u_{(\alpha _1,ji)}(z)\omega ^{(il)\alpha _1}.`$ Set $`^{ab}(z)=R_{0a+b}^{ab}(z)`$. Then we have $$_n^{ab}(z,v)=\underset{\alpha _1_n}{}u_{(\alpha _1,a)}^{(n)}(z,v)\omega ^{b\alpha _1}.$$ (2.12) The transformation property of $`^{ab}(z)`$ and the initial condition $`^{ab}(0)=\delta ^{b0}`$ imply that $$^{ab}(z)=0,\text{at }z=c\tau \text{ (}cb\text{, mod }n\text{) and }z=(ab)\tau w\text{, mod }\mathrm{\Lambda }_{n\tau }.$$ (2.13) Hence $`^{ab}(z)`$ has the form $`^{ab}(z)=C^{ab}(w)\theta ^{(ab)}(z+w){\displaystyle \underset{cb}{}}\theta ^{(c)}(z).`$ By substituting $`z=b\tau `$ we have $`C^{ab}(w)^1=\theta ^{(a)}(w){\displaystyle \underset{c0}{}}\theta ^{(c)}(0),`$ which concludes that (2.11) holds. $`\mathrm{}`$ As a corollary of Proposition 2.11 we have $$PR(w)=R(w),R(w)P=R(w).$$ (2.14) Now we assume that $`0<t<q<u<1`$, where $`t:=\mathrm{exp}(\pi \sqrt{1}\tau ),q:=\mathrm{exp}(\pi \sqrt{1}w)`$, and $`u:=\mathrm{exp}(\pi \sqrt{1}z)`$. Following Baxter we call such domain of parameters the principal regime. Note that (2.11) is weights of the eight-vertex model when $`n=2`$. ### 2.3 Unitarity and crossing symmetry Belavin’s $`R`$-matrix satisfies the unitarity and crossing symmetry relations . ###### Proposition 2.2 Belavin’s $`R`$-matrix satisfies the following unitarity relation or the first inversion relation: $$\overline{R}_{21}(z)\overline{R}_{12}(z)=\rho _1(z,w)II,$$ (2.15) where $$\rho _1(z,w)=\frac{\sigma (z+w)\sigma (z+w)}{\sigma ^2(w)}.$$ (2.16) \[Proof\] Note that $`\overline{R}_{21}(z)\overline{R}_{12}(z)`$ $`=`$ $`{\displaystyle \underset{𝜶G_n}{}}u_𝜶^{(n)}(z,w)I_𝜶^1I_𝜶{\displaystyle \underset{𝜷G_n}{}}u_𝜷^{(n)}(z,w)I_𝜷I_𝜷^1`$ $`=`$ $`{\displaystyle \underset{𝜶𝜷G_n}{}}u_𝜶^{(n)}(z)u_𝜷^{(n)}(z)I_𝜶^1I_𝜷I_𝜶I_𝜷^1`$ $`=`$ $`{\displaystyle \underset{𝒂G_n}{}}f_𝒂^{(n)}(z,w)I_𝒂I_𝒂^1,`$ where $$f_𝒂^{(n)}(z,w)=\underset{𝜶G_n}{}\omega ^{𝜶,𝒂}u_𝜶^{(n)}(z,w)u_{𝒂+𝜶}^{(n)}(z,w),$$ (2.17) and $`𝜶,𝒂=\alpha _1a_2\alpha _2a_1`$. Proposition 2.16 is thus reduced to $$f_𝒂^{(n)}(z,w)=\rho _1(z,w)\delta _{𝒂\mathrm{𝟎}}.$$ (2.18) Concerning the proof of (2.18), see Theorem 3.3 and Lemma 3.2 in . $`\mathrm{}`$ Next we describe the crossing symmetry for Belavin’s $`_n`$-symmetric model. For that purpose let us recall the $`R`$-matrix on $`KL`$, where $`K=V_{z_1}\mathrm{}V_{z_k}`$ and $`L=V_{z_1^{}}\mathrm{}V_{z_l^{}}`$: $$\begin{array}{cc}\hfill R^{K,V_z^{}}:=& R_{1;k+1}^{V_{z_1},V_z^{}}\mathrm{}R_{k;k+1}^{V_{z_k},V_z^{}},\hfill \\ \hfill R^{K,L}:=& R_{1\mathrm{}k;k+l}^{K,V_{z_l^{}}}\mathrm{}R_{1\mathrm{}k;k+1}^{K,V_{z_1^{}}}.\hfill \end{array}$$ YBE holds for $`R^{K,L}`$ by virtue of YBE for $`R^{V,V}`$ (2.8) $`R_{12}^{K,L}R_{13}^{K,M}R_{23}^{L,M}=R_{23}^{L,M}R_{13}^{K,M}R_{12}^{K,L},`$ (2.19) as a linear map on $`KLM`$. For special $`K_z^k=V_{z_1}\mathrm{}V_{z_k}`$ such that $`z_j=z+(k+1j)w`$ ($`1jk`$), the fusion operator $`\pi `$ associated with $`K_z^k`$ is given as follows : $$\pi :=R_{k1;k}^{V_{z_1},V_{z_2}}R_{k2,k1;k}^{V_{z_1}V_{z_2},V_{z_3}}\mathrm{}R_{1,\mathrm{},k1;k}^{V_{z_1}\mathrm{}V_{z_{k1}},V_{z_k}}.$$ (2.20) From the first equation of (2.14) and the Yang–Baxter equation (2.8) we have $$\pi (K_z^k)=\mathrm{\Lambda }^k(V)=\text{Anti}(K_z^k).$$ (2.21) Let $`V^{}`$ be the dual space of $`V`$ and $`\{v_i^{}\}_{i_n}`$ be the dual basis of $`\{v_i\}_{i_n}`$. Then we have the isomorphism $`C:V_{z+nw/2}^{}\text{Anti}(K_z^{n1})`$ $`Cv_i^{}={\displaystyle \underset{i_1,\mathrm{},i_{n1}}{}}{\displaystyle \frac{ϵ_i^{i_1\mathrm{}i_{n1}}}{\sqrt{(n1)!}}}v_{i_1}\mathrm{}v_{i_{n1}},`$ (2.22) where $`ϵ_i^{i_1\mathrm{}i_{n1}}`$ is the $`n`$-th order completely antisymmetric tensor. The spectral parameter $`z+nw/2`$ associated with the dual space $`V^{}`$ refers to the mean value of $`n1`$ spectral parameters $`z+(n1)w`$, $`\mathrm{}`$, $`z+w`$ of $`V`$<sup>1</sup><sup>1</sup>1Note that the spectral parameter of $`V^{}`$ is shifted by $`nw/2`$ from the one in . . Then the $`R`$-matrices on $`VV^{}`$ and $`V^{}V`$ are defined as follows: $$\begin{array}{ccc}\hfill R^{V_{z_1},V_{z_2+nw/2}^{}}& =& (IC)^1R^{V_{z_1},V_{z_2+(n1)w}\mathrm{}V_{z_2+w}}(IC),\hfill \\ \hfill R^{V_{z_1+nw/2}^{},V_{z_2}}& =& (CI)^1R^{V_{z_1+(n1)w}\mathrm{}V_{z_1+w},V_{z_2}}(CI).\hfill \end{array}$$ (2.23) The un-normalized $`\overline{R}`$ on $`VV^{}`$ and $`V^{}V`$ are also defined in a similar manner. ###### Proposition 2.3 The $`R`$-matrix on $`VV^{}`$ and $`V^{}V`$ defined in (2.23) meet the crossing symmetry : $$\begin{array}{ccc}\hfill \overline{R}_{21}^{V_{z_2},V_{z_1+nw/2}^{}}& =& (\overline{R}_{12}^{V_{z_1},V_{z_2}})^{t_1}\underset{p=2}{\overset{n1}{}}\frac{h(z_1+z_2+pw)}{h(w)},\hfill \\ \hfill \overline{R}_{12}^{V_{z_1+nw/2}^{},V_{z_2}}& =& (\overline{R}_{21}^{V_{z_2},V_{z_1+nw}})^{t_1}\underset{p=1}{\overset{n2}{}}\frac{h(z_1+z_2pw)}{h(w)},\hfill \end{array}$$ (2.24) where $`t_i`$ denotes the transposition of the $`i`$-th space. \[Proof\] Let $$\overline{R}_{21}^{V_{z_2},V_{z_1+nw/2}^{}}(v_jv_l^{})=\underset{i,k}{}(v_iv_k^{})a_{jl}^{ik}(z_1+z_2)$$ Because of the initial condition (2.10) and the second equation of (2.14), the element $`a_{jl}^{ik}(z)`$ vanishes at $`z=pw`$, where $`p=2,\mathrm{},n1`$. Thus we have an entire function $`b_{jl}^{ik}(z)`$ from $`a_{jl}^{ik}(z)`$ devided by $`h(z2w)\mathrm{}f(z(n1)w)`$. The transformation property of $`b_{jl}^{ik}(z)`$ are the same as $`\overline{R}_{kj}^{li}(z)`$. It follows from the second equation of (2.14) that $`b_{jl}^{ik}(z)=0`$ at $`z=c\tau `$ for $`cjk`$ and at $`z=(ik)\tau w`$, which coincide the zeros of $`\overline{R}_{kj}^{li}(z)`$ (2.13). Thus $`b_{jl}^{ik}(z)`$ equals $`\stackrel{ˇ}{R}_{ij}^{kl}(z)`$ up to a scalar factor, which is determined by substituting $`z=(ki)\tau `$. The second equation of (2.24) can be shown in a similar way. $`\mathrm{}`$ From (2.16) and (2.24), we have the following second inversion relation $$\underset{jl}{}\overline{R}_{12}^{t_1}(z)\overline{R}_{21}^{t_1}(znw)=\rho _2(z,w)I,$$ (2.25) where $$\rho _2(z,w)=\frac{h(z)h(z+nw)}{h^2(w)}.$$ (2.26) Imposing the unitarity and crossing symmetry condition with respect to the normalized $`R`$-matrix: $$R_{21}(z)R_{12}(z)=II,$$ (2.27) $$R_{21}^{V_{z_2},V_{z_1+nw/2}^{}}=(R_{12}^{V_{z_1},V_{z_2}})^{t_1},R_{12}^{V_{z_1+nw/2}^{},V_{z_2}}=(R_{21}^{V_{z_2},V_{z_1+nw}})^{t_1},$$ (2.28) the normalization factor $`\kappa (z)`$ should obey the following functional equations: $$\begin{array}{ccc}\hfill \kappa (z)\kappa (z)& =& \rho _1(z,w),\hfill \\ \hfill \kappa (z)\kappa (znw)& =& \rho _2(z,w).\hfill \end{array}$$ (2.29) Hereafter $`\kappa (z)`$ is often denoted by $`\kappa (u)`$ through the relation $`u=\mathrm{exp}(\pi \sqrt{1}z)`$. In the principal regime using (2.3) the following expression solves (2.29) $$\kappa (u)=u^{(n2)/n}\frac{(u^2;t^2)_{\mathrm{}}(t^2u^2;t^2)_{\mathrm{}}}{(q^2;t^2)_{\mathrm{}}(t^2q^2;t^2)_{\mathrm{}}}\overline{\kappa }(u),$$ (2.30) where $$\overline{\kappa }(u)=\frac{(q^2u^2;t^2,q^{2n})_{\mathrm{}}(q^{2n}u^2;t^2,q^{2n})_{\mathrm{}}(t^2q^2u^2;t^2,q^{2n})_{\mathrm{}}(t^2q^{2n}u^2;t^2,q^{2n})_{\mathrm{}}}{(q^{2+2n}u^2;t^2,q^{2n})_{\mathrm{}}(u^2;t^2,q^{2n})_{\mathrm{}}(t^2q^{2+2n}u^2;t^2,q^{2n})_{\mathrm{}}(t^2u^2;t^2,q^{2n})_{\mathrm{}}}.$$ From $`\kappa (1)=1`$ the initial condition for $`R`$ also holds: $$R(0)=P.$$ (2.31) ## 3 Boundary Belavin model ### 3.1 Reflection equation for the boundary Belavin model In this section we consider the following reflection equation or the boundary Yang–Baxter equation: $$K_2(z_2)R_{21}(z_1+z_2)K_1(z_1)R_{12}(z_1z_2)=R_{21}(z_1z_2)K_1(z_1)R_{12}(z_1+z_2)K_2(z_2).$$ (3.1) The reflection equation (3.1) is valid when $`z_1=z_2`$ because $`R(0)=P`$. Furthermore, the following Lemma holds: ###### Lemma 3.1 The reflection equation $`(\text{3.1})`$ is valid when $`(1)`$ $`z_1=0;`$ $`(2)`$ $`z_1=z_2`$ provided $$\begin{array}{cc}(1)\text{ Boundary initial condition}:\hfill & K(0)=I;\hfill \\ (2)\text{ Boundary unitarity relation}:\hfill & K(z)K(z)=I,\hfill \end{array}$$ (3.2) respectively. \[Proof\] It is evident from the unitarity (2.27) and the initial condition (2.31) for $`R`$-matrix. $`\mathrm{}`$ Here we notice that Belavin’s $`R`$-matrix have the following quasi-periodic properties $$\begin{array}{ccc}\hfill \overline{R}(z+1)& =& (gI)^1\overline{R}(z)(gI)=(Ig)\overline{R}(z)(Ig)^1,\hfill \\ \hfill \overline{R}(z+\tau )& =& (hI)^1\overline{R}(z)(hI)\times \mathrm{exp}\left\{2\pi \sqrt{1}\left(z+\frac{\tau }{2}+\frac{w}{n}\right)\right\}\hfill \\ & =& (Ih)\overline{R}(z)(Ih)^1\times \mathrm{exp}\left\{2\pi \sqrt{1}\left(z+\frac{\tau }{2}+\frac{w}{n}\right)\right\}.\hfill \end{array}$$ (3.3) Thus we have the following Proposition: ###### Proposition 3.2 Let $$K(z)=\frac{1}{\lambda (z)}\overline{K}(z),$$ where $`\lambda (z)`$ is a scalor function. Suppose (3.2) and the following quasi transformation property: $$\begin{array}{ccc}\hfill \overline{K}(z+1)& =& g\overline{K}(z)g,\hfill \\ \hfill \overline{K}(z+\tau )& =& h\overline{K}(z)h\times \mathrm{exp}\left\{2\pi \sqrt{1}\left(z+\frac{\tau }{2}+c\right)\right\},\hfill \end{array}$$ (3.4) where $`c`$ is a constant. Then $`\overline{K}(z)`$ solves (3.1). \[Proof\] Let $`F(z_1,z_2)`$ stand for the difference of the LHS and the RHS of (3.1). Then we have $$\begin{array}{ccc}\hfill F(z_1+1,z_2)& =& (gI)F(z_1,z_2)(gI),\hfill \\ \hfill F(z_1+\tau ,z_2)& =& (hI)F(z_1,z_2)(hI)\times \mathrm{exp}(2\pi \sqrt{1}B),\hfill \end{array}$$ (3.5) where $`B=3z_1+3\tau /2+2w/n+c`$. The second equation of (3.5) implies that the $`(ik,jl)`$-th element of $`F(z_1,z_2)`$ satisfies $$F(z_1+\tau ,z_2)_{jl}^{ik}=F(z_1+\tau ,z_2)_{j1l}^{i+1k}\times (2\pi \sqrt{1}B).$$ (3.6) Thus we find that $`F(p\tau ,z_2)_{jl}^{ik}F(0,z_2)_{jpl}^{i+pk}=0`$ for $`0pn1`$ from Lemma 3.1. Similarly, we have $`F(z_2+p\tau ,z_2)_{jl}^{ik}=F(z_2+p\tau ,z_2)_{jl}^{ik}=0`$ for $`0pn1`$: $$F(p\tau ,z_2)_{jl}^{ik}=F(z_2+p\tau ,z_2)_{jl}^{ik}=F(z_2+p\tau ,z_2)_{jl}^{ik}=0,(0pn1).$$ (3.7) Assume that $`F(z_1+\tau ,z_2)_{jl}^{ik}`$ is not identically zero. From Richey-Tracy’s lemma (see section 3 in or Lemma 2.4 in ) we conclude that $`F(z_1+\tau ,z_2)_{jl}^{ik}`$ has $`3n`$ zeros in $`E_{n\tau }`$ whose sum is equal to $`nc2w3n(n1)\tau (i+j)\tau `$. The contradiction to (3.7) implies the claim of this Proposition. $`\mathrm{}`$ ### 3.2 Solutions of the reflection equation Under the assumption of the quasi periodicity (3.4) compatible to (3.3) we find that the $`K(z)`$ is not a diagonal matrix for $`n>2`$. When $`n=2`$ we can take $`K(z)`$ diagonal because of $`g^1=g`$ and $`h^1=h`$. The most general and non-diagonal solution for $`n=2`$ is given in . Other non-diagonal solutions for $`D_n^{(2)}`$-vertex model are given in . In this paper we consider the following two solutions of (3.1), which can be also found in . #### 3.2.1 Constant $`K`$-matrix ###### Proposition 3.3 Let $$𝒦_0v_j=v_{nj},$$ (3.8) where $`v_n=v_0`$. Then $`𝒦_0`$ solves (3.1). \[Proof\] It is easy to see $`g𝒦_0g=h𝒦_0h=𝒦_0`$. Hence we have $$\begin{array}{cc}& K_2(z_2)R_{21}(z_1+z_2)K_1(z_1)R_{12}(z_1z_2)\hfill \\ =& I𝒦_0\underset{𝜶}{}u_𝜶(z_1+z_2)(I_𝜶^1I_𝜶)(𝒦_0I)\underset{𝜷}{}u_𝜷(z_1z_2)(I_𝜷I_𝜷^1)\hfill \\ =& 𝒦_0𝒦_0\underset{𝜶}{}\omega ^{\alpha _1\alpha _2}u_𝜶(z_1+z_2)I_𝜶I_𝜶\underset{𝜷}{}\omega ^{\beta _1\beta _2}u_𝜷(z_1z_2)I_𝜷I_𝜷\hfill \\ =& 𝒦_0𝒦_0\underset{𝜷}{}\omega ^{\beta _1\beta _2}u_𝜷(z_1z_2)I_𝜷I_𝜷\underset{𝜶}{}\omega ^{\alpha _1\alpha _2}u_𝜶(z_1+z_2)I_𝜶I_𝜶\hfill \\ =& \underset{𝜷}{}\omega ^{\beta _1\beta _2}u_𝜷(z_1z_2)(I_𝜷I_𝜷)(𝒦_0I)\underset{𝜶}{}\omega ^{\alpha _1\alpha _2}u_𝜶(z_1+z_2)(I_𝜶I_𝜶)(I𝒦_0)\hfill \\ =& R_{21}(z_1z_2)K_1(z_1)R_{12}(z_1+z_2)K_2(z_2),\hfill \end{array}$$ that implies this Proposition. $`\mathrm{}`$ #### 3.2.2 Elliptic $`K`$-matrix Let $$m=\{\begin{array}{cc}n\hfill & \text{if }n\text{ is odd,}\hfill \\ n/2\hfill & \text{if }n\text{ is even},\hfill \end{array}$$ and let $$𝒦(z)=\underset{𝜶G_m}{}\omega ^{2\alpha _1\alpha _2}u_{2𝜶}^{(n)}(z,v)I_{2𝜶}=\underset{𝜶G_m}{}u_{2𝜶}^{(n)}(z,v)J_𝜶,$$ (3.9) where $$J_𝜶=h^{\alpha _2}g^{2\alpha _1}h^{\alpha _2}$$ for $`𝜶=(\alpha _1,\alpha _2)`$, and $`v(0\text{mod }\mathrm{\Lambda }_\tau )`$ is a constant. Using the identity $$\frac{1}{m}\underset{\alpha _1=0}{\overset{m1}{}}\omega ^{2\alpha _1(i\alpha _2)}=\{\begin{array}{cc}\delta _{\alpha _2,i}\hfill & \text{if }n\text{ is odd,}\hfill \\ \delta _{\alpha _2,i}+\delta _{\alpha _2,im}\hfill & \text{if }n\text{ is even,}\hfill \end{array}$$ we have $`𝒦(0)=𝒦_0`$. ###### Lemma 3.4 The following quasi transformation property holds: $$\begin{array}{ccc}\hfill 𝒦(z+1)& =& g^1𝒦(z)g;\hfill \\ \hfill 𝒦(z+\tau )& =& h^1𝒦(z)h\times \mathrm{exp}\left\{2\pi \sqrt{1}\left(z+\frac{\tau }{2}+\frac{v}{m}\right)\right\}.\hfill \end{array}$$ (3.10) \[Proof\] This is based on the transformation properties of the elliptic theta function. $`\mathrm{}`$ ###### Lemma 3.5 Let $`\overline{K}(z)=𝒦_0𝒦(z)`$. Then the boundary inversion relation holds: $$\overline{K}(z)\overline{K}(z)=\rho _1(z,v)I.$$ (3.11) \[Proof\] Direct calculation shows $$\begin{array}{ccc}\hfill \overline{K}(z)\overline{K}(z)& =& 𝒦_0\underset{𝜶G_m}{}u_{2𝜶}^{(n)}(z,v)J_𝜶𝒦_0\underset{𝜷G_m}{}u_{2𝜷}^{(n)}(z,v)J_𝜷\hfill \\ & =& \underset{𝜶G_m}{}u_{2𝜶}^{(n)}(z,v)J_𝜶\underset{𝜷G_m}{}u_{2𝜷}^{(n)}(z,v)J_𝜷\hfill \\ & =& \underset{𝜶G_m}{}\underset{𝜷G_m}{}\omega ^{2𝜶,𝜷}u_{\text{2}𝜶}^{(n)}(z,v)u_{2𝜷}^{(n)}(z,v)J_{𝜶𝜷}\hfill \\ & =& \underset{𝒂G_m}{}g_𝒂^{(n)}(z,v)I_𝒂,\hfill \end{array}$$ where $$g_𝒂^{(n)}(z,v)=\underset{𝜶G_m}{}\omega ^{2𝜶,𝒂}u_{2𝜶}^{(n)}(z,v)u_{2(𝒂+𝜶)}^{(n)}(z,v).$$ (3.12) By comparing $`g_𝒂^{(n)}(z,v)`$ with $`f_𝒂^{(n)}(z,w)`$ defined in (2.17), we easily have $`g_𝒂^{(n)}(z,v)=f_𝒂^{(m)}(z,v)`$ and hence (3.11) holds for even $`n`$. Repeating the similar argument in Proposition 2.16 we can also obtain (3.11) for odd $`n`$. $`\mathrm{}`$ ###### Theorem 3.6 Let $`\overline{K}(z)=𝒦_0𝒦(z)`$. Then $`\overline{K}(z)`$ solves the reflection equation (3.1). \[Proof\] From Lemma 3.10 we find that $`\overline{K}(z)`$ satisfies (3.4) with $`c=v/m`$. Since $`\overline{K}(0)=𝒦_0{}_{}{}^{2}=I`$, the $`\overline{K}(z)`$ also satisfies the first equation of (3.2). It follows from Lemma 3.11 that $`\overline{K}(z)`$ satisfies the second one of (3.2). Thus $`\overline{K}(z)`$ is a solution to the reflection equation (3.1) from Proposition 3.2. $`\mathrm{}`$ Remark. Our $`K`$-matrix for $`n=2`$ is different from the one used in so that the readers should be careful to compare our results with those of $`n=2`$. ### 3.3 Matrix elements of $`K`$-matrix In this subsection we calculate the $`(j,k)`$-th element of $`\overline{K}(z)`$: $$\overline{K}(z)v_k=\underset{j_n}{}v_j\overline{K}(z)_k^j.$$ Note that $$\overline{K}(z)_k^j=𝒦(z)_k^{nj}=\underset{\alpha _2_m}{}\delta _{j+k}^{2\alpha _2}\underset{\alpha _1_m}{}u_{(2\alpha _1,j+k)}^{(n)}(z,v)\omega ^{(jk)\alpha _1}$$ When $`n`$ is even, thanks to the sum over $`\alpha _2`$, $`\overline{K}(z)_k^j=0`$ if $`j+k`$ is odd. By comparing (2.12) we obtain $$\overline{K}(z)_k^j=\{\begin{array}{cc}_m^{\frac{j+k}{2},\frac{jk}{2}}(z,v)\hfill & \text{if }j+k\text{ is even},\hfill \\ 0\hfill & \text{if }j+k\text{ is odd},\hfill \end{array}$$ (3.13) for even $`n`$, $$\overline{K}(z)_k^j=\{\begin{array}{cc}_n^{jk,\frac{jk}{2}}(z,v)\hfill & jk\text{ is even},\hfill \\ _n^{jk,\frac{jk+n}{2}}(z,v)\hfill & jk\text{ is odd},\hfill \end{array}$$ (3.14) for odd $`n`$. We are now in a position to determine the normalization factor $`\lambda (z)`$. The boundary inversion relation (3.11) implies $$\lambda (z)\lambda (z)=\rho _1(z,v).$$ (3.15) Furthermore, the boundary crossing symmetry holds for $`n=2`$ : $$K(z)_k^j=\underset{j^{},k^{}}{}R(2zw)_{1jk}^{j^{}\mathrm{\hspace{0.17em}1}k^{}}K(zw)_j^{}^k^{},$$ (3.16) which implies that $$\frac{\lambda (zw)}{\lambda (z)}=\frac{1}{\overline{\kappa }(u^2)}\frac{(q^2u^2;t^2)_{\mathrm{}}(t^2q^2u^2;t^2)_{\mathrm{}}}{(u^2;t^2)_{\mathrm{}}(t^2u^2;t^2)_{\mathrm{}}}.$$ (3.17) Since $`V^{}\mathrm{\Lambda }^{n1}(V)\overline{)}V`$ for $`n>2`$, the LHS of (3.16) for higher $`n`$ should be replaced by the $`(j,k)`$-th element of the dual $`K`$-matrix. We wish to discuss this point again in section 4. Here we assume the following functional relation holds for $`n2`$: $$\frac{\lambda (z\frac{n}{2}w)}{\lambda (z)}=\frac{1}{\overline{\kappa }(u^2)}\frac{(q^nu^2;t^2)_{\mathrm{}}(t^2q^nu^2;t^2)_{\mathrm{}}}{(u^2;t^2)_{\mathrm{}}(t^2u^2;t^2)_{\mathrm{}}}.$$ (3.18) Not (3.18) but (3.15) is important to calculate the spontaneous polarization in section 5, so that we proceed further under the assumption (3.18). By solving (3.15) and (3.17) we obtain $$\lambda (z)=\frac{1}{(r^2;t^2)_{\mathrm{}}(t^2r^2;t^2)_{\mathrm{}}}\frac{(r^2u^2;t^2,q^{2n})_{\mathrm{}}(t^2r^2u^2;t^2,q^{2n})_{\mathrm{}}}{(r^2q^{2n}u^2;t^2,q^{2n})_{\mathrm{}}(t^2r^2q^{2n}u^2;t^2,q^{2n})_{\mathrm{}}}\frac{\varphi (u^2)}{\varphi (u^2)},$$ (3.19) where $`r=\mathrm{exp}(\pi \sqrt{1}v)`$, and $$\begin{array}{ccc}\hfill \varphi (x)& =& \frac{(q^nx;t^2,q^{2n})_{\mathrm{}}(t^2q^nx;t^2,q^{2n})_{\mathrm{}}}{(q^{2n}x;t^2,q^{2n})_{\mathrm{}}(t^2x;t^2,q^{2n})_{\mathrm{}}(r^2q^nx;t^2,q^{2n})_{\mathrm{}}(t^2r^2q^nx;t^2,q^{2n})_{\mathrm{}}}\hfill \\ & \times & \frac{(q^{2n+2}x^2;t^2,q^{4n})_{\mathrm{}}(t^2q^{2n2}x^2;t^2,q^{4n})_{\mathrm{}}}{(q^{2n}x^2;t^2,q^{4n})_{\mathrm{}}(t^2q^{2n}x^2;t^2,q^{4n})_{\mathrm{}}}.\hfill \end{array}$$ ### 3.4 Comments on boundary weights for the boundary $`A_{n1}^{(1)}`$ face model In this subsection we wish to discuss the boundary analogue of the vertex-face correspondence. Concerning the case $`n=2`$, see . Let us consider the bulk $`A_{n1}^{(1)}`$-face model whose local state takes on values of $`P`$, the weight lattice of $`A_{n1}^{(1)}`$ . An ordered pair $`(a,b)P^2`$ is called admissible if $`b=a+\widehat{j}`$, for a certain $`j_n`$, where $$\widehat{j}=v_j\frac{1}{n}\underset{k=0}{\overset{n1}{}}v_k.$$ Let be the local Boltzmann weight for a state configuration $`(a,b,c,d)`$ round a face. Then $`W\left(\begin{array}{cc}a& b\\ d& c\end{array}\|z\right)=0`$ unless all the four pairs $`(a,b),(a,d),(b,c)`$ and $`(d,c)`$ are admissible. Non-zero Boltzmann weights are given as follows: $$W\left(\begin{array}{cc}a& b\\ d& c\end{array}\|z\right)=\frac{1}{w(z,w)}\overline{W}\left(\begin{array}{cc}a& b\\ d& c\end{array}\|z\right),$$ (3.20) where $`w(z,w)`$ is a scalar function and $$\begin{array}{ccc}\hfill \overline{W}\left(\begin{array}{cc}a& a+\widehat{j}\\ a+\widehat{j}& a+2\widehat{j}\end{array}\|z\right)& =& \frac{h(z+w)}{h(w)},\hfill \\ & & \\ \hfill \overline{W}\left(\begin{array}{cc}a& a+\widehat{j}\\ a+\widehat{j}& a+\widehat{j}+\widehat{k}\end{array}\|z\right)& =& \frac{h(a_{jk}wz)}{h(a_{jk}w)}(jk),\hfill \\ & & \\ \hfill \overline{W}\left(\begin{array}{cc}a& a+\widehat{k}\\ a+\widehat{j}& a+\widehat{j}+\widehat{k}\end{array}\|z\right)& =& \frac{h(z)}{h(w)}\frac{h(a_{jk}w+w)}{h(a_{jk}w)}(jk).\hfill \end{array}$$ (3.21) Here $$a_{jk}=\overline{a_j}\overline{a_k},\overline{a_j}=(a+\rho ,v_j),$$ and $`\rho ={\displaystyle \underset{j=0}{\overset{n1}{}}}(n1j)\widehat{j}`$ is the half sum of the positive roots. Jimbo, Miwa and Okado introduced the intertwining vectors to show the equivalence between the $`_n`$-symmetric model and the $`A_{n1}^{(1)}`$ model. Let $$\begin{array}{cc}\hfill t_b^a(z):=& {}_{}{}^{t}(t_b^{a(0)}(z),\mathrm{},t_b^{a(n1)}(z)),\hfill \\ \hfill t_b^{a(i)}(z):=& \{\begin{array}{cc}\theta ^{(i)}(z+\delta nw\overline{a}_j)\hfill & \text{if }b=a+\overline{ϵ}_j,\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}\hfill \end{array}$$ (3.22) where $`\delta `$ is an arbitrary constant. Then we have the so-called vertex-face correspondence : $$\overline{R}(z_1z_2)t_d^a(z_1)t_c^d(z_2)=\underset{b}{}\overline{W}\left(\begin{array}{cc}a& b\\ d& c\end{array}\|z_1z_2\right)t_c^b(z_1)t_b^a(z_2).$$ (3.23) Thanks to (3.23) the Boltzmann weights (3.21) solve the face-type Yang-Baxter equation : $$\begin{array}{cc}& \underset{g}{}\overline{W}\left(\begin{array}{cc}b& c\\ g& d\end{array}\|z_1z_2\right)\overline{W}\left(\begin{array}{cc}a& b\\ f& g\end{array}\|z_1z_3\right)\overline{W}\left(\begin{array}{cc}f& g\\ e& d\end{array}\|z_2z_3\right)\\ & \\ =& \underset{g}{}\overline{W}\left(\begin{array}{cc}a& b\\ g& c\end{array}\|z_2z_3\right)\overline{W}\left(\begin{array}{cc}g& c\\ e& d\end{array}\|z_1z_3\right)\overline{W}\left(\begin{array}{cc}a& g\\ f& e\end{array}\|z_1z_2\right).\end{array}$$ (3.24) Let us now consider the boundary $`A_{n1}^{(1)}`$-face model. By analogy with the bulk case, we find the following Proposition: ###### Proposition 3.7 Assume that the existence of boundary weights $`V`$’s satisfying $$\overline{K}(z)t_c^a(z)=\underset{b}{}\overline{V}\left(a\begin{array}{c}b\\ c\end{array}\|z\right)t_b^a(z).$$ (3.25) Then $`\overline{V}`$ solves the face-type reflection equation $$\begin{array}{cc}& \underset{b,e}{}\overline{V}\left(f\begin{array}{c}g\\ e\end{array}\|z_2\right)\overline{W}\left(\begin{array}{cc}a& f\\ b& e\end{array}\|z_1+z_2\right)\overline{V}\left(b\begin{array}{c}e\\ c\end{array}\|z_1\right)\overline{W}\left(\begin{array}{cc}a& b\\ d& c\end{array}\|z_1z_2\right)\hfill \\ & \\ =& \underset{b,e}{}\overline{W}\left(\begin{array}{cc}a& f\\ b& g\end{array}\|z_1z_2\right)\overline{V}\left(b\begin{array}{c}g\\ e\end{array}\|z_1\right)\overline{W}\left(\begin{array}{cc}a& b\\ d& e\end{array}\|z_1+z_2\right)\overline{V}\left(d\begin{array}{c}e\\ c\end{array}\|z_2\right).\hfill \end{array}$$ (3.26) In order to solve (3.25), let us recall the dual intertwining vectors $$\begin{array}{cc}\hfill t_a^b(z):=& (t_{a(0)}^b(z),\mathrm{},t_{a(n1)}^b(z)),\hfill \\ \hfill t_{a(i)}^{a+\widehat{j}}(z):=& (\stackrel{~}{\mathrm{\Phi }}^a(z))_j^i/det\mathrm{\Phi }^a(z).\hfill \end{array}$$ (3.27) Here $`\mathrm{\Phi }^a(z)`$ is a matrix whose $`(i,j)`$-component is $`t_{a+\widehat{j}}^{a(i)}(z)`$, and $`\stackrel{~}{\mathrm{\Phi }}^a(z)`$ is a cofactor matrix of $`\mathrm{\Phi }^a(z)`$. Note that $`t_b^a(z)`$ is a column vector while $`t_a^b(z)`$ is a row vector. Thus by the rule of multiplication of matrices, $`t_a^b(z)t_d^c(z^{})`$ represents a scalar function while $`t_b^a(z)t_d^c(z^{})`$ does a function-valued matrix. Since $`t_b^a(z)`$ and $`t_a^b(z)`$ enjoy the following orthogonal properties $`t_a^{a+\widehat{j}}(z)t_{a+\widehat{k}}^a(z)`$ $`=`$ $`\delta _{jk},`$ (3.28) $`{\displaystyle \underset{j=0}{\overset{n1}{}}}t_{a+\overline{ϵ}_j}^a(z)t_a^{a+\widehat{j}}(z)`$ $`=`$ $`I_n,`$ (3.29) the boundary analogue of the vertex-face correspondence (3.25) is equivalent to $$\overline{V}\left(a\begin{array}{c}b\\ c\end{array}\|z\right)=t_a^b(z)K(z)t_c^a(z)=\underset{j,k}{}t_{a(j)}^b(z)K(z)_k^jt_c^{a(k)}(z).$$ (3.30) ###### Proposition 3.8 Let $$V\left(a\begin{array}{c}b\\ c\end{array}\|z\right)=\frac{1}{\lambda (z)}\overline{V}\left(a\begin{array}{c}b\\ c\end{array}\|z\right),$$ where $`\lambda (z)`$ is the same scalar function as for $`K(z)`$, and $`\overline{V}`$ is defined by (3.30). Then the boundary weights $`V`$’s satisfy the initial condition $$V\left(a\begin{array}{c}b\\ c\end{array}\|0\right)=\delta _c^b,$$ (3.31) and the inversion relation $$\underset{g}{}V\left(a\begin{array}{c}b\\ g\end{array}\|z\right)V\left(a\begin{array}{c}g\\ c\end{array}\|z\right)=\delta _c^b.$$ (3.32) \[Proof\] The initial condition (3.31) follows from that for $`K(z)`$ and (3.28). The inversion relation (3.32) follows from (3.29), (3.11) and (3.28). $`\mathrm{}`$ The boundary weights $`V\left(a\begin{array}{c}b\\ c\end{array}\|z\right)`$ are non-diagonal in the sense that they do not vanish even for $`bc`$ as a function of $`z`$. Hence (3.30) does not coincide with the diagonal solution of (3.26) involving the bulk Boltzmann weights for the $`A_{n1}^{(1)}`$-face model given in for $`n2`$. Such disagreement indicates that there may exist unknown solution to (3.1) corresponding to the solution given in and also unknown solution to (3.26) corresponding to our $`K`$-matrix, throughout the boundary vertex-face correspondence. ### 3.5 Commuting transfer matrix The transfer matrix with $`L`$ columns, is expressed in terms of $`R`$ and $`K`$-matrices as follows : $$\begin{array}{ccc}\hfill T_L(z_1,z_2)& =& \text{Tr}_0K_+(z_1)𝒯(z_1,z_2)\hfill \\ \hfill 𝒯(z_1,z_2)& =& \text{Tr}_0𝒯(z_1z_2)^1K_{}(z_1)𝒯(z_1z_2).\hfill \end{array}$$ (3.33) Here $$\begin{array}{ccc}\hfill 𝒯(z_1z_2)& =& R_{01}^{V_{z_1},V_{z_2}}\mathrm{}R_{0L}^{V_{z_1},V_{z_2}}\text{End}(V_0V_1\mathrm{}V_L),\hfill \\ \hfill 𝒯(z_1z_2)^1& =& R_{L0}^{V_{z_2},V_{z_1}}\mathrm{}R_{10}^{V_{z_2},V_{z_1}}\text{End}(V_0V_1\mathrm{}V_L),\hfill \end{array}$$ are monodromy matrices satisfying $$R_{12}(z_1z_2)𝒯_1(z_1)𝒯_2(z_2)=𝒯_2(z_2)𝒯_1(z_1)R_{12}(z_1z_1^{}),$$ (3.34) and $`\text{T}r_0`$ signifies the trance on the auxiliary space associated with the spectral parameters $`z_1`$ and $`z_1`$. Note that the boundary monodromy matrix $`𝒯(z,z^{})`$ is a solution to the reflection equation: $$𝒯_2(z_1^{},z_2)R_{21}(z_1+z_2)𝒯_1(z_1,z_2)R_{12}(z_1z_2)=R_{21}(z_1^{}z_1)𝒯_1(z_1,z_2)R_{12}(z_1+z_1^{})𝒯_2(z_1^{},z_2).$$ (3.35) ###### Proposition 3.9 If one takes $$K_{}(z)=K(z,v),K_+(z)=K(z\frac{n}{2}w,v^{})\text{End}(V_0),$$ (3.36) where $`v`$ and $`v^{}`$ are arbitrary parameters, the transfer matrices (3.33) commute each other : $$[T_L(z_1,z_2),T_L(z_1^{},z_2)]=0.$$ (3.37) \[Proof\] From the crossing symmetry (2.28) and the unitarity (2.27) we have $$\begin{array}{cc}& T_L(z_1,z_2)T_L(z_1^{},z_2)\hfill \\ =& \text{Tr}_1K_1(z_1\frac{n}{2}w)𝒯_1(z_1,z_2)\text{Tr}_2K_2(z_1^{}\frac{n}{2}w)𝒯_2(z_1^{},z_2)\hfill \\ =& \text{Tr}_1\text{Tr}_2K_2(z_1^{}\frac{n}{2}w)K_1^{t_1}(z_1\frac{n}{2}w)𝒯_1^{t_1}(z_1,z_2)𝒯_2(z_1^{},z_2)\hfill \\ =& \text{Tr}_1\text{Tr}_2K_2(z_1^{}\frac{n}{2}w)K_1^{t_1}(z_1\frac{n}{2}w)R_{21}^{t_1}(z_1z_2nw)R_{12}^{t_1}(z_1+z_2)𝒯_1^{t_1}(z_1,z_2)𝒯_2(z_1^{},z_2)\hfill \\ =& \text{Tr}_1\text{Tr}_2K_2(z_1^{}\frac{n}{2}w)(R_{21}(z_1z_2nw)K_1(z_1\frac{n}{2}w))^{t_1}(𝒯_1(z_1,z_2)R_{12}(z_1+z_2))^{t_1}𝒯_2(z_1^{},z_2)\hfill \\ =& \text{Tr}_1\text{Tr}_2K_2(z_1^{}\frac{n}{2}w)R_{21}(z_1z_2nw)K_1(z_1\frac{n}{2}w)R_{12}(z_2z_1)\hfill \\ \times & R_{21}(z_1z_2)𝒯_1(z_1,z_2)R_{12}(z_1+z_2)𝒯_2(z_1^{},z_2),\hfill \end{array}$$ where we use $`\text{Tr}AB=\text{Tr}A^tB^t`$. Furthermore, from (3.35) we have $$\begin{array}{cc}=& \text{Tr}_1\text{Tr}_2R_{21}(z_2z_1)K_1(z_1\frac{n}{2}w)R_{12}(z_1z_2nw)K_2(z_1^{}\frac{n}{2}w)\hfill \\ \times & 𝒯_2(z_1^{},z_2)R_{21}(z_1+z_2)𝒯_1(z_1,z_2)R_{12}(z_1z_2)\hfill \\ =& \text{Tr}_1\text{Tr}_2K_1(z_1\frac{n}{2}w)(K_2(z_1^{}\frac{n}{2}w)R_{12}(z_1z_2nw))^{t_2}(R_{21}(z_1+z_2)𝒯_2(z_1^{},z_2))^{t_2}𝒯_1(z_1,z_2)\hfill \\ =& \text{Tr}_1\text{Tr}_2K_1(z_1\frac{n}{2}w)(R_{12}(z_1z_2nw)K_2(z_1^{}\frac{n}{2}w))^{t_2}(𝒯_2(z_1^{},z_2)R_{21}(z_1+z_2))^{t_2}𝒯_1(z_1,z_2)\hfill \\ =& \text{Tr}_1\text{Tr}_2K_1(z_1\frac{n}{2}w)K_2^{t_2}(z_1^{}\frac{n}{2}w)𝒯_2^{t_2}(z_1^{},z_2)𝒯_1(z_1,z_2)\hfill \\ =& T_L(z_1^{},z_2)T_L(z_1,z_2),\hfill \end{array}$$ that implies the commutativity (3.37). $`\mathrm{}`$ ## 4 Boundary CTM bootstrap In this section we construct lattice realization of vertex operators and the boundary vacuum states for the boundary Belavin model. ### 4.1 Partition function Let us consider the inhomogeneous lattice $`_{LM}`$ with $`2M`$ horizontal lines carrying alternating spectral parameters $`z_1`$ and $`z_1`$ and $`L(0`$ mod $`n`$) vertical lines carrying the spectral parameter $`z_2`$ as below: The lattice $`_{LM}`$ and the $`i`$-th ground state. The arrows stand for the orientation of the spectral parameters. The dots $``$’s stand for the boundary interaction $`K(z)`$. For the sake of simplicity, we here denote the state $`i\pm 1`$ and $`i\pm 2`$ by $`i_\pm `$ and $`i_{\pm \pm }`$, respectively. A zigzag line on which the state variables take $`i+1`$ is presented for transparency. In this paper we restrict ourselves to the principal regime $`0<t<q<r<u_\pm <1`$, where $`u_\pm =\mathrm{exp}\left(\pi \sqrt{1}(z_1\pm z_2)\right)`$. In this regime of parameters, the bulk Boltzmann weights of the type $`R(z)_{j,j+1}^{j+1,j}`$ dominates the others; and the boundary Boltzmann weight $`K_i^i(z)`$ is the largest among $`K_i^j(z)`$ for fixed $`i`$. Thus in the low temperature limit $`t,q0`$, only the configuration such that the spin variables take the same value along the zigzag line (see the above figure) and increase by one in the direction from West to East, is possible. We call it a configuration of the ground state labeled by the boundary state $`i_n`$. Actually, the boundary weight $`K_0^0(z)`$ (and $`K_m^m(z)`$ if $`n`$ is even) are the largest among $`K_i^i(z)`$. We therefore have only one real ground state for odd $`n`$ and two for even $`n`$. Nevertheless, we regard all $`n`$ kinds of configurations as the ground states. In what follows, we fix one of them (say, $`i`$) and define all the correlation functions in terms of the low-temperature series expansion (i.e., the formal power series of $`t`$ and $`q`$). Then the lowest order of them comes from the $`i`$-th ground state configuration. Furthermore, any finite order contribution is derived from the configurations which differ from that of the $`i`$-th ground state by altering a finite number of spins. It is equivalent to taking the GNS representation obtained from the $`i`$-th ground state ($`i`$-th GNS representation) as the Hilbert space. It is expected that the correlation function defined in such a way is an analytic function which has a finite convergence radius if there exists the phase transition at a finite temperature. Following we conjecture the partition function $`Z_{LM}^{(i)}(z_1,z_2)`$ of this model behaves in the thermodynamic limit $`L,M\mathrm{}`$ as $$\begin{array}{ccc}\hfill \mathrm{log}Z_{LM}^{(i)}(z_1,z_2)& & LM\left(\mathrm{log}\mu ^{(i)}(z_1z_2)+\mathrm{log}\mu ^{(i)}(z_1+z_2)\right)\hfill \\ & +& M\left(\mathrm{log}\nu ^{(i)}(z_1)+\mathrm{log}\nu ^{(i)}(z_1\frac{n}{2}w)\right).\hfill \end{array}$$ (4.1) Here $`\mu ^{(i)}(z)`$ is the partition function per cite for the bulk theory, and $`\nu ^{(i)}(z)`$ is the that per boundary cite, which are normalized as follows: $$\mu ^{(i)}(z)=1,\nu ^{(0)}(z)=1,\nu ^{(m)}(z)=1,\text{if }n\text{ is even}.$$ (4.2) Next we consider the boundary CTM lattice as below: The inhomogeneous CTM lattice split into four sections. We denote the SW and NW corner transfer matrices by $`A_{SW}^{(i)}(z_1,z_2)`$ and $`A_{NW}^{(i)}(z_1,z_2)`$, respectively; and also denote the upper and lower lines of $`K(z)`$ by $`{}_{i}{}^{}B\|`$ and $`\|B_i`$, respectively. Let $$\begin{array}{ccc}\hfill ^{(i)}& :=& \{\mathrm{}v_{p(3)}v_{p(2)}v_{p(1)}\|p(j)_n,p(j)=i+1j\text{ (mod }n\text{) for }j1\},\hfill \\ \hfill \overline{}^{(i)}& :=& \{\mathrm{}v_{p(3)}v_{p(2)}v_{p(1)}\|p(j)_n,p(j)=i\text{ (mod }n\text{) for }j1\},\hfill \end{array}$$ (4.3) and $`^{(i)}`$ and $`\overline{}^{(i)}`$ be their dual spaces. Then in the infinite lattice limit we conclude that $`\|B_i\overline{}`$, $`{}_{i}{}^{}B\|\overline{}^{(i)}`$, and $$\begin{array}{cc}\hfill A_{SW}^{(i)}(z_1,z_2):& \overline{}^{(i)}^{(i)},\hfill \\ \hfill A_{NW}^{(i)}(z_1,z_2):& ^{(i)}\overline{}^{(i)}.\hfill \end{array}$$ (4.4) The partition function is given as follows: $$Z^{(i)}(z_1,z_2)={}_{i}{}^{}B\|A_{NW}^{(i)}(z_1,z_2)A_{SW}^{(i)}(z_1,z_2)\|B_{i}^{}.$$ (4.5) ### 4.2 Vertex operators Let us introduce the type I vertex operators where the sub/superscripts $`(i\pm 1,i)`$ specify the spaces intertwined by the vertex operators. We often suppress these sub/superscripts when we have no fear of confusion. It follows from the Yang–Baxter equation that these vertex operators satisfy the following commutation relations : $$\begin{array}{ccc}\hfill \varphi ^{j_2}(z_2)\varphi ^{j_1}(z_1)& =& \underset{j_1^{},j_2^{}}{}(R^{V_{z_1},V_{z_2}})_{j_1^{}j_2^{}}^{j_1j_2}\varphi ^{j_1^{}}(z_1)\varphi ^{j_2^{}}(z_2),\hfill \\ \hfill \varphi ^{j_2}(z_2)\varphi ^{j_1}(z_1)& =& \underset{j_1^{},j_2^{}}{}(R^{V_{z_1},V_{z_2}^{}})_{j_1^{}j_2^{}}^{j_1j_2}\varphi ^{j_1^{}}(z_1)\varphi ^{j_2^{}}(z_2),\hfill \\ \hfill \varphi ^{j_2}(z_2)\varphi ^{j_1}(z_1)& =& \underset{j_1^{},j_2^{}}{}(R^{V_{z_1}^{},V_{z_2}^{}})_{j_1^{}j_2^{}}^{j_1j_2}\varphi ^{j_1^{}}(z_1)\varphi ^{j_2^{}}(z_2).\hfill \end{array}$$ (4.6) Furthermore, the unitarity relations for $`R`$-matrices imply the inversion relation of the vertex operators: $$\underset{j}{}\varphi _j(z)\varphi ^j(z)=1,\underset{j}{}\varphi _j^{}(z)\varphi ^j(z)=1.$$ (4.7) From the crossing symmetry we have $$\varphi ^j(z)=\varphi _j(z\frac{n}{2}w),\varphi _j^{}(z)=\varphi ^j(z\frac{n}{2}w).$$ (4.8) Using these vertex operators, the transfer matrix for the semi-infinite lattice is defined as follow: $$\begin{array}{ccc}\hfill T_B(z_1,z_2)& =& \underset{j,k}{}\varphi _j(z_1+z_2)K_k^j(z_1)\varphi ^k(z_1z_2)\hfill \\ & =& \underset{j,k}{}\varphi ^j(z_1\frac{n}{2}wz_2)K_k^j(z_1)\varphi ^k(z_1z_2).\hfill \end{array}$$ (4.9) If the $`i`$-th vauuam states $`\|\text{vac}_i`$ and $`{}_{i}{}^{}\text{vac}\|`$ satisfy the following reflection properties: $$\begin{array}{ccc}\hfill \underset{k}{}K_k^j(z)\varphi ^k(z)\|\text{vac}_i& =& \nu ^{(i)}(z)\varphi ^j(z)\|\text{vac}_i,\hfill \\ \hfill {}_{i}{}^{}\text{vac}\|\underset{k}{}\varphi _k(z)K_j^k(z)& =& \nu ^{(i)}(z){}_{i}{}^{}\text{vac}\|\varphi _j(z),\hfill \end{array}$$ (4.10) these vacuums are the eigenstates of $`T_B(z,0)`$ associated with the eigenvalues $`\nu ^{(i)}(z)`$, respectively: $$T_B(z,0)\|\text{vac}_i=\nu ^{(i)}(z)\|\text{vac}_i,_i\text{vac}\|T_B(z,0)=\nu ^{(i)}(z){}_{i}{}^{}\text{vac}\|.$$ For $`n=2`$, it is suffices to consider only two types vertex operators $`\varphi ^j(z)`$ and $`\varphi _j(z)`$ because of $`\varphi ^j(z)=\varphi _{1j}(zw)`$ and $`\varphi _j^{}(z)=\varphi ^{1j}(zw)`$ . Furthermore, from $`T_B(z_1,z_2)=T_B(z_1w,z_2)`$ for $`n=2`$, we have $$\begin{array}{cc}& \underset{j,k}{}\varphi ^{1j}(z_1wz_2)K_k^j(z_1)\varphi ^k(z_1z_2)\hfill \\ =& \underset{j^{},k^{}}{}\varphi ^{1k^{}}(z_1z_2)K_j^{}^k^{}(z_1w)\varphi ^j^{}(z_1wz_2)\hfill \\ =& \underset{\genfrac{}{}{0pt}{}{j,k}{j^{},k^{}}}{}R(2z_1w)_{1jk}^{j^{}\mathrm{\hspace{0.17em}1}k^{}}\varphi ^{1j}(z_1wz_2)\varphi ^k(z_1z_2)K_k^{}^j^{}(z_1w),\hfill \end{array}$$ (4.11) which implies the boundary crossing symmetry (3.16). The crucial point in (4.11) consists in the self-duality $`\varphi _j^{}(z)=\varphi ^{1j}(z)`$ for $`n=2`$. Thus the boundary crossing symmetry (3.16) does not have a simple generalization for $`n>2`$. We should rather regard the RHS of (3.16) for general $`n`$ as the definition of the dual $`K`$-matrix. In order to see that, let us repeat the reduction (4.11) for general $`n`$. Using eqs. (4.8), (4.10), (4.6) and (4.7) we have $$\begin{array}{ccc}\hfill \nu ^{(i)}(z)& =& \underset{j^{},k^{}}{}{}_{i}{}^{}\text{vac}\|\varphi ^k^{}(z\frac{n}{2}w)K_j^{}^k^{}(z)\varphi ^j^{}(z)\|\text{vac}_{i}^{}\hfill \\ & =& \underset{\genfrac{}{}{0pt}{}{j,k}{j^{},k^{}}}{}{}_{i}{}^{}\text{vac}\|\varphi ^j(z)(R^{V_z,V_{znw/2}^{}})_{jk}^{j^{}k^{}}K_j^{}^k^{}(z)\varphi ^k(z\frac{n}{2}w)\|\text{vac}_{i}^{}\hfill \\ & =& \underset{\genfrac{}{}{0pt}{}{j,k}{j^{},k^{}}}{}{}_{i}{}^{}\text{vac}\|\varphi _j^{}(z\frac{n}{2}w)(R^{V_z,V_{znw/2}^{}})_{jk}^{j^{}k^{}}K_j^{}^k^{}(z)\varphi ^k(z\frac{n}{2}w)\|\text{vac}_{i}^{}.\hfill \end{array}$$ Thus, if we define the dual $`K`$-matrix by $$K^{}(z\frac{n}{2}w)_k^j:=\underset{j^{},k^{}}{}(R^{V_z,V_{znw/2}^{}})_{jk}^{j^{}k^{}}K(z)_j^{}^k^{},$$ (4.12) then the following dual reflection properties hold: $$\begin{array}{ccc}\hfill \underset{k}{}K^{}(z)_k^j\varphi ^k(z)\|\text{vac}_i& =& \nu ^{(i)}(z\frac{n}{2}w)\varphi ^j(z)\|\text{vac}_i,\hfill \\ \hfill {}_{i}{}^{}\text{vac}\|\underset{k}{}\varphi _k^{}(z)K^{}(z)_j^k& =& \nu ^{(i)}(z\frac{n}{2}w){}_{i}{}^{}\text{vac}\|\varphi _j^{}(z).\hfill \end{array}$$ (4.13) The associativity condition of the algebra (4.6) and (4.13) implies the reflection equations involving $`K^{}`$-matrices: $$\begin{array}{ccc}\hfill K_2(z_2)R_{21}^{V_{z_2},V_{z_1}^{}}K_1^{}(z_1)R_{12}^{V_{z_1}^{},V_{z_2}}& =& R_{21}^{V_{z_2},V_{z_1}^{}}K_1^{}(z_1)R_{12}^{V_{z_1}^{},V_{z_2}}K_2(z_2),\hfill \\ \hfill K_2^{}(z_2)R_{21}^{V_{z_2}^{},V_{z_1}^{}}K_1^{}(z_1)R_{12}^{V_{z_1}^{},V_{z_2}^{}}& =& R_{21}^{V_{z_2}^{},V_{z_1}^{}}K_1^{}(z_1)R_{12}^{V_{z_1}^{},V_{z_2}^{}}K_2^{}(z_2).\hfill \end{array}$$ (4.14) ### 4.3 Derivation of the reflection properties In this subsection we derive the reflection properties (4.10,4.13). For that purpose we introduce following further two types of vertex operators: where the sub/superscripts $`(i\pm 1,i)`$ specify the spaces intertwined by the vertex operators. Hereafter we also suppress these sub/superscripts. From the reflection equation (3.1) we have the following relation: $`{\displaystyle \underset{k}{}}K(z_3)_k^j\phi ^k(z_1,z_3)\|B_i`$ $`=`$ $`\nu ^{(i)}(z_3)\phi ^j(z_1,z_3)\|B_i.`$ (4.15) By similar argument we have $`{\displaystyle \underset{k}{}}{}_{i}{}^{}B\|\phi _k(z_1,z_3)K(z_3)_j^k`$ $`=`$ $`\nu ^{(i)}(z_3){}_{i}{}^{}B\|\phi _j(z_1,z_3),`$ (4.16) Furthermore, we have the relations $`A_{SW}^{(i1)}(z_1,z_2)\phi ^j(z_1,z_3)\|B_i`$ $`=`$ $`\varphi ^j(z_3z_2)A_{SW}^{(i)}(z_1,z_2)\|B_i,`$ (4.17) $`{}_{i}{}^{}B\|\phi _j(z_1,z_3)A_{NW}^{(i1)}(z_1,z_2)`$ $`=`$ $`{}_{i}{}^{}B\|A_{NW}^{(i)}(z_1,z_2)\varphi _j(z_2z_3).`$ (4.18) These are based on the unitarity and Yang-Baxter relation of $`R`$-matrix in the thermodynamic limit. The unitarity (2.28) allows us to obtain Using the Yang–Baxter equation (2.8) we get By successive use of the YBE and the unitarity we can bring the line associated with the spectral parameter $`z_3`$ to the directions pointed by dotted lines in the above figure as far as we like. Thus we find Those manipulation implies (4.17) because the contribution of Boltzmann weights along the tail graphically represented in the figure by the dotted line is unity in the thermodynamic limit. The relation (4.18) can be similarly obtained. Applying $`A_{SW}^{(i1)}(z_1,z_2)`$ (resp. $`A_{NW}^{(i1)}(z_1,z_2)`$) from the left (resp. right) to both sides of (4.15) (resp. (4.16)) and using (4.17) (resp. (4.18)) we obtain $`{\displaystyle \underset{k}{}}K(z_3)_k^j\varphi ^k(z_3z_2)A_{SW}^{(i)}(z_1,z_2)\|B_i`$ $`=`$ $`\nu ^{(i)}(z_3)\varphi ^j(z_3z_2)A_{SW}^{(i)}(z_1,z_2)\|B_i,`$ (4.19) $`{\displaystyle \underset{k}{}}{}_{i}{}^{}B\|A_{NW}^{(i)}(z_1,z_2)\varphi _k(z_2+z_3)K(z_3)_j^k`$ $`=`$ $`\nu ^{(i)}(z_3){}_{i}{}^{}B\|A_{NW}^{(i)}(z_1,z_2)\varphi _j(z_2z_3).`$ (4.20) Taking account of (4.19) and (4.20) with (4.10) we find the following identification $$\|\text{vac}_i=A_{SW}^{(i)}(z_1,z_2=0)\|B_i,_i\text{vac}\|={}_{i}{}^{}B\|A_{SW}^{(i)}(z_1,z_2=0).$$ (4.21) From the identification (4.21) and the definition of the dual $`K`$-matrix (4.12) we obtain $`{\displaystyle \underset{k}{}}K^{}(z_3)_k^j\varphi ^k(z_3z_2)A_{SW}^{(i)}(z_1,z_2)\|B_i`$ $`=`$ $`\nu ^{(i)}(z_3\frac{n}{2}w)\varphi ^j(z_3z_2)A_{SW}^{(i)}(z_1,z_2)\|B_i,`$ (4.22) $`{\displaystyle \underset{k}{}}{}_{i}{}^{}B\|A_{NW}^{(i)}(z_1,z_2)\varphi _k^{}(z_2+z_3)K^{}(z_3)_j^k`$ $`=`$ $`\nu ^{(i)}(z_3\frac{n}{2}w){}_{i}{}^{}B\|A_{NW}^{(i)}(z_1,z_2)\varphi _j^{}(z_2z_3).`$ (4.23) ## 5 Correlation functions and difference equations The relations appeared in the previous section are not rigorous because all the objects are defined on the infinite lattice. Nevertheless we assume that eqs. (4.1–4.23) are exactly correct on the basis of the CTM bootstrap method, which is supported by some numerical calculations and consistency with the vertex operator method . ### 5.1 Local state probabilities Let us consider the correlation function on the dislocated CTM lattice Thanks to (4.17) and (4.18) we have $$\begin{array}{cc}& G_N^{(i)}(z,z^{}\|z_1^{},\mathrm{},z_N^{},z_N,\mathrm{},z_1)^{j_1^{},\mathrm{},j_N^{},j_N,\mathrm{},j_1}\hfill \\ =& {}_{i}{}^{}B\|A_{SW}^{(i)}(z,z^{})\varphi _{j_1^{}}(z^{}z_1^{})\mathrm{}\varphi _{j_N^{}}(z^{}z_N^{})\varphi ^{j_N}(z_Nz^{})\mathrm{}\varphi ^{j_1}(z_1z^{})A_{SW}^{(i)}(z,z^{})\|B_{i}^{}.\hfill \end{array}$$ (5.1) Thus the correlation function $`G_N^{(i)}(z,z^{}\|z_1^{},\mathrm{},z_N^{},z_N,\mathrm{},z_1)^{j_1^{},\mathrm{},j_N^{},j_N,\mathrm{},j_1}`$ normalized by the partition function (4.5) is called the $`N`$-point local state probability of the boundary Belavin model if we set $`z_l=z_l^{}=z^{}=0`$, $`j_l=j_l^{}`$ ($`1lN`$). Owing to the unitarity (4.7) we have $$Z^{(i)}(z_1,z_2)=\underset{j_1,\mathrm{}j_N}{}G_N^{(i)}(z,z^{}\|z_1,\mathrm{},z_N,z_N,\mathrm{},z_1)^{j_1,\mathrm{},j_N,j_N,\mathrm{},j_1}.$$ (5.2) Thus we obtain the expression of the $`n`$-point local state probability: $$P_N^{(i)}(j_1,\mathrm{},j_N)=\frac{G_N^{(i)}(z,0\|0,\mathrm{},0)^{j_1,\mathrm{},j_N,j_N,\mathrm{},j_1}}{{\displaystyle \underset{j_1,\mathrm{}j_N}{}}G_N^{(i)}(z,0\|0,\mathrm{},0)^{j_1,\mathrm{},j_N,j_N,\mathrm{},j_1}}$$ (5.3) ### 5.2 Boundary analogue of the quantum Knizhnik–Zamolodchikov equation Only for $`n=2`$, the $`N`$-point function (5.1) is reduced to the following $`2N`$-point function of the form $$\begin{array}{cc}& F_{2N}^{(i)}(z,z^{}\|y_1,\mathrm{},y_N,z_N,\mathrm{},z_1)^{j_1^{},\mathrm{},j_N^{},j_N,\mathrm{},j_1}\hfill \\ =& {}_{i}{}^{}B\|A_{SW}^{(i)}(z,z^{})\varphi ^{k_1}(y_1z^{})\mathrm{}\varphi ^{k_N}(y_Nz^{})\varphi ^{j_N}(z_Nz^{})\mathrm{}\varphi ^{j_1}(z_1z^{})A_{SW}^{(i)}(z,z^{})\|B_{i}^{},\hfill \end{array}$$ (5.4) by putting $`y_l=z_l^{}w`$ and $`k_l=1j_l^{}`$ for $`1lN`$ . It is nothing to do with any local state probabilities for $`n>2`$, however, we can consider the correlation function of (5.4)-type: $$\begin{array}{ccc}\hfill F_N^{(i)}(z,z^{}\|z_1,\mathrm{},z_N)& =& \underset{j_1,\mathrm{},j_N}{}v_{j_1}\mathrm{}v_{j_N}F_N^{(i)}(z,z^{}\|z_1,\mathrm{},z_N)^{j_1,\mathrm{},j_N},\hfill \\ \hfill F_N^{(i)}(z,z^{}\|z_1,\mathrm{},z_N)^{j_1,\mathrm{},j_N}& =& {}_{i}{}^{}B\|A_{SW}^{(i)}(z,z^{})\varphi ^{j_1}(z_1z^{})\mathrm{}\varphi ^{j_N}(z_Nz^{})A_{SW}^{(i)}(z,z^{})\|B_{i}^{}.\hfill \end{array}$$ (5.5) Here we assume that $`N0`$ mod $`n`$ for simplicity. From the same discussion as in , we obtain ###### Proposition 5.1 The correlation function (5.5) satisfies the following relations: $`1.R\text{-matrix symmetry}:\text{ }`$ $`P_{jj+1}F_N^{(i)}(z,z^{}\|\mathrm{},z_{j+1},z_j,\mathrm{})`$ $`=`$ $`R_{jj+1}^{V_{z_j},V_{z_{j+1}}}F_N^{(i)}(z,z^{}\|\mathrm{},z_j,z_{j+1},\mathrm{}),`$ (5.6) $`2.\text{Reflection property }I:`$ $`K_N(z_N)F_N^{(i)}(z,z^{}\|z_1,\mathrm{},z_{N1},z_N)`$ $`=`$ $`\nu ^{(i)}(z_N)F_N^{(i)}(z,z^{}\|z_1,\mathrm{},z_{N1},z_N),`$ (5.7) $`3.\text{Reflection property }II:`$ $`\widehat{K}_1(z_1)F_N^{(i)}(z,z^{}\|z_1,z_2,\mathrm{},z_N)`$ $`=`$ $`\nu ^{(i)}(z_1)F_N^{(i)}(z,z^{}\|z_1nw,z_2,\mathrm{},z_N),`$ (5.8) where $$\widehat{K}(z)v_k=\underset{j}{}v_jK^{}(z\frac{n}{2}w)_j^k.$$ \[Proof\] The first equation (5.6) follows from the commutation relation (4.6), while the second one (5.7) follows from (4.19). Finally, from the crossing relation (4.8) and (4.23) $$\begin{array}{cc}& \widehat{K}_1(z_1)F_N^{(i)}(z,z^{}\|z_1,z_2,\mathrm{},z_N)\hfill \\ \hfill =& \underset{j_1^{},j_1,\mathrm{},j_N}{}v_{j_1}\mathrm{}v_{j_N}{}_{i}{}^{}B\|A_{SW}^{(i+N)}(z,z^{})\varphi _{j_1^{}}^{}(z^{}z_1\frac{n}{2}w)\mathrm{}A_{SW}^{(i)}(z,z^{})\|B_{i}^{}K_1^{}(z_1\frac{n}{2}w)_{j_1}^{j_1^{}}\hfill \\ \hfill =& \nu ^{(i)}(z_1)\underset{j_1,\mathrm{},j_N}{}v_{j_1}\mathrm{}v_{j_N}{}_{i}{}^{}B\|A_{SW}^{(i+N)}(z,z^{})\varphi _{j_1}^{}(z^{}+z_1+\frac{n}{2}w)\mathrm{}A_{SW}^{(i)}(z,z^{})\|B_{i}^{},\hfill \end{array}$$ we obtain the last equation (5.8). $`\mathrm{}`$ Owing to the equations (5.65.8) we obtain ###### Theorem 5.2 The correlation function (5.5) satisfies the following difference equation: $$\begin{array}{ccc}\hfill T_jF_N^{(i)}(z,z^{}\|z_1,\mathrm{},z_N)& =& R_{jj1}^{V_{z_jnw},V_{z_{j1}}}\mathrm{}R_{j1}^{V_{z_jnw},V_{z_1}}\widehat{K}_j(z_j)\hfill \\ & \times & R_{1j}^{V_{z_1,V_{z_j}}}\mathrm{}R_{j1j}^{V_{z_{j1},V_{z_j}}}R_{j+1j}^{V_{z_{j+1},V_{z_j}}}\mathrm{}R_{Nj}^{V_{z_N,V_{z_j}}}\hfill \\ & \times & K_j(z_j)R_{jN}^{V_{z_j},V_{z_N}}\mathrm{}R_{jj+1}^{V_{z_j},V_{z_{j+1}}}F_N^{(i)}(z,z^{}\|z_1,\mathrm{},z_N),\hfill \end{array}$$ (5.9) where $$T_jf(z,z^{}\|z_1,\mathrm{},z_j,\mathrm{},z_N)=f(z,z^{}\|z_1,\mathrm{},z_jnw,\mathrm{},z_N).$$ Using the crossing symmetries we have another expression of the correlation function on the dislocated CTM lattice for general $`n2`$: $$\begin{array}{cc}& G_N^{(i)}(z,z^{}\|z_1^{},\mathrm{},z_N^{},z_N,\mathrm{},z_1)^{j_1^{},\mathrm{},j_N^{},j_N,\mathrm{},j_1}\hfill \\ =& {}_{i}{}^{}B\|A_{SW}^{(i)}(z,z^{})\varphi ^{j_1^{}}(z_1^{}z^{})\mathrm{}\varphi ^{j_N^{}}(z_N^{}z^{})\varphi ^{j_N}(z_Nz^{})\mathrm{}\varphi ^{j_1}(z_1z^{})A_{SW}^{(i)}(z,z^{})\|B_{i}^{},\hfill \end{array}$$ (5.10) where $`z_l^{}=z_l^{}\frac{n}{2}w`$ for $`1lN`$. We thus introduce the $`V^nV^n`$-valued correlation function $$\begin{array}{cc}& G_N^{(i)}(z,z^{}\|z_1^{},\mathrm{},z_N^{},z_N,\mathrm{},z_1)\hfill \\ =& \underset{\genfrac{}{}{0pt}{}{j_1,\mathrm{},j_N}{j_1^{},\mathrm{}j_n^{}}}{}v_{j_1^{}}^{}\mathrm{}v_{j_N^{}}^{}v_{j_N}\mathrm{}v_{j_1}G_N^{(i)}(z,z^{}\|z_1^{},\mathrm{},z_N^{},z_N,\mathrm{},z_1)^{j_1^{},\mathrm{},j_N^{},j_N,\mathrm{},j_1}.\hfill \end{array}$$ (5.11) Let us describe the $`R`$-matrix symmetry corresponding to (5.6). ###### Proposition 5.3 Let $$\begin{array}{cc}& G_N^{(\sigma i)}(z,z^{}\|x_{\sigma (1)},\mathrm{},x_{\sigma (2N)})\hfill \\ =& \underset{\genfrac{}{}{0pt}{}{j_1,\mathrm{},j_N}{j_1^{},\mathrm{}j_n^{}}}{}v_{j_1^{}}^{}\mathrm{}v_{j_N^{}}^{}v_{j_N}\mathrm{}v_{j_1}G_N^{(i)}(z,z^{}\|x_{\sigma (1)},\mathrm{},x_{\sigma (2N)})^{k_{\sigma (1)},\mathrm{},k_{\sigma (N)}},\hfill \\ & G_N^{(\sigma i)}(z,z^{}\|x_{\sigma (1)},\mathrm{},x_{\sigma (2N)})^{k_{\sigma (1)},\mathrm{},k_{\sigma (2N)}}\hfill \\ =& {}_{i}{}^{}B\|A_{SW}^{(i)}(z,z^{})\mathrm{\Phi }^{\sigma (1)}\mathrm{}\mathrm{\Phi }^{\sigma (2N)}A_{SW}^{(i)}(z,z^{})\|B_{i}^{}.\hfill \end{array}$$ (5.12) Here $`\sigma `$ be the permutation of $`(1,\mathrm{},2N)`$, and $$x_l=\{\begin{array}{cc}z_l^{}=z_l^{}\frac{n}{2}w,\hfill & (1lN);\hfill \\ z_{2N+1l},\hfill & (N+1l2N);\hfill \end{array}k_l=\{\begin{array}{cc}j_l^{},\hfill & (1lN);\hfill \\ j_{2N+1l},\hfill & (N+1l2N);\hfill \end{array}$$ and $$\mathrm{\Phi }^l=\{\begin{array}{cc}\varphi ^{k_l}(x_lz^{}),\hfill & (1lN);\hfill \\ \varphi ^{k_l}(x_lz^{}),\hfill & (N+1l2N).\hfill \end{array}$$ Then the following $`R`$-matrix symmetry holds: $$G_N^{(\sigma _ji)}(z,z^{}\|\mathrm{},x_{\sigma (j+1)},x_{\sigma (j)},\mathrm{})=R_{\sigma (j),\sigma (j+1)}^{V^{\sigma (j)},V^{\sigma (j+1)}}G_N^{(\sigma i)}(z,z^{}\|\mathrm{},x_{\sigma (j)},x_{\sigma (j+1)},\mathrm{}),$$ (5.13) where $$V^l=\{\begin{array}{cc}V_{x_l}^{}\hfill & (1lN);\hfill \\ V_{x_l}\hfill & (N+1l2N);\hfill \end{array}$$ and $`\sigma _j`$ is the permutation of $`(1,\mathrm{},2N)`$ obtained from $`\sigma `$ by transposing $`\sigma (j)`$ and $`\sigma (j+1)`$. The reflection properties can be similarly shown as before: ###### Proposition 5.4 The following relations holds: $`K_{2N}(z_1)G_N^{(\pi i)}(z,z^{}\|\mathrm{},z_1)`$ $`=`$ $`\nu ^{(i)}(z_1)G_N^{(\pi i)}(z,z^{}\|\mathrm{},z_1),`$ (5.14) $`\widehat{K}_{2N}(z_1)G_N^{(\rho i)}(z,z^{}\|z_1,\mathrm{})`$ $`=`$ $`\nu ^{(i)}(z_1)T_1G_N^{(\rho i)}(z,z^{}\|z_1,\mathrm{}),`$ (5.15) $`\widehat{K}_1^{}(z_1^{})G_N^{(\varsigma i)}(z,z^{}\|z_1^{},\mathrm{})`$ $`=`$ $`\nu ^{(i)}(z_1^{}\frac{n}{2}w)T_1G_N^{(\varsigma i)}(z,z^{}\|z_1^{},\mathrm{}),`$ (5.16) $`K_1^{}(z_1^{})G_N^{(\tau i)}(z,z^{}\|\mathrm{},z_1^{})`$ $`=`$ $`\nu ^{(i)}(z_1^{}\frac{n}{2}w)G_N^{(\tau i)}(z,z^{}\|\mathrm{},z_1^{}),`$ (5.17) Here, $$\widehat{K}^{}(z)v_k^{}=\underset{j}{}v_j^{}K(z\frac{n}{2}w)_j^k,$$ and $`\pi `$, $`\rho `$, $`\varsigma `$, $`\tau 𝔖_{2N}`$ such that $$\pi (2N)=2N,\rho (1)=2N,\varsigma (1)=1,\tau (2N)=1.$$ \[Proof\] The relation (5.13) is evident from the commutation relations (4.6). The last two (5.16) and (5.17) follow from (4.20), (4.8) and (4.22). $`\mathrm{}`$ From Propositions 5.3 and 5.4, we have ###### Theorem 5.5 Let $`V_1^l=V_{x_l}`$, $`V_2^l=V_{x_lnw}`$. Then the following difference equations holds $$\begin{array}{ccc}\hfill T_lG_N^{(i)}(z,z^{}\|x_1,\mathrm{},x_{2N})& =& R_{ll1}^{V_2^l,V^{l1}}\mathrm{}R_{l\mathrm{\hspace{0.17em}1}}^{V_2^l,V^1}\widehat{K}_l^{}(x_l)\hfill \\ & \times & R_{1l}^{V^1,V_1^l}\mathrm{}R_{l1l}^{V^{l1},V_1^l}R_{l+1l}^{V^{l+1},V_1^l}\mathrm{}R_{2Nl}^{V^{2N},V_1^l}\hfill \\ & \times & K_l^{}(x_l)R_{l\mathrm{\hspace{0.17em}2}N}^{V^l,V^{2N}}\mathrm{}R_{ll+1}^{V^l,V^{l+1}}G_N^{(i)}(x_1,\mathrm{},x_{2N}),\hfill \end{array}$$ (5.18) for $`1lN`$, and $$\begin{array}{ccc}\hfill T_lG_N^{(i)}(z,z^{}\|x_1,\mathrm{},x_{2N})& =& R_{ll1}^{V_2^l,V^{l1}}\mathrm{}R_{l\mathrm{\hspace{0.17em}1}}^{V_2^l,V^1}\widehat{K}_l(x_l)\hfill \\ & \times & R_{1l}^{V^1,V_1^l}\mathrm{}R_{l1l}^{V^{l1},V_1^l}R_{l+1l}^{V^{l+1},V_1^l}\mathrm{}R_{2Nl}^{V^{2N},V_1^l}\hfill \\ & \times & K_l(x_l)R_{l\mathrm{\hspace{0.17em}2}N}^{V^l,V^{2N}}\mathrm{}R_{ll+1}^{V^l,V^{l+1}}G_N^{(i)}(x_1,\mathrm{},x_{2N}),\hfill \end{array}$$ (5.19) for $`N+1l2N`$. Theorem 5.5 gives an elliptic generalization of the corresponding difference equations for the boundary $`U_q(\widehat{sl_n})`$-symmetric model . ### 5.3 Boundary spontaneous polarization Applying the similar argument as in (5.9) to the simplest case $`N=1`$ we obtain the following difference equations: $$\begin{array}{ccc}\hfill T_1G_1^{(i)}(z,z^{}\|z_1^{},z_2)& =& \widehat{K}_1^{}(z_1^{})R_{21}^{V_{z_2},V_{z_1^{}}^{}}K_1^{}(z_1^{})R_{12}^{V_{z_1^{}}^{},V_{z_2}}G_1^{(i)}(z,z^{}\|z_1^{},z_2),\hfill \\ \hfill T_2G_1^{(i)}(z,z^{}\|z_1^{},z_2)& =& R_{21}^{V_{z_2nw},V_{z_1^{}}^{}}\widehat{K}_2(z_2)R_{12}^{V_{z_1^{}}^{},V_{z_2}}K_2(z_2)G_1^{(i)}(z,z^{}\|z_1^{},z_2),\hfill \end{array}$$ (5.20) where $`z_1^{}=z_1\frac{n}{2}w`$. It is difficult to get each element of $`G_1^{(i)}(z,z^{}\|z_1,z_2)`$, however, it is possible to obtain the expression of the following sums: $$P_m^{(i)}(z,z^{}\|z_1,z_2)=\underset{j=0}{\overset{n1}{}}\omega ^{mj}G_1^{(i)}(z,z^{}\|z_1\frac{n}{2}w,z_2)^{jj}.$$ (5.21) Note that the boundary spontaneous polarization as the vacuum expectation value of the operator $`g`$ at boundary is expressed in terms of (5.21) as follows: $$g^{(i)}=\frac{P_1^{(i)}(z,z^{}=0\|z_1,z_2)}{P_0^{(i)}(z,z^{}=0\|z_1,z_2)}\|_{z_1=z_2=z^{}}.$$ (5.22) Now we restrict ourselves to the free boundary condition $`r1`$ for simplicity. Since $`\underset{r1}{lim}𝒦(0)𝒦_0`$, the initial condition does not hold if we take $`\overline{K}(z)=𝒦_0𝒦(z)`$. Thus we should regard the $`K`$-matrix in this limit as $`\overline{K}(z)=𝒦(0)𝒦(z)`$. Under this identification the $`K`$-matrix behaves as $$K(z)k(z)I_n,$$ where $`k(z)`$ is a scalar function of $`z`$. Here we cite the following sum formula from <sup>2</sup><sup>2</sup>2 Note that there are typographical errors in the formula . $$\underset{j=0}{\overset{n1}{}}\omega ^{mj}\frac{\theta ^{(j)}(z+w)}{\theta ^{(j)}(w)}=n\frac{h((zm)/n+w)_{lm}h((z+l)/n)}{h(w)_{l0}h(l/n)},$$ (5.23) Then we see the dual $`K`$-matrix in the free boundary limit $`r1`$ behaves as $$K^{}(z\frac{n}{2}w)k(z)f_0(u^2q^n)I_n,$$ where $$\begin{array}{ccc}\hfill f_m(u)& \hfill :=& \underset{j=0}{\overset{n1}{}}\omega ^{mj}R(z)_{0j}^{j0}\hfill \\ & \hfill =& \frac{1}{\overline{\kappa }(u)}\frac{(\omega ^mq^2u^{2/n};t^2)_{\mathrm{}}(t^2\omega ^mq^2u^{2/n};t^2)_{\mathrm{}}}{(\omega ^mu^{2/n};t^2)_{\mathrm{}}(t^2\omega ^mu^{2/n};t^2)_{\mathrm{}}}.\hfill \end{array}$$ (5.24) The difference equations (5.20) are therefore reduced to $$\begin{array}{ccc}\hfill T_1G_1^{(i)}(z,z^{}\|z_1^{},z_2)^{jj}& =& f_0(u_1^2q^n)\underset{k,l}{}R_{12}(z_1z_2)_{jk}^{kj}R_{21}(z_2z_1)_{kl}^{lk}G_1^{(i)}(z,z^{}\|z_1^{},z_2)^{ll},\hfill \\ \hfill T_2G_1^{(i)}(z,z^{}\|z_1^{},z_2)^{jj}& =& f_0(u_2^2q^n)\underset{k,l}{}R_{12}(z_1z_2)_{jk}^{kj}R_{21}(z_1z_2)_{kl}^{lk}G_1^{(i)}(z,z^{}\|z_1^{},z_2)^{ll},\hfill \end{array}$$ (5.25) where $`z_1^{}=z_1\frac{n}{2}w`$, and we use (2.28) and (3.11). Substituting (5.25) into (5.21) we obtain $$P_m^{(i)}(z,z^{}\|z_1,z_2)=C_m^{(i)}A(u_1)A(u_2)B_m(u_+)B_m(u_{}).$$ (5.26) Here $`C_m^{(i)}`$ is a constant, and $`A(u)`$ and $`B_m(u)`$ are solutions to the following difference equations: $$\frac{A(uq^n)}{A(u)}=f_0(u^2q^n),\frac{B_m(uq^n)}{B_m(u)}=f_m(u).$$ (5.27) By solving these difference equations we obtain $$A(u)=\psi (u^2)\frac{(q^2u^{4/n};t^2,q^4)_{\mathrm{}}(q^4u^{4/n};t^2,q^4)_{\mathrm{}}}{(t^2u^{4/n};t^2,q^4)_{\mathrm{}}(t^2q^2u^{4/n};t^2,q^4)_{\mathrm{}}},$$ (5.28) where $$\psi (u):=g_0(uq^{n/2})g(u^1q^{n/2}),g_0(u):=\frac{(q^{2+3n}u^2;t^2,q^{2n},q^{4n})_{\mathrm{}}(t^2q^{2+3n}u^2;t^2,q^{2n},q^{4n})_{\mathrm{}}}{(q^{3n}u^2;t^2,q^{2n},q^{4n})_{\mathrm{}}(t^2q^{3n}u^2;t^2,q^{2n},q^{4n})_{\mathrm{}}};$$ and $$B_m(u)=\phi (u)\frac{(t^2\omega ^mu^{2/n};t^2)_{\mathrm{}}(t^2\omega ^mu^{2/n};t^2)_{\mathrm{}}}{(q^2\omega ^mu^{2/n};q^2)_{\mathrm{}}(q^2\omega ^mu^{2/n};q^2)_{\mathrm{}}},$$ (5.29) where $$\phi (u):=g(uq^{n/2})g(u^1q^{n/2}),g(u):=\frac{(q^{3n}u^2;t^2,q^{2n},q^{2n})_{\mathrm{}}(t^2q^{3n}u^2;t^2,q^{2n},q^{2n})_{\mathrm{}}}{(q^{2+n}u^2;t^2,q^{2n},q^{2n})_{\mathrm{}}(t^2q^{2+n}u^2;t^2,q^{2n},q^{2n})_{\mathrm{}}}.$$ Note that $`B_m(u)`$ is essentially the same as $`G^{(m)}(u)`$ in , which corresponds to the quantity (5.21) in the bulk theory. From (5.26) we have $`{\displaystyle \frac{P_1^{(i)}(z,z^{}=0\|z_1,z_2)}{P_0^{(i)}(z,z^{}=0\|z_1,z_2)}}`$ $`=`$ $`{\displaystyle \frac{C_1^{(i)}}{C_0^{(i)}}}{\displaystyle \frac{B_1(u_+)B_1(u_{})}{B_0(u_+)B_0(u_{})}}.`$ (5.30) Taking the low temperature limit $`t,q0`$, we find that the ratio $`C^{(i)}/C_0^{(i)}`$ should be equal to $`\omega ^i`$. We therefore obtain the boundary spontaneous polarization from (5.30) and (5.29) by putting $`u_+=u_{}=1`$ $$g^{(i)}=\omega ^i\frac{(q^2;q^2)_{\mathrm{}}^4}{(t^2;t^2)_{\mathrm{}}^4}\frac{(t^2\omega ;t^2)_{\mathrm{}}^2(t^2\omega ^1;t^4)_{\mathrm{}}^2}{(q^2\omega ;q^2)_{\mathrm{}}^2(q^2\omega ^1;q^4)_{\mathrm{}}^2}.$$ (5.31) When $`n=2`$ this expression coincides with the previous result obtained in . We also emphasize that the boundary spontaneous polarization for the boundary Belavin model is exactly the square of that for the bulk Belavin model obtained in , up to a phase factor. ## 6 Summary and discussion In this paper we have obtained two non-diagonal solutions of the reflection equation associated with Belavin’s $`_n`$-symmetric elliptic model. Unfortunately, our elliptic $`K`$-matrix is not connected with the diagonal boundary Boltzmann weights for the $`A_{n1}^{(1)}`$-face model but the non-diagonal ones. It is thus an open problem to obtain the $`K`$-matrix corresponding to the boundary Boltzmann weights given in . On the basis of the boundary CTM bootstrap we have derived a set of difference equations for correlation functions of the boundary Belavin model. By solving the simplest difference equations, we have obtained the boundary spontaneous polarization of the boundary Belavin model. Our result is consistent with the one given in when $`n=2`$. The boundary spontaneous polarization is equal to the square of the bulk spontaneous polarization up to a phase factor. The same phenomena were observed in . In this paper we have shown that correlation functions of the boundary model satisfy the $`R`$-matrix symmetry and the reflection properties, which are the boundary analogue of Smirnov’s first two axioms . It may be interesting to construct integral formulae for correlation functions such that the integrand possesses the determinant structure as in Smirnov’s integral . In integral formulae for correlation functions of the boundary $`XYZ`$ model by using bosonization of vertex operators . In order to obtain the higher $`n`$ generalization of , the construction of free field realization of the boundary Belavin model is required. It is a very hard but important work. ## Acknowledgement The author would like to thank A. Kuniba and A. Nakayashiki for discussion.
warning/0003/hep-ph0003204.html	ar5iv	text	# Untitled Document McGILL-00-02 UA/NPPS-01-00 Polarized Photoproduction of Heavy Quarks in Next-to-Leading Order<sup>+</sup> A.P. Contogouris<sup>a,b,1</sup>, Z. Merebashvili<sup>a,∗,2</sup> and G. Grispos<sup>b,3</sup> a. Department of Physics, McGill University, Montreal, Qc., H3A 2T8, Canada b. Nuclear and Particle Physics, University of Athens, Athens 15771, Greece ABSTRACT The results of a next-to-leading order calculation of heavy quark production in longitudinally polarized photon-nucleon collisions are presented. At c.m. energy $`\sqrt{S}=10`$ GeV, for $`\stackrel{}{\gamma }+\stackrel{}{p}c+X`$, cross sections differential in the transverse momentum and rapidity of the charmed quark $`c`$ and the corresponding asymmetries are presented; also, as functions of $`\sqrt{S}`$, integrated cross sections, $`K`$-factors and the corresponding asymmetries are given. Errors in the asymmetries are estimated and the possibility to distinguish between three scerarios differing essentially in the polarized gluon distribution is discussed. <sup>+</sup><sup>+</sup>footnotetext: Also supported by the Secretariat of Research and Technology of Greece and by the Natural Sciences and Engineering Research Counsil of Canada. <sup></sup><sup></sup>footnotetext: Present address: High Energy Physics Institute, Tbilisi State University, University St. 9, 380086 Tbilisi, Republic of Georgia. <sup>1</sup><sup>1</sup>footnotetext: e-mail: apcont@physics.mcgill.ca, acontog@cc.uoa.gr <sup>2</sup><sup>2</sup>footnotetext: e-mail: mereb@sun20.hepi.edu.ge <sup>3</sup><sup>3</sup>footnotetext: e-mail: ggrispos@cc.uoa.gr In spite of the welth of data on polarized deep inelastic scattering, the size and shape of the polarized gluon distribution $`\mathrm{\Delta }g(x)`$ remains a central problem in Spin Physics. Important progress needs experiments on processes with longitudinally polarized initial particles dominated by subprocesses with initial gluons. Such a process is $$\stackrel{}{\gamma }+\stackrel{}{p}Q\overline{(Q)}+X,$$ (1) where $`Q\overline{(Q)}`$ denotes heavy quark (antiquark), and, in general, it is dominated by the subprocess $$\stackrel{}{\gamma }+\stackrel{}{g}Q+\overline{Q}$$ (2) Proposals on experiments closely related to, or in, (1) exist in various stages of approval $`[`$1$`]`$. On the other hand, the importance of determining higher order corrections (HOC) is well known. In this Letter we report the essential results of a calculation of the (next-to-leading) HOC. The subprocesses contributing to (1) are as follows: * At leading order (LO, $`\alpha \alpha _s`$) $`[`$2, 3$`]`$ * The Born subprocess of (2). * The resolved $`\gamma `$ via $`\stackrel{}{q}\stackrel{}{\overline{q}}Q\overline{Q}`$ and $`\stackrel{}{g}\stackrel{}{g}Q\overline{Q}`$. These involve the polarized photon structure functions $`\mathrm{\Delta }F_{q/\gamma }`$ and $`\mathrm{\Delta }F_{g/\gamma }`$, known only theoretically; hopefully, more information will eventually come from experiments on (1). * At next-to-leading order (NLO, $`\alpha \alpha _s^2`$) * The loop and Bremsstrahlung (Brems, i.e. $`\stackrel{}{\gamma }\stackrel{}{g}Q\overline{Q}g`$) associated with (2). * The subprocess $`\stackrel{}{\gamma }\stackrel{}{q}Q\overline{Q}q`$ At NLO, scheme independent cross sections can be rigorously obtained only for the sum of resolved and direct contributions. Notice, that the Abelian part of (B1) provides the HOC to $$\stackrel{}{\gamma }\stackrel{}{\gamma }Q\overline{Q};$$ (3) these HOC have been determined $`[`$4, 5$`]`$. HOC to the process (3) are of interest in themselves: In searches of the Higgs boson (mass $`m_H`$), for $`90`$ GeV$`m_H2m_W`$, in future $`\gamma \gamma `$ colliders, the dominant decay mode is $`Hb\overline{b}`$. With polarized $`\gamma `$’ s, the Born contribution to the background $`\stackrel{}{\gamma }\stackrel{}{\gamma }b\overline{b}`$ is much suppressed; however, due to gluon Brems, HOC have an important effect $`[`$4, 5$`]`$. A calculation of NLO corrections for (1) already exists $`[`$6$`]`$. Our work, however, makes use of a different regularization scheme; also, in the treatment of the soft and collinear contributions, contrary to $`[`$6$`]`$ which separates them from the hard parts via a cut parameter, we apply more conventional methods $`[`$4$`]`$. Thus, in view of the importance of (1), we believe that an independent calculation is worthwhile. Comparisons, as much as possible, with the results of $`[`$6$`]`$ will be also reported. Most conveniently, singularities are eliminated by dimensional methods. For polarized reactions this requires extension of the Dirac matrix $`\gamma _5`$ in $`n=42\epsilon `$ dimensions. There are several schemes for this, and, as in $`[`$4$`]`$, we follow that of dimensional reduction. This scheme violates the Ward identity between the vertex and quark self energy functions, but care has been taken by introducing a (finite) counterterm, as discussed in $`[`$4$`]`$. The wave function and mass renormalizations are carried on shell $`[`$4$`]`$. In the present case charge renormalization is also required. We introduce $$A_\epsilon \left(m\right)=\left(\frac{g}{4\pi }\right)^2\left(\frac{4\pi \mu ^2}{m^2}\right)^\epsilon \mathrm{\Gamma }\left(1+\epsilon \right)$$ (4) where $`\mu `$ is an arbitrary mass scale entering in $`n=42\epsilon `$ dimensions via the change $`gg\mu ^\epsilon `$. Let $`N_{lf}=`$ number of light flavors and $`b\left(11N_c2N_{lf}\right)/6`$. Then we carry charge renormalization by introducing the counterterm $$\frac{1}{\epsilon }\left(A_\epsilon \left(\mu _R\right)b\frac{1}{3}A_\epsilon \left(m\right)\right)$$ (4a) where $`\mu _R`$ is a regularization mass. In this scheme graphs containing internal loops of the heavy quark $`Q`$ are subtracted out, so that $`Q`$ is decoupled. This is consistent with parton distributions of which the evolution is determined from split functions involving only light quarks, as is our case. Finally, the $`gQ\overline{Q}`$ vertex was renormalized via the Slavnov-Taylor identities $`[`$7$`]`$. Now, loop contributions to (B1) were determined via Passarino-Veltman techniques $`[`$8$`]`$; and $`23`$ parton contributions by going to the c.m. (Gottfried-Jackson) frame of $`\overline{Q}(Q)`$ and final $`g`$ $`[`$9, 4$`]`$. Also, (B2) was treated by going to the c.m. frame of $`\overline{Q}(Q)`$ and final $`q`$. The integrals listed in the Appendices A and C of $`[`$9$`]`$ have been very useful. Some remaining integrals are given in an Appendix of $`[`$10$`]`$. Singularities $`1/\epsilon ^2`$ appearing in the cross sections of (B1) are cancelled by adding loops and Brems. The cancellation of singularities $`1/\epsilon `$ appearing in (B1) and (B2) required the addition of factorization counterterms. With $`p_1,p_2,p_3`$ the 4-momenta of $`\gamma `$, initial parton and observed $`Q(\overline{Q})`$, define $$s=(p_1+p_2)^2,t=(p_3p_1)^2m^2,u=(p_3p_2)^2m^2.$$ Our counterterms are defined in the $`\overline{MS}`$ scheme and have the general form: $$\mathrm{\Delta }\frac{d\sigma _{cter}}{dtdu}=\frac{1}{\epsilon }\frac{K(\epsilon )}{s}\mathrm{\Delta }P_{ab}(x)\left(\frac{m^2}{M^2}\right)^\epsilon \mathrm{\Delta }\frac{d\widehat{\sigma }_B}{dt}$$ (5) where $`\mathrm{\Delta }P_{ab}(x)`$ split function (in $`n=4`$ dimensions, see below), $`x`$ proper dimensionless variable, $`M`$ the factorization scale, $`K(\epsilon )`$ a kinematic factor determined from phase-space and color and $`\mathrm{\Delta }d\widehat{\sigma }_B/dt`$ the Born cross section of a $`22`$ subprocess with $`s`$ and either $`t`$ or $`u`$ replaced by $`xs`$ and either $`xt`$ or $`xu`$. For $`\stackrel{}{\gamma }\stackrel{}{g}Q\overline{Q}g`$ the $`22`$ subprocess is (2); for $`\stackrel{}{\gamma }\stackrel{}{q}Q\overline{Q}q`$ one needs two counterterms, one involving (2) and another involving $`\stackrel{}{\overline{q}}\stackrel{}{q}Q\overline{Q}`$ $`[`$11$`]`$. The resulting finite cross sections are convoluted with polarized parton distributions whose evolution is determined by 2-loop anomalous dimensions $`[`$12, 13$`]`$. On their basis, several groups have constructed sets of such distributions, differing mainly in the shape and size of $`\mathrm{\Delta }g`$. We use throughout the NLO sets of one group $`[`$14$`]`$. Also, we use the NLO expression of the running coupling constant $`\alpha _s(\mu )`$ with $`\mathrm{\Lambda }=231`$ MeV and $`N=N_{lf}+1=4`$ flavours. The scheme for extending $`\gamma _5`$ in $`n`$ dimensions used in $`[`$12, 13$`]`$ is not dimensional reduction, so the addition of certain conversion terms is necessary. The form of the conversion terms is easily given in terms of Eq. (5): In $`n`$ dimensions the split functions have the form: $$\mathrm{\Delta }P_{ba}^n(x,\epsilon )=\mathrm{\Delta }P_{ba}\left(x\right)+\epsilon \mathrm{\Delta }P_{ba}^\epsilon \left(x\right);$$ (6) it is $`\mathrm{\Delta }P_{ba}^\epsilon \left(x\right)`$ that depends on the scheme. The conversion terms are determined from the difference of $`\mathrm{\Delta }P_{ba}^\epsilon \left(x\right)`$ in the different schemes. In dimensional reduction $`\mathrm{\Delta }P_{ba}^\epsilon \left(x\right)\left(=P_{ba}^\epsilon \left(x\right)\right)`$ $`=0`$. Refs. $`[`$12, 13$`]`$ use the t’ Hooft-Veltman scheme, modified so that $`\mathrm{\Delta }P_{qq}^n(x,\epsilon )=P_{qq}^n(x,\epsilon )`$; then $`\mathrm{\Delta }P_{ab}^\epsilon (x,)0`$. In terms of these $`\mathrm{\Delta }P_{ab}^\epsilon \left(x\right)`$, the conversion terms are $$\frac{d\sigma _{conv}}{dtdu}=\frac{K(0)}{s}\mathrm{\Delta }P_{ab}^\epsilon (x)\frac{d\widehat{\sigma }_B}{dt}$$ (7) where $`d\widehat{\sigma }_B/dt`$ the Born cross section of (5) scaled in the same way. For the polarized parton distributions we use the sets A, B and C of Ref. 14 and for the unpolarized the most recent version CTEQ5 $`[`$15$`]`$. Finally, in the absence of any experimental information, as an estimate of the resolved $`\gamma `$ contributions, we have used the maximal and minimal saturation LO sets of $`\mathrm{\Delta }F_{q/\gamma }`$ and $`\mathrm{\Delta }F_{g/\gamma }`$ of $`[`$16$`]`$, as well as the set of the asymptotic solutions $`[`$17$`]`$. In brief, those of $`[`$17$`]`$ give the largest contributions, whereas those of the minimal saturation set give the smallest. The analytical calculations were carried with REDUCE and to some extent with FORM. Subsequently we present results for $`Q=`$c-quark with $`m=1.5`$ GeV only at $`\sqrt{S_{\gamma p}}=\sqrt{S}=10`$ GeV, relevant to the experiments (a) and (b) of $`[`$1$`]`$. Higher energies and $`Q=`$b-quark are considered elsewhere $`[`$10$`]`$. Also, in relation with (4a), we take $`\mu _R=\mu `$. Fig. 1 presents results related with the differential cross sections $`\mathrm{\Delta }d\sigma /dp_T`$ versus $`x_T2p_T/\sqrt{S}`$, where $`p_T`$ the transverse momentum of $`Q`$; measurement of such cross sections is possible in (b) of $`[`$1$`]`$. Subsequently we denote by $`\mathrm{\Delta }d\sigma _B/dp_T`$, $`\mathrm{\Delta }d\sigma _{res}/dp_T`$, and $`\mathrm{\Delta }d\sigma _{\gamma q}/dp_T`$ the contributions to the physical cross section of (A1), (A2), and (B2) correspondingly, and by $`\mathrm{\Delta }d\sigma /dp_T`$ that of the sum (A1), (B1) and (B2). We use the scale $`\mu =M=\left(p_T^2+m^2\right)^{1/2}`$; the stability of our results against variations of $`\mu `$ and $`M`$ is studied in $`[`$10$`]`$. In Fig. 1(a) the cross sections $`\mathrm{\Delta }d\sigma /dp_T`$ and the corresponding $`\mathrm{\Delta }d\sigma _B/dp_T`$ (denoted by a $``$) are determined for sets A, B and C of $`[`$14$`]`$. The presented $`\mathrm{\Delta }d\sigma _{res}/dp_T`$ and $`\mathrm{\Delta }d\sigma _{\gamma q}/dp_T`$ correspond to set B and the former to the maximal saturation set of $`[`$16$`]`$. In calculating $`\mathrm{\Delta }d\sigma _B/dp_T`$ and $`\mathrm{\Delta }d\sigma _{res}/dp_T`$ we use the NLO sets of $`[`$14$`]`$.<sup>(a)</sup> <sup>(a)</sup><sup>(a)</sup>footnotetext: In this way the effect of the perturbative HOC, as it is reflected e.g. in the magnitude of $`K`$-factors (see Eqs. (10)), is made more clear. This is particularly true for polarized reactions, in which the LO and NLO $`\mathrm{\Delta }g`$ differ significantly. Fig. 1(b) presents the asymmetries $$A_{LL}(p_T)=\frac{\mathrm{\Delta }d\sigma /dp_T}{d\sigma /dp_T}$$ (8) for sets A, B and C. The resolved $`\gamma `$ contributions have been left out in view of their smallness and of the fact that the stage of their present knowledge does not permit a scheme independent calculation. Following usual practice in calculating $`d\sigma /dp_T`$ we average over $`n2`$ spin degrees of freedom for every incoming boson. Finally, the errors in Fig. 1(b) have been estimated from: $$\delta A_{LL}=\frac{1}{P_BP_T\sqrt{L\sigma ϵ}}$$ (9) In (9) we use unpolarized cross section $`\sigma `$ integrated over a bin of $`x_T`$ corresponding to $`\mathrm{\Delta }p_T=0.5`$ GeV and the conditions of $`[`$1$`a]`$ ($`P_B=80\%`$, $`P_T=25\%`$, $`ϵ=0.014`$ and $`L=2`$ $`fb^1`$). Note that the proposal $`[`$1$`b]`$ amounts to better conditions and thus smaller $`\delta A_{LL}`$. On the basis of our calculations and Fig. 1 we remark the following: * Defining by $`\mathrm{\Delta }d\sigma _{\gamma g}/dp_T`$ the contribution of (A1) and (B1) we may introduce the $`K`$-factors $$K=\mathrm{\Delta }\frac{d\sigma }{dp_T}/\mathrm{\Delta }\frac{d\sigma _B}{dp_T},K_{\gamma g}=\mathrm{\Delta }\frac{d\sigma _{\gamma g}}{dp_T}/\mathrm{\Delta }\frac{d\sigma _B}{dp_T}$$ (10) For sets A and B and $`x_T0.3`$, where perturbative QCD is more trustworthy, the $`K`$-factors are $`>1`$ and, in particular $`K_{\gamma g}`$, fairly large. The largeness of $`K_{\gamma g}`$ is partly due to the fact that $`\sqrt{S}=10`$ GeV is fairly low, so $`\alpha _s(\mu )`$ is fairly large. * $`\mathrm{\Delta }d\sigma _{\gamma q}/dp_T`$ are found to change little between sets A, B and C; the reason is that the valence distributions $`\mathrm{\Delta }u`$ and $`\mathrm{\Delta }d`$ vary little. Thus Fig. 1(a) presents $`\mathrm{\Delta }d\sigma _{\gamma q}/dp_T`$ for only one set (B). * In general, $`\mathrm{\Delta }d\sigma _{\gamma q}/dp_T`$ are significantly smaller than $`\mathrm{\Delta }d\sigma _{\gamma g}/dp_T`$ for sets A and B. Note that also in unpolarized photoproduction we find that $`\gamma qQ\overline{Q}q`$ contributes much less than the NLO $`\gamma gQ\overline{Q}`$, in full accord with $`[`$18$`]`$. * Most importantly, Fig. 1(b) shows that at $`x_T0.4`$ one can distinguish sets A and C and perhaps also all A, B, C. Fig. 2 presents results related with the distributions $`\mathrm{\Delta }d\sigma /dY`$ where $`Y`$ the c.m. rapidity of $`Q`$ with respect to the photon. The subsequent notation is equivalent to the above and parts (a) and (b) are equivalent to these of Fig. 1. Here we use the scale $`\mu =M=2m`$. The errors in Fig. 2(b) have been estimated using Eq. (9) with the conditions of $`[`$1$`a]`$ and the unpolarized cross section integrated over a bin $`\mathrm{\Delta }Y=1`$. On the basis of Fig. 2 we remark: * At $`Y>0`$, $`\mathrm{\Delta }d\sigma /dY`$ for sets A, B are significantly larger than for C and than $`\mathrm{\Delta }d\sigma _{\gamma q}/dY`$ and $`\mathrm{\Delta }d\sigma _{res}/dY`$. (Fig. 2(a)) * Taking into account the errors, Fig. 2(b) suggests that $`Y1.251.5`$ is the best region to distinguish set C from A or B. The region $`Y<0`$ leads to large errors because $`d\sigma /dY`$, like $`\mathrm{\Delta }d\sigma /dY`$ (Fig. 2(a)), is small. * Figs 2(a) and 2(b) suggest that integrating the cross sections in the range $`1Y1.5`$ offers perhaps the best possibility to distinguish A or B form C. Finally, Fig. 3 presents results related with the integrated cross sections $`\mathrm{\Delta }\sigma `$ versus the c.m. energy $`\sqrt{S_{\gamma p}}=\sqrt{S}`$ in a range including $`\sqrt{S}=10`$ GeV. This quantity will be measured in both experiments $`[`$1$`a]`$ and $`[`$1$`b]`$. Here the scale is again $`\mu =M=2m`$. To show clearly the NLO effects, we present $`K`$-factors (Fig. 3(a)) in addition to integrated cross sections (Fig. 3(b)). At $`\sqrt{S}=10`$ GeV the $`K`$-factors for all sets exceed $`K=1`$. In the range $`7<\sqrt{S}<14`$ GeV for sets A and B, $`K`$ are smooth, but for C, $`K`$ is discontinuous due to the vanishing of $`\mathrm{\Delta }\sigma _B`$ at $`\sqrt{S}11`$ GeV. Most important is Fig. 3(c), which presents NLO asymmetries. The error at $`\sqrt{S}=10`$ GeV is estimated using in (9) again the conditions of $`[`$1$`a]`$. Under these conditions our results show that sets A and C can be distinguished, but not sets A and B or B and C. Perhaps the proposed SLAC experiment $`[`$1$`b]`$, which will give results at somewhat lower $`\sqrt{S}`$ and, as stated, amounts to better conditions, can distinguish also B and C. At $`\sqrt{S}10`$ GeV of importance is the precise knowledge of the charmed quark mass $`m`$. We understand that recently the uncertainity in $`m`$ has been somewhat reduced<sup>(b)</sup>. Using the set B, we find that varying $`m`$ in the range $`1.35m1.65`$ changes $`\mathrm{\Delta }\sigma `$ in the range $`27.11\mathrm{\Delta }\sigma 13.34`$ nb and the asymmetry in $`6.26\%A_{LL}8.88\%`$. <sup>(b)</sup><sup>(b)</sup>footnotetext: We would like to thank A. Despande for this information. Ref. $`[`$6$`]`$ works in the scheme of Ref. 19 and presents most of its results using the ”standard” set of $`[`$20$`]`$ (GRSVst); for integrated cross sections (Fig. 3), it also presents results for sets A and C of $`[`$14$`]`$. Also, as stated, they treat their soft and collinear gluon parts by the phase space slicing method of $`[`$18$`]`$, which separates them from the hard gluon parts via a cut parameter. First, we have checked that our virtual, soft and collinear contributions are in complete agreement (Appendix C of the first reference of $`[`$6$`]`$). Second, with respect to numerical comparisons, care is needed in convoluting the LO cross sections with proper (LO or NLO) parton distributions. Using also GRSVst, and comparing with their corrected results, regarding differential cross sections we find differences not exceeding $`3\%`$; however, regarding asymmetries of integrated cross sections, although in fair agreement, our differences are somewhat larger. In conclusion, the COMPASS experiment $`[`$1$`a]`$ is expected to provide useful information on $`\mathrm{\Delta }g`$. This may be said with more emphasis for the proposed SLAC experiment $`[`$1$`b]`$. * We thank I. Bojak for providing us with the results for several quantities we are comparing and for taking his part in doing comparisons. Thanks are also due to G. Bunce, D. de Florian, B. Kamal and J. Körner for discussions, to W. Vogelsang for discussions and for providing us the sets of $`[`$20$`]`$, to P. Bosted for several communications, to A. Despande for useful information and remarks and to V. Spanos and G. Veropoulos for participating in part of the calculation. REFERENCES * (a) G. Baum et al, COMPASS Collaboration: CERN/SPLC 96-14 and 96-30; (b) P. Bosted, SLAC-PROPOSAL-E156, 1997; (c) W.-D. Nowak, DESY 96-095; (d) A. de Roeck and T. Gehrmann, DESY-Proceedings-1998-1. * M. Gluck and E. Reya, Z. Phys. C39 (1988) 569; M. Stratmann and W. Vogelsang, Z. Phys. C74 (1997) 641; A. Watson, ibid C12 (1982) 123. * B. Lampe and E. Reya, MPI-PhT/98-23 and DO-TH $`98/02^{(+)}`$. <sup>(+)</sup><sup>(+)</sup>footnotetext: This is also the most up to date review of polarized particle reactions * B. Kamal, Z. Merebashvili and A.P. Contogouris, Phys. Rev. D51 (1995) 4808; ibid D55 (1997) 3229 (E). * G. Jikia and A. Tkabladze, ibid D54 (1996) 2030. * I. Bojak and M. Stratmanm: Nucl. Phys. B540 (1999) 345 and Erratum (to be published); Phys. Lett. B433 (1998) 411. * T. Muta, ”Foundations of Quantum Chromodynamics” (World Scientific, 1987) * G. Passarino and M. Veltman, Nucl. Phys. B160 (1979) 151. * W. Beenakker et al, Phys. Rev. D40 (1989) 54. * Z. Merebashvili, A.P. Contogouris and G. Grispos, in preparation. See also proceedings of ”Spin 99” International Workshop, Prague, 8-12 September 1999. * A.P. Contogouris, S. Papadopoulos and B. Kamal, Phys. Lett. B246 (1990) 523. * R. Mertig and W. van Neerven, Z. Phys. C70 (1996) 637. * W. Vogelsang, Phys. Rev. D54 (1996) 2023; Nucl. Phys. B475 (1996) 47. * T. Gehrmann and W. Stirling, Phys. Rev. D53 (1996) 6100. * H. Lai et al, Eur. Phys. J. C12 (2000) 375. * M. Gluck and W. Vogelsang, Z. Phys. C55 (1992) 353 and C57 (1993) 309; M. Gluck, M. Stratmann and W. Vogelsang, Phys. Lett. B187 (1994) 373. * J. Hassan and D. Pilling, Nucl. Phys. B187 (1981) 563. * J. Smith and W.L. van Neerven, Nucl. Phys. B374 (1992) 36. * G. t’ Hooft and M. Veltman, ibid B44 (1972) 189. * M. Gluck, E. Reya, M. Stratmann and W. Vogelsang, Phys. Rev. D53 (1996) 4775. FIGURE CAPTIONS * Quantities related with the $`p_T`$-distributions versus $`x_T=2p_T/\sqrt{S}`$: (a) Polarized differential cross sections. The LO cross sections are indicated by $``$. (b) NLO asymmetries for sets A, B and C. * Quantities related with the rapidity distributions: (a) and (b) as in Fig. 1. * Quantities related with the integrated cross sections for sets A, B and C: (a) Factors $`K=\mathrm{\Delta }\sigma /\mathrm{\Delta }\sigma _B`$ (b) LO (indicated by $``$) and NLO cross sections. The presented cross section for $`\stackrel{}{\gamma }\stackrel{}{q}Q\overline{Q}q`$ corresponds to set B. (c) NLO asymmetries.
warning/0003/cs0003056.html	ar5iv	text	# A note on the Declarative reading(s) of Logic Programming ## Introduction Logic’s fundamental role in the area of computing and artificial intelligence, is its use for knowledge representation. There may be innumerable ways in which some domain knowledge can be encoded in a logic theory; however, there is one principle which most consider as the canonical way of using logic. Declarative knowledge representation operates according to the following principle: > the expert represents his domain knowledge by a set of formal statements that are true in the problem domain. This principle relies on the ability of the expert to interpret a logical axiom as a clear and precise statement about the domain of discourse. The ability of interpreting the formulas of a logic as meaningful statements about the problem domain -given some interpretation of the user-defined symbols- is the declarative reading of the logic. It is based on a clear understanding of the connectors and quantifiers and of how composed axioms combine these meanings. A declarative semantics of the logic can be defined as a formal study of its -intuitive- declarative reading. As such it should contain the following parts: * a clear account of some declarative reading of the formulas and theories in a logic; * a mathematical characterisation of a formal semantics; * a justification why and how this semantics characterises this declarative reading. In a sense, a declarative semantics should relate a logic to some part of the human cognition and understanding. Note that a simple way of defining a declarative semantics of a logic is by providing an embedding in another logic with a well-established declarative semantics. Such semantics are sometimes called transformational semantics. The declarative reading of logical axioms gives logic a decisive advantage over other languages. It allows the expert to compare a formal statement with his or her knowledge and to evaluate its truth, without going through the painful process of explicitly constructing the mathematical semantics of the axioms. This ideal is reached most clearly in classical logic. So, our common expert knowledge allows us to recognize the statement $$x.person(x)male(x)female(x)$$ as true (given the obvious intended interpretation of the predicate symbols), and the statement $$x.male(x)female(x)$$ as false. We can do so on the basis of our intuitive understanding, without having to construct the semantics of these sentences, i.e. their class of models. It is a fact that to be able to represent knowledge in a declarative way and to benefit the potential advantages of this type of knowledge representation, a human expert must have acquired a deep and precise understanding of the declarative reading of the logic that he is using. Ambiguities and unclarities on the level of declarative reading cause ambiguities and unclarities on the level of knowledge representation. For this reason, studies of declarative reading are of key importance in the development of knowledge representation methodologies. This paper is a study of the declarative reading(s) of logic programming. The history of the declarative semantics of logic programming is well-known. Originally, the picture was simple and clear: the declarative reading of a Horn logic program was the declarative reading of a set of classical logic implications. The introduction of negation as failure blurred this simple view. On the one hand, negation as failure derived conclusions with a strong commonsense appeal and turned out to be very useful and natural in many practical situations. On the other hand, the negation as failure inference rule was unsound with respect to the declarative reading of a program as a set of classical implications. As expressed by Przymusinski (?), “we really do not want classical logic semantics for logic programs. .. We want the semantics of a logic program to be determined more by its commonsense meaning.”. All main semantic investigations since the end of the seventies (least model semantics (?), completion (?), perfect model semantics (??), stable model semantics (?), well-founded model semantics (?)) attempted to formalise and explicitate this commonsense meaning of logic programs. The question considered in this paper is to what extent the commonsense meaning referred to by Przymusinski has been identified: what is or are the declarative reading(s) of logic programming and what are the corresponding meanings of its symbols $`not,`$. What are the semantics that correspond best to these declarative readings? The analysis of the declarative reading of logic programming is complicated by at least two factors. One complication is that different formal semantics exist (in 2-, 3- and 4-valued versions). In particular, stable and well-founded semantics are generally considered as the main ones. However, there is a second complication which is more subtle and much more dangerous. Let me try to pinpoint this problem. It is well-known that logic programming can be used for representing many different sorts of knowledge (?): (inductive) definitions, defaults, reflective knowledge of experts, etc.., both under well-founded semantics and stable model semantics. Remarkable though is that the same logic program or rule in a logic program can be used to represent knowledge with a decidedly distinct commonsense flavor. For example, we might represent the definition that dead means not alive by the rule: $$deadnotalive$$ The interpretation of this rule as a definition of dead is explicitated by the completion semantics (?), Clark’s embedding of logic programming in classical logic. If the rule defining $`dead`$ is the only rule with $`dead`$ in the head, the completion will contain the equivalence $$dead\neg alive$$ Since stable and well-founded models are models of the completion, this rule is satisfied also in these models. On the other hand, consider the reflective knowledge of the distrustful man who is unhappy if he does not know that his wife is faithful to him. His reflective knowledge can be represented by the same rule (modulo renaming) as in the above scenario: $$unhappynotwife\mathrm{\_}faithful$$ Indeed, Gelfond’s embedding of logic programs in autoepistemic logic (?) maps this rule to the autoepistemic formula: $$unhappy\neg Kwife\mathrm{\_}faithful$$ which directly represents the knowledge of the distrustful man. Gelfond’s embedding defines a transformational declarative semantics for logic programming. It formed the knowledge theoretical foundation for stable model semantics (?); therefore we are entitled to assume that stable semantics is conceived to represent this sort of knowledge. Note here that negation as failure is interpreted as a modal operator and the implication operator as classical implication. This is a common feature of all embeddings of logic programming in autoepistemic logic (AEL) and in default logic (DL). Another important embedding is Marek and Truszczynski’s one to DL (?). It maps he above rule to: $$\begin{array}{c}:\neg wife\mathrm{\_}faithful\\ \\ unhappy\end{array}$$ For an overview of different embeddings see (?)). From the commonsense point of view, there is a definite distinction between both pieces of knowledge. The definition of dead expresses that in the actual state of the world, dead and alive are mutually exclusive: the exclusive “or” holds between them. If the expert does not know whether alive is true or false, then he does not know whether dead is true or false. On the other hand, the knowledge of the distrustful man does not imply any relationship between the unhappiness of the distrustful man and the loyalty of the spouse in the actual state of the world. It is perfectly possible that in the actual state of the world, she is not loyal but he does not know and he is happy. Consider the logic program program $`\{pnotq\}`$. All main semantics coincide for it; its formal semantics (i.e. the set of its models) is the unique model $`\{p\}`$. Yet, as illustrated above, Clark’s embedding and Gelfond’s embedding assign two different commonsense meanings to this program: * In the completion, the program states that $`q`$ is false (because of the empty definition) and that $`p`$ is true iff $`q`$ is false. The model represents the state of the world in which $`p`$ is true and $`q`$ is false. According to the completion, this is the only possible state of the world. * The program in Gelfond’s embedding states that $`p`$ holds if $`q`$ is not known to be true. Since in this interpretation, the program has no knowledge about $`q`$, $`q`$ is unknown, hence $`p`$ is entailed. Note that contrary to the first reading, here $`q`$ is not known to be false. What is proven in (?), is that the stable model represents the set of believed atoms of the above theory: $`\{p\}`$ means that $`p`$ is believed, and that the truth of $`q`$ is not believed. Based on our intuition, what can be said about the possible states of the world in this reading? Obviously, since $`q`$ is unknown, there should be states in which $`q`$ is true and states in which $`q`$ is false; $`p`$ should be true always. Hence, the two possible states are $`\{p\}`$ and $`\{q,p\}`$. Moore (?) defined a possible world semantics for AEL. Intuitively, a possible world model is a set of possible states of the world according to the expert’s knowledge. The unique possible world model of the theory $`\{p\neg Kq\}`$ is exactly the set $`\{\{p\},\{q,p\}`$; this confirms our intuition. The above discussion leads to a key point of this paper: even if we know the models of a logic program, we still cannot decide the intended declarative reading. We can only know -to some extent- what is the declarative reading if we know also what is the role of the model. A mathematical definition of some collection of models of theories in a logic cannot define a declarative reading. In this sense, a model theory is not a declarative semantics in its own right. In the sequel, a model semantics will be called a possible state semantics for some declarative reading if it characterises the possible states of the world; a model theory will be called an atomic belief set semantics of some declarative reading if it characterises the sets of believed atoms in this declarative reading. It is not a simple task to search for the declarative reading(s) of logic programming. In the first place, many of the early semantical studies in the context of logic programming are not primarily concerned with finding and formalising a commonsense meaning, but are more concerned with finding a mathematical justification for the reasoning techniques in Prolog. Other studies are more focussed on the commonsense meaning, but fail to give a clear account of the formalised declarative reading and the role of the models. There is an enormous amount of mathematical results on the relationships between different model semantics. However, because atomic belief sets and possible states are incomparable objects, it is a priori not clear what these relations mean on the level of declarative reading. Moreover, in many knowledge representation examples, one can observe that the same semantics is used once as an atomic belief set semantics, once as an possible state semantics. In order to clarify the role of logic programming for Knowledge representation, the question of declarative reading of logic programming cannot be circumvented. In the rest of the paper I will try to pinpoint the main ideas on declarative reading and the confusions on this topic. ## Declarative readings of logic programming An investigating of the transformational semantics that have been proposed for logic programming, gives some insight in the possible declarative readings of logic programming. It seems to me that at least in the non-monotonic and A.I. oriented part of the logic programming community, logic programming is now routinely seen as a sub-logic of default logic or autoepistemic logic. In this view, the negation as failure is interpreted as a modal negation. This view has a natural motivation: a Prolog system is said to infer a negative literal $`notp`$ when it is unable to prove $`p`$. A natural way of modeling failure to prove in semantics is as not knowing. From here, it was natural to interpret $`notp`$ as $`\neg Kp`$. All main logic programming semantics - least model, supported model, 3-valued supported model, perfect model, stable model, 3-valued stable model and well-founded semantics - have been justified as atomic belief set semantics of diverse modal interpretations of logic programming. The methods that have been used are analogous as Gelfond and Lifschitz’s justification for stable semantics: one defines an embedding of logic programming to some non-monotone modal logic and the models of the logic program are shown to be the set of believed atoms. In these transformational semantics, models of logic programming semantics systematically play the role of a set of believed atoms. For an overview of these results, I refer to (?). On the other hand, I believe that there is also a persistent and strong intuition that among all classical models of a logic program, there is a canonical one (or at most a small number of canonical ones) which represents the unique possible state of the world. The Clark completion semantics was an early, weak attempt to identify this canonical model. In this view, negation as failure does not need to be interpreted as a modal operator: it is classical objective negation denoting falsity in the canonical model. The main questions are what declarative reading of logic programming could explain the existence of a unique possible state and how could this unique model be mathematically characterized? These questions were considered in (??). The idea is to read a logic program as an inductive definition. From the very start, logic programs were considered as definitions. This was Clark’s basic idea with the completion semantics. Note that Clark’s completed definitions are identical to the way (non-recursive) definitions are expressed, for example in Beth’s studies. The evident problem with Clark completion semantics is that it does not deal well with inductive definitions. On the other hand, least model semantics is known to deal right with positive inductive definitions such as transitive closure. Using the above terminology, the least model semantics can be said to be possible state semantics for the reading of Horn programs as inductive definitions. In fact, as pointed out in (?), the methods that have been used to characterize monotone induction are identical to those that were used to characterize least model semantics of Horn programs. Can we extend the view of Horn logic programs as inductive definitions to programs with negation? In (?), I have argued that the use of induction in mathematics is not restricted to positive induction. An example is the induction in a well-founded set. To some extend, a generalised form of non-monotone induction have been studied in the area of inductive definitions, the so called Iterated Inductive Definitions (?). As I showed, this formalism is isomorphic modulo syntactic sugar with stratified logic programs under perfect model semantics. Further on, I have pointed to several intolerable weaknesses of this stratified approach as a formalisation of generalised induction and have argued that these problems are solved by the well-founded semantics. Or, the argument there was that the well-founded semantics is a possible state semantics of the declarative reading of logic programs as generalised inductive definitions<sup>1</sup><sup>1</sup>1Note that well-founded semantics has been motivated both as an atomic belief set semantics and as a possible state semantics.. Consequently, logic programming has not a unique declarative reading. The modal view and the definition view are both consistent ways of interpreting logic programs; moreover they lead to very similar model theories, though these theories have different roles. The above phenomenon, the existence of different consistent declarative views on the same formalism is a potential source of considerable confusion. In the remaining sections, I investigate possible confusions. ## Distinguishing between declarative readings ### Comparing declarative semantics The mathematical relations between the least model, models of the completion, the perfect model, stable models and well-founded model are understood quite well. The question I consider here is what they tell about the relations between the two types of declarative readings. A naive comparison of different semantics of logic programming is misleading. For example, the collection of stable models is known to be a subset of the collection of models of the completion. What does this result mean on the level of the declarative readings underlying both semantics? Not much it seems. E.g. the possible world model of Gelfond’s embedding and Marek and Truszczynski’s embedding of the program $`\{pnotq\}`$ is the set $`\{\{p\},\{q,p\}\}`$; this is a proper superset of the (singleton) set of models of the completion. Consequently, in the case of this particular program, the completion is strictly stronger than the default or AEL reading. ### Expressing knowledge declaratively Consider a logic with an atomic belief set semantics for its declarative reading. Assume that the models of a logic theory are exactly the possible states of the world according to the expert’s knowledge. This theory encodes the expert’s knowledge but obviously, it is not necessarily a declarative representation of the expert’s knowledge. In general the declarative reading of the theory does not justify that the models are the only possible states of the world. In fact, it may well be that part of the axioms of the theory are false statements about the domain of discourse. The above phenomenon can be illustrated in LP. Since stable semantics is an extension of least model semantics, it is suitable to encode positive inductive definitions such as the one of transitive closure: $$\begin{array}{c}p(a,a)\hfill \\ p(b,c)\hfill \\ tr(X,Y)p(X,Y)\hfill \\ tr(X,Y)p(X,Z),tr(Z,Y)\hfill \end{array}$$ The unique stable model indeed represents the unique possible state of the graph and its transitive closure. But this does not necessarily mean that inductive definitions can be expressed under the default reading of logic programs. For example, the meaning of this program under Gelfond’s embedding is identical to its classical logic meaning (since AEL is a conservative extension of classical propositional logic). This problem gives rise to unsoundness. For example, with the inductive definition, the expert knows that $`p(a,b)`$ and $`tr(a,b)`$ are false. Yet, the AEL embedding entails $`\neg Kp(a,b)`$ and $`\neg Ktr(a,b)`$. In the case of the transitive closure, the modal declarative reading (as expressed under the above mentioned AEL and DL embeddings) of the axioms is true but too weak to justify the unique model. An example where the declarative reading would be plainly false can be given by a variant of the $`dead`$ and $`alive`$ example. Assume that the expert wants to represent the definition that $`dead`$ means not $`alive`$, as represented by $`dead\neg alive`$. A possible way to do this using stable semantics is: $$\begin{array}{c}deadnotalive\hfill \\ alivenotalive^{}\hfill \\ alive^{}notalive\hfill \end{array}$$ The two stable models (after projection on the two atoms $`alive,dead`$ are identical to the models of the equivalence. Yet, it is easy to see that both the AEL and DL embedding assigns false meaning to the first axiom. Indeed, it is not true that $`dead`$ is true if one does not know that $`alive`$ is true. ### Mixing different declarative readings The presence of multiple declarative readings and multiple roles of models raises complications on the level of methodology. This becomes obvious when different roles are mixed. Consider what happens if the different sorts of declarative readings are used in the same program. Reconsider the example of the distrustful man. Assume that we want to add the definition that to be $`happy`$ means not to be $`unhappy`$. According to the completion, this knowledge is correctly represented by the rule: $$happynotunhappy$$ What happens if we combine this rule with the reflective knowledge of the jealous husband. Consider the program: $$\left\{\begin{array}{c}unhappynotwife\mathrm{\_}faithful\hfill \\ happynotunhappy\hfill \end{array}\right\}$$ This program consists now of two isomorphic statements which are both true, but under different declarative readings. Note that again all semantics for this program coincide. The model is $`\{unhappy\}`$. What does this model mean? Consider the two options that arose earlier: * If this model is to be interpreted that $`unhappy`$ is true and $`wife\mathrm{\_}faithful`$ and $`happy`$ are false, then there is a mismatch with our understanding because neither the truth of $`unhappy`$ nor even the knowledge that $`unhappy`$ is true is a sufficient condition for happy being false. I.e. $`happy`$ should be unknown. * On the other hand, if the model is interpreted as a belief set, there is again mismatch with our understanding because since $`wife\mathrm{\_}faithful`$ is unknown, then also $`unhappy`$ should be unknown. Whether the model is interpreted as a possible state or as the belief set, it contains an error. ## Extending Logic Programming This paper focuses on the original logic programming formalism with only the original negation as failure, and does not investigate its extensions with classical or strong negation. At this point, it is clear that there is natural reason for that. A point of this paper is exactly that in one commonsense view on logic programming, negation as failure is classical negation. It is clear then than in this view, the formalism cannot be further extended with another classical or strong negation. Extensions of logic programming with classical negation (including disjunctive logic programming) make sense only in the modal negation view on logic programs. Logic programming extensions were designed as a way to solve a number of serious disadvantages of logic programming for knowledge representation. However, as argued in (?), the analysis of what are these disadvantages exactly, depends on the declarative reading one takes. Different declarative readings lead to different conclusions and more importantly they subsequently lead to different ways to different extensions of the formalism. In the case of the modal views, the problem was the absence of classical negation. However, as argued in (?), the problem of logic programming under the definition view is that a logic program contains definitions for all predicates. For this reason, one should extend the formalism with the possibility of leaving certain predicates open. This view is further explored in (?).

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